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# On cooling flows and the Sunyaev-Zel’dovich effect
## 1. Introduction
Galaxy clusters are extensively observed in optical, X-ray and radio bands. In the radio band, a cluster can be observed in the the Rayleigh-Jeans side of the cosmic microwave background spectrum, as a dip in the brightness temperature, due to Sunyaev-Zel’dovich effect (SZ) (Sunyaev & Zel’dovich 1972; for a comprehensive review, see Birkinshaw 1999). The SZ distortion appears as a decrement for wavelengths $`1.44`$ mm (frequencies $`218`$ GHz) and as an increment for wavelengths $`1.44`$ mm. SZ effect has the advantage that the SZ intensity, unlike that of the X-ray, does not suffer from the $`(1+z)^4`$ cosmological dimming. As discussed by numerous authors (Birkinshaw & Hughes 1994; Silverberg et.al. 1997), one can combine the X-ray and radio observations for clusters to determine cosmological parameters. This has been done in the recent years to determine the Hubble constant $`H_{}`$ (Birkinshaw 1999). The SZ signal is, however, weak and difficult to detect. Recent high signal to noise detection have been made over a wide range in wavelengths using single dish observations: at radio wavelengths (Herbig et al, 1995, Hughes & Birkinshaw 1998), millimeter wavelengths (Holzapfel et al. 1997, Pointecouteau et al. 1999) and submillimeter wavelengths (Komatsu et al. 1999). Interferometric observations have also been carried out to image the SZ effect (Jones et al. 1993, Saunders et al. 1999, Reese et al. 1999, Grego et al. 2000). Other than estimating $`H_{}`$ combining SZ and X-ray data, SZ effect alone can also be used to determine the cosmological mass density $`\mathrm{\Omega }_{}`$ of the universe (Bartlett & Silk 1994, Oukbir & Blanchard 1992, 1997, Blanchard & Bartlett 1998). However, these procedures generally assume the cluster gas is spherical, unclumped and isothermal. Almost all clusters, however, show departures from these simplistic assumptions with some to a large extent.
Departures from these simple assumptions can lead to systematic errors in the determination of the different cosmological parameters (Inagaki et al 1995), especially, it was seen that non-isothermality of the cluster can lead to a substantial error in values of the cosmological parameters. Temperature structure in a cluster can be the result of the shape of the gravitational potential (Navarro et al 1997, Makino et al 1998), or it can arise due to the fact that the initial falling gas in the cluster potential is less shock heated than the later falling gas (Evrard 1990). In fact hydrodynamical simulations of isolated clusters also show a definite temperature structure and can introduce error in the value of $`H_{}`$ when compared to the traditional isothermal $`\beta `$-models (Yoshikawa 1998). Roettiger et al (1997) have shown that cluster mergers can result in deviations from both sphericity and isothermality. Observationally, the main handicap arises from the fact that the thermal structure of clusters are hard to measure, and the temperatures generally taken in analyses are the X-ray emission weighted temperature, usually measured over a few core radii. So, in general an isothermal description of the cluster is taken (or sometimes a phenomenological temperature model based on the Coma cluster: for example see Eq 73 of Birkinshaw 1999).
In this paper, we study another important phenomenon that can substantially change the temperature structure, viz. a cooling flow. Cooling flows in clusters of galaxies (for an introduction, see Fabian et al 1984) is a well established fact by now, and it is seen that around $`6090\%`$ of clusters exhibit cooling flows in their core with $`40\%`$ of them having cooling flows of more than $`100M_{}yr^1`$ (Markevitch et al 1998, Peres et al 1998, Allen et al 1999). In the largest systems, the mass deposition rate can be as high as $`1000M_{}yr^1`$ (Allen 2000). The idealised picture of a cooling flow is as following: Initially when the cluster forms, the infalling gas is heated from gravitational collapse. With time this gas cools slowly and a quasi hydrostatic state emerges. However in the central region, where energy is lost due to radiation faster than elsewhere, an inward ‘cooling flow’ initially arises due to the pressure gradient (Fabian 1994). This can modify the SZ decrement and act as a systematic source of error in the determination of the cosmological parameters.
Schlickeiser (1991) has shown that free-free emission from cold gas in the cooling flow can actually lead to an apparent decrease of the SZ effect at the centre. Since the central cooling flow region is generally very small, the isothermal $`\beta `$-model of cluster gas can still be used for the majority of cluster region even for cooling flow cluster, with the extra precaution of excluding the central X-ray spike from the X-ray fit, and a corresponding change made in the fitting of the SZ decrement. This is only possible for nearby clusters, however with well resolved cluster cores.
Naively, the change in the central SZ decrement $`y(0)`$ can be seen as follows : For a non cooling flow cluster, the central decrement is given by the line of sight integral of the electron pressure through the cluster centre along the full extent of the cluster. If the cluster has a maximum radius $`r_{cl}`$, then the central SZ decrement at RJ wavelengths can be written as $`y(0)=4\frac{\sigma _T}{m_ec^2}_0^{r_{cl}}p_e𝑑l`$. For a cluster with a cooling flow, let us suppose that the electron pressure $`p_e`$, drastically falls below a certain radius $`r_s`$, which is typically well inside the core of the cluster. The resulting central decrement is then $`y(0)4\frac{\sigma _T}{m_ec^2}_{r_s}^{r_{cl}}p_e𝑑l`$. Depending on the distance of $`r_s`$ from the cluster centre \[$`r_s(0.1`$ to $`0.3)r_{core}`$\],there wiil be a change in the value of $`y(0)`$ by $`5\%25\%`$. However, this simplistic view may not be true. This estimate assumes that the pressure profile remains a $`\beta `$ profile outside the radius $`r_s`$. The pressure profile, however, need not follow the $`\beta `$ profile once cooling flow starts and it can deviate from it substantially even for radii much larger than $`r_s`$. As a matter of fact, there can actually be an increase in the pressure for a large region inside the cooling flow, before a sudden drop inside $`r_s`$. Since the usual proecedure for estimating the Hubble constant depends on fitting $`\beta `$ profiles to the SZ and X-ray profiles, to estimate $`r_{core}`$, this change in the pressure profile due to cooling flow can distort the estimation of $`r_{core}`$ and hence, the value of $`H_{}`$ in a non-trivial way. We study this effect in detail in later sections.
In this paper, we have investigated the problem of cooling flow induced change in the temperature and density profile, its effect on the SZ effect, and its subsequent effect on the determination of cosmological parameters. In §2 we briefly review the physics of Sunyaev Zel’dovich effect; §3 is devoted to the physics of cooling flows and discussing the cooling flow solutions; in §4 we look at the effect of cooling flow solutions on SZ effect, both on the determination of $`H_{}`$ and $`\mathrm{\Omega }_{}`$; we conclude in §5 with a brief comment on how this work differs from other work and the relevance of this paper. For the SZ effect, our notation and approach mainly follows that described in Barbosa et al(1996).
## 2. Determining Hubble constant with Sunyaev-Zel’dovich effect
### 2.1. The Sunyaev-Zel’dovich Effect
The integral of the electron pressure along any line-of-sight through the cluster determines the magnitude of the distortion of the apparent brightness temperature of the cosmic microwave background (CMB) due to SZ effect. This is quantified in terms of the Compton $`y`$-parameter:
$$y=𝑑l\frac{k_BT_e}{m_ec^2}n_e\sigma _T,$$
(1)
where $`k_B`$ is the Boltzman constant, $`T_e`$ is the gas temperature, $`m_e`$ is the electron rest mass, $`n_e`$ is the electron number density, $`c`$ is the velocity of light and $`\sigma _T`$ is the Thomson scattering cross-section. This occurs through the inverse-Compton scattering, by the hot intracluster gas, of the CMB photons propagating through the cluster medium, and the energy transfer in this interaction between hot electrons and CMB photons resulting in a distortion to the CMB spectrum. The SZ surface brightness at a position $`\theta `$ of the cluster with respect to the mean CMB intensity is given by
$$\delta \text{i}_\nu (\theta )=\text{y}(\theta )\text{j}_\nu (x),$$
(2)
x is a dimensionless frequency parameter
$$x=\frac{h\nu }{kT_o},$$
(3)
where $`h`$ is the Planck constant, $`\nu `$ is the observing frquency and $`T_0`$ is the CMB temperature at the present epoch: $`T_02.73`$K. The function $`j_\nu (x)`$ describes the spectral shape of the effect
$$\text{j}_\nu (x)=\frac{2(kT_{})^3}{(hc)^2}\frac{x^4e^x}{\left(e^x1\right)^2}\left[\frac{x}{\mathrm{tanh}\left(x/2\right)}4\right].$$
(4)
Since the total photon number is conserved in the inverse Compton scattering process, upscattering of the photons, the spectral dependance gets an unique shape, through a decrement in the brightness temperature at lower frequencies while an increase is observed at higher frequencies.
The Sunyaev-Zel’dovich effect provides a unique observing approach to traditional methods — which use X-ray temperature and X-ray luminosity. The X-ray studies have a major disadvantage because of ‘cosmological dimming’, the surface brightness of distant X-ray sources falls off as $`(1+z)^4`$, and for this reason, obtaining samples of clusters at cosmological distances is challenging. The SZ effect has the distinct advantage of being independent of the distance to the cluster. The SZ flux density from a cluster will diminish with distance to the cluster as the square of the angular-size distance; in contrast to X-ray flux densities from clusters which diminish as the square of the luminosity distance to the cluster. If the observations resolve clusters, particularly at lower redshifts, the observed sky SZ temperature distribution will be sensitive to the thermal electron temperature structure within the clusters; once again this may be contrasted with X-ray emission images of cluster gas distributions which are mainly sensitive to the gas density distribution.
### 2.2. Determination of Hubble constant
The method for the determination of Hubble constant using Sunyaev-Zel’dovich effect uses two observable quantities : 1) $`\mathrm{\Delta }T/T`$ of the CMB due to SZ effect; 2) the X-ray surface brightness $`S_X`$ of the cluster. These can be written as
$$\frac{\mathrm{\Delta }T_{SZ}}{T}(r)=2_{l_{min}}^{l_{max}}\frac{k_BT_e}{m_ec^2}\sigma _Tn_e𝑑l,$$
(5)
$$S_X(r)=\frac{1}{4\pi (1+z)^4}_{l_{min}}^{l_{max}}\frac{dL_X}{dV}𝑑l,$$
(6)
where $`r`$ is the distance to the line of sight from the cluster centre, $`l_{max}`$ and $`l_{min}`$ give the extension of the cluster along the line of sight, $`\frac{dL_X}{dV}`$ is the X-ray emissivity and $`dl`$ the line element along the line of sight. The X-ray emissivity in the frequency band $`\nu =\nu _1`$ to $`\nu _2`$ can be written as
$$\frac{dL_X}{dV}=n_{e}^{}{}_{}{}^{2}\alpha (T_e;\nu _1,\nu _2,z),$$
(7)
where
$$\alpha (T_e;\nu _1,\nu _2,z)=\frac{2}{1+X}\left[\frac{2\pi }{3m_ec^2}\right]^{1/2}\frac{16e^6}{3\mathrm{}m_ec^2}A(T_e;\nu _1,\nu _2,z),$$
(8)
where
$`A(T_e;\nu _1,\nu _2,z)={\displaystyle _{u_1(1+z)}^{u_2(1+z)}}(k_BT_e)^{1/2}e^u`$
$`[Xg_{ff}(T_e,u,1)+(1X)g_{ff}(T_e,u,2)]du.`$ (9)
In the above equations we have assumed primordial abundance of hydrogen and helium and have set $`X=0.76`$, $`e`$ is the electron charge, $`\mathrm{}=h/(2\pi )`$, $`u2\pi \mathrm{}\nu /k_Bt_e`$, and $`g_{ff}(T_e,u,Z)`$ is the velocity averaged Gaunt factor for the ion of charge Ze (Kellog 1975). Traditionally, to model the cluster gas distribution one takes the following density and temperature profiles (Cavaliere & Fusco-Femiano 1978)
$$n_e(r)=n_{e0}\left[1+\left(\frac{r}{r_{core}}\right)^2\right]^{3\beta /2},$$
(10)
$$T_e(r)=T_{iso}=constant,$$
(11)
where $`n_{e0}`$ is the central electron density and $`r_{core}`$ the core radius of the cluster. The above expressions are used as an empirical fitting model, and the parameter $`\mathrm{`}\beta ^{}`$ is regarded as the fitting parameter. The equation holds for $`0<r<R_{cluster}`$, where $`R_{cluster}`$ is the maximum ‘effective’ extension of the cluster. Conventionally, $`R_{cluster}=\mathrm{}`$, and then from equations (5),(6),(10),(11), we get
$`{\displaystyle \frac{\mathrm{\Delta }T_{SZ}}{T}}(\theta )`$ $`=`$ $`{\displaystyle \frac{2\sqrt{\pi }\sigma _Tk_BT_{iso}}{m_ec^2}}n_{e0}r_{core}{\displaystyle \frac{\mathrm{\Gamma }(3\beta /21/2)}{\mathrm{\Gamma }(3\beta /2)}}`$ (12)
$`\times \left[1+({\displaystyle \frac{d_A\theta }{r_{core}}})^2\right]^{1/23\beta /2},`$
$`S_X(\theta )`$ $`=`$ $`{\displaystyle \frac{\sqrt{\pi }}{4\pi (1+z)^4}}\alpha n_{e0}^2r_{core}{\displaystyle \frac{\mathrm{\Gamma }(3\beta 1/2)}{\mathrm{\Gamma }(3\beta )}}`$ (13)
$`\times \left[1+({\displaystyle \frac{d_A\theta }{r_{core}}})^2\right]^{1/23\beta },`$
where $`\mathrm{\Gamma }(x)`$ is the gamma function. Since both the central CMB decrement and the X-ray surface brightness are observed, one can then combine equations (12) and (13) to estimate the core radius as
$`r_{c,est}`$ $`=`$ $`{\displaystyle \frac{\left[\frac{\mathrm{\Delta }T(\theta )}{T}\right]_{obs}^2}{S_X(\theta )_{obs}}}{\displaystyle \frac{\mathrm{\Gamma }(3\beta 1/2)\mathrm{\Gamma }(3\beta /2)^2}{\mathrm{\Gamma }(3\beta /21/2)^2\mathrm{\Gamma }(3\beta )}}`$
$`{\displaystyle \frac{m_e^2c^4\alpha }{16\pi ^{3/2}(1+z)^4\sigma _Tk_BT_{e,fit}^2}}\left[1+({\displaystyle \frac{\theta }{\theta _{X,core}}})^2\right]^{1/2}.`$
In the above equation $`\theta _{X,core}`$ is the angular core radius observed in X-ray, and $`T_{e,fit}`$ is the X-ray flux averaged temperature (obtained from fitting the observed X-ray spectrum to the theoretical spectrum expected from isothermal case). This X-ray emission weighted temperature is given by
$$T_{iso}\frac{_0^{r_{vir}}T_e(r)\alpha n_{e}^{}{}_{}{}^{2}r^2𝑑r}{_0^{r_{vir}}\alpha n_{e}^{}{}_{}{}^{2}r^2𝑑r}.$$
(15)
The point to be noted is that $`r_{vir}`$ is the virial radius of the cluster and its choice depends on the observer. If the temperature has a spatial structure then the $`T_{iso}`$ inferred from such a procedure may give different values depending on how much of the cluster is taken in making the above average. It has been seen (Yoshikawa et al 1998), that this can lead to a substantial change in SZ effect inferred, and thus to the value $`H_{}`$.
The angular diameter distance $`d_A`$ can be approximated for nearby ($`z1`$) clusters as
$$d_A=\frac{cz}{H_{}}\left[1+\frac{2\mathrm{\Lambda }\mathrm{\Omega }_{}6}{4}z+O(z^2)\right].$$
(16)
Thus finally we have got an estimate of $`d_A(z)`$ as $`r_{c,est}/\theta _{X,core}`$. Now, if from other observations we know the cosmological parameters $`\mathrm{\Omega }_{}`$ and $`\mathrm{\Lambda }`$, then one can estimate the Hubble constant
As can be seen from the above equations, the value of $`H_{}`$ depends crucially on the many assumptions of isothermality and $`\beta `$-model density distribution of the cluster. Cooling flow changes both of these and so it can significantly affect the value of the Hubble constant.
## 3. Cooling flows in clusters
### 3.1. Preliminaries
From X-ray spectra of clusters, it is known that the continuum emission is thermal Bremsstrahlung in nature and originates from diffuse intracluster gas with densities $`10^210^4cm^3`$ and temperatures around $`10^710^8`$ K. The gas is usually believed to be in hydrostatic equilibrium. If, however, the density in the inner region is large enough, so that the cooling time is less than the age of the cluster, then there is a ‘cooling flow’ (Fabian et al 1984 and references therein). Of course there would be a flow only when the dynamical time is also shorter than the cooling time scale ($`t_{age}>t_{cool}>t_{dyn}`$).
The basic equations of cooling flow are:
$`{\displaystyle \frac{du}{dr}}`$ $`=`$ $`{\displaystyle \frac{u}{\rho _b}}{\displaystyle \frac{d\rho _b}{dr}}{\displaystyle \frac{2u}{r}}{\displaystyle \frac{\dot{\rho _b}}{\rho _b}},`$
$`u{\displaystyle \frac{du}{dr}}+{\displaystyle \frac{1}{\rho _b}}{\displaystyle \frac{d(\rho _b\theta )}{dr}}`$ $`=`$ $`{\displaystyle \frac{GM_t(r)}{r^2}},`$
$`u{\displaystyle \frac{d}{dr}}\left({\displaystyle \frac{3\theta }{2}}\right){\displaystyle \frac{\theta u}{\rho _b}}{\displaystyle \frac{d\rho _b}{dr}}`$ $`=`$ $`{\displaystyle \frac{\rho \mathrm{\Lambda }(\theta )}{(\mu m_p)^2}},`$ (17)
where $`\theta =2k_BT/\mu m_p`$, $`\mu `$ is the mean molecular weight, and $`m_p`$ is the proton mass. For steady flows with constant mass flux, $`\dot{\rho }=0`$. This implies $`u=\dot{m}/4\pi \rho r^2`$ for steady flows. (Note that in cooling flows both $`u`$ and $`\dot{m}`$ are negative. The subscripts $`b`$ refers to baryons and $`t`$ refers to total i.e baryons + dark matter. However, we assume the baryonic contribution to the total mass negligible w.r.t to the dark matter contribution.)
$`M(r)`$ describes the distribution of the total mass and depends on the details of dark matter density profiles (see below). $`\mathrm{\Lambda }(\theta )`$ is the cooling function defined so that $`n_en_p\mathrm{\Lambda }(\theta )`$ is the rate of cooling per unit volume. We use an analytical fit to the optically thin cooling function as given by Sarazin & White (1987),
$`\left({\displaystyle \frac{\mathrm{\Lambda }(\theta )}{10^{22}\mathrm{ergcm}^3\mathrm{s}^1}}\right)=4.7\times \mathrm{exp}\left[\left({\displaystyle \frac{T}{3.5\times 10^5\mathrm{K}}}\right)^{4.5}\right]`$ (18)
$`+`$ $`0.313\times T^{0.08}\mathrm{exp}\left[\left({\displaystyle \frac{T}{3.0\times 10^6\mathrm{K}}}\right)^{4.4}\right]`$
$`+`$ $`6.42\times T^{0.2}\mathrm{exp}\left[\left({\displaystyle \frac{T}{2.1\times 10^7\mathrm{K}}}\right)^{4.0}\right]`$
$`+`$ $`0.000439\times T^{0.35}.`$
This fit is accurate to within $`4`$% accuracy, for a plasma with solar metallicity, within $`10^5T10^8`$ K. For $`10^8T<10^9`$ K, it underestimates cooling by a factor of order unity (compared to the exact cooling function , as in, e.g., Schmutzler & Tscharnuter, 1993) , and therefore is a conservative fit to use, as far as the effect of cooling is concerned.
For non-steady flows, we adopt the formalism of White & Sarazin (1987), where the mass deposition rate, $`\dot{\rho }`$, is characterised by a ‘gas-loss efficiency’ parameter $`q`$. One writes $`\dot{\rho }=q(\rho /t_{cool})`$ where $`t_{cool}`$ is the local isobaric cooling rate ($`t_{cool}=5k_BT\mu m_p/\rho \mathrm{\Lambda }`$). It has been found that $`q3`$ models can reproduce the observed variation of mass flux ($`\dot{m}r`$) (Sarazin & Graney 1991). Fabian (1994) has noted that these models of White & Sarazin (1987) yield good approximations to the emission weighted mean temperature and density profiles for cooling flow clusters. We also note that Rizza et al (2000) have used the steady flow models of White & Sarazin (1987) to simulate cooling flows.
We first discuss cooling flows with $`\dot{m}=`$constant. With $`q=0`$, one can eliminate the density from Eq. 17 to get two differential equations:
$`{\displaystyle \frac{du}{dr}}`$ $`=`$ $`{\displaystyle \frac{u}{\left[r^2(5\theta 3u^2)\right]}}\left[3GM10r\theta +{\displaystyle \frac{\dot{m}}{2\pi }}{\displaystyle \frac{\mathrm{\Lambda }(\theta )}{uM^2}}\right]`$
$`{\displaystyle \frac{d\theta }{dr}}`$ $`=`$ $`{\displaystyle \frac{2}{\left[r^2(5\theta 3u^2)\right]}}\left[\theta (2u^2rGM)(u^2\theta ){\displaystyle \frac{\dot{m}}{4\pi }}{\displaystyle \frac{\mathrm{\Lambda }(\theta )}{uM^2}}\right]`$ (19)
These equations have singularities at the sonic radius $`r_s`$ where $`5\theta _s=3u_s`$. A necessary condition of singularity is that the numerators of Eq. 19 vanish at the sonic radius. Therefore (Mathews & Bregman 1978)
$$r_s=(3/10\theta _s)\left[GM(r)+\frac{\dot{m}\mathrm{\Lambda }(\theta _s)}{10\pi \theta _sM^2}\right]$$
(20)
We have used two different dark matter profiles for the cluster. The first model ( Model A) has been discussed earlier in the literature in the context of cooling flows in cluster (White & Sarazin 1987; Wise & Sarazin 1993) with a density profile,
$$\rho _d=\{\begin{array}{cc}\frac{\rho _o}{1+(r/r_{core})^2}+\frac{\rho _{o,g}}{1+(r/r_{c,g})^2}\hfill & \text{if }r<237\text{ kpc}\hfill \\ \frac{\rho _o}{1+(r/r_{core})^2}\hfill & \text{if }r>237\text{ kpc}\hfill \end{array}$$
(21)
Here $`\rho _o=1.8\times 10^{25}`$ gm cm<sup>-3</sup> and $`r_{core}=250`$ kpc describe the profile of the cluster mass, and $`\rho _{o,g}=4.1\times 10^{22}`$ gm cm<sup>-3</sup> and $`r_{c,g}=1.69`$ kpc describe the profile of the galaxy in the centre of the cluster.
Model B does not have the galaxy in the center, and so it is described simply by $`\rho =\rho _0/[1+(r/r_{core})^2]`$.
With these assumptions, the solutions for steady cooling flows, $`\dot{m}=`$ constant) are fully characterized by (1) the inflow rate, $`\dot{m}`$, and (2) the temperature of the gas $`T_s`$ at the sonic radius $`r_s`$. Obviously, the cooling flow solutions are only valid within the cooling radius $`r_{cool}`$ where $`t_{cool}=t_{age}`$. We assume a value of $`t_{age}=10`$ Gyr for all models. We assume that outside the cooling radius, gas obeys quasi-hydrostatic equilibrium (Sarazin 1986). Although this means matching the cooling flow solutions with nonzero $`u`$ to $`u=0`$ solutions outside, in reality the velocity of gas at the cooling radius is very small (for a $`\dot{m}=100`$ M yr<sup>-1</sup> with $`\rho 10^{26}`$ gm cm<sup>-3</sup>, at $`r=250`$ kpc implies a velocity of $`30`$ km s<sup>-1</sup>), which is close to the limit of turbulence in the cluster gas (Jaffe 1980), and smaller than the sound velocity ($`1.5\times 10^3(T/10^8K)^{1/2}`$ km s<sup>-1</sup>). The velocity of the flow at the cooling radius is, therefore, for all practical purposes, sufficiently small to be matched to the solution of hydrostatic equilibrium outside. (In this approach, we avoid the time consuming search for the critical value of $`\dot{m}`$ for which the flow solutions behave isothermally at $`r\mathrm{}`$ (see Sulkanen et al. 1993).)
As in the usual assumptions for the interpretation of SZ effect, we assume that the gas outside the cooling radius is isothermal, with a constant temperature profile. The density, therefore, obeys $`\rho [1+(r/r_{core})^2]^{3\beta /2}`$, where $`\beta =\mu m_p\sigma ^2/k_BT_{iso}`$, and $`T_{iso}`$ is the temperature of the gas at and outside the cooling radius.
For models with non-zero $`q`$ (Model C: has the same mass profiles as Model A), the solutions are characterized by $`T_s`$ and the value of $`\dot{m}`$ at the cooling radius, $`\dot{m}_{cool}`$. Since a fraction of mass drops out of the flow in this case, the inflow velocity need not rise fast and so it is possible to find completely subsonic solutions.
### 3.2. Cooling flow solutions
We numerically solved the flow equations for the parameters listed in Table 1. The density, temperature and pressure profiles for three cases are presented in Figures 1, 2 & 3 . We mark the position of $`r_{cool}`$ in each case, and we mark $`r_s`$ for the cases of transonic flows (when $`\dot{m}=`$ constant). Beyond $`r_{cool}`$ we match a hydrostatic solution, as explained above, for the respective potentials. We also present, for comparison, the behaviour if the solutions outside $`r_{cool}`$ are assumed to extend inwards (that is, if no cooling flow is assumed). We will postpone the discussion on the effect of these profiles on the SZ decrement to a later section, and only discuss the qualitative aspects of the solutions here.
The solution A1 is similar to that presented by Wise & Sarazin (1993) (their Figure 1; although they chose to characterise the solutions by the temperature at $`r_{cool}`$ and not $`T_s`$ as we have done here). It is also similar (qualitatively) to the solution for A262 presented by Sulkanen et al. (1989). As the latter authors have noted, the effect of having a galactic potential in the center is to have a flatter temperature profile for $`r>r_s`$, than in the case of no galactic potential. This aspect is clearly seen while comparing our solutions with and without galactic potentials in the center. Our calculations for the case without the central galaxy are admittedly flawed in the very inner regions where the gas density is larger than the dark matter density, which results in an incorrect determination of the graviational potential in the inner region. However, this happens only inside a region $`25`$ kpc from the centre, and should not influence our final results by a large extent.
A word of explanation for the pressure profiles is in order here. Naively speaking, it would appear that the pressure profile inside the cooling radius should have lower values than the corresponding case of hydrostatic equilibrium. The fact that it is not always so has been noted in the literature (e.g., Soker & Sarazin 1988, Fig 1 of Sulkanen et al. 1989 ). The reason for the pressure bump just outside of the sonic radius is that the flow in this inner region is not pressure driven, but rather by gravity (see also Soker & Sarazin 1988). This is why the bump in the profile depends on the presence and absence of the galaxy in the centre. And this profile leads to the curious result that the presence of cooling flow can lead to the overestimation of the Hubble constant as discussed in the next section
The model with mass deposition (C1) is shown in Figure 3. The local mass flux is found to be almost proportional to the radius, consistent with various observations (Fabian, Nulsen & Canizares 1984; Thomas et al. 1987), and, therefore, is probably a realistic model for cooling flow clusters. In this case the temperature drops gradually all the way through, since the velocity does not rise too fast. The deposited mass is assumed to impart negligible pressure and the pressure refers only to that of gas taking part in the flow.
## 4. Determination of Hubble constant
In this section we discuss the SZ and X-ray profiles of clusters with cooling flows. We compare these with profiles from cluster having gas in hydrostatic equilibrium, and comment on the reliability of measuring Hubble constant. The effect of cooling flow and the subsequent increased Bremsstrahlung emmison is seen in the sudden increase in the X-ray flux in the innermost region of the cluster (Figure 5). The signature of the cooling flow is seen in the form of the central spike in the X-ray profile. The X-ray profile is only affected slightly by the drop in temperature and it is the dependence on the gas density that holds . The temperature dependance becomes important only near the sonic point. Outside $`r_{cool}`$, the X-ray profile is the same as that in the hydrostatic cases.
The SZ distortion is proportional to the line of sight integral of the pressure, and the sudden increase of the gas density inside the cooling radius is moderated by the decrease of the gas temperature. As a result there is a gradual increase in the gas pressure. Near the sonic point the temperature drops drastically by orders of magnitude, and results in sudden decrease in pressure. However, since this change in pressure becomes acute only within $`5\%`$ of the core radius, it contributes negligibly to the line-of-sight integral of the gas pressure, and leads to an increase in the SZ distortion inside the cooling radius for all models considered (see Figure 4). Like the X-ray profiles, the SZ profiles outside $`r_{cool}`$ is the same as that for the corresponding hydrosstatic cases.
The SZ profiles have been calculated in the Rayleigh-Jeans limit ($`x1`$) where $`j_\nu (x)`$ of Eqn. 4 goes to $`2`$. In general, however, the profiles should be calculated using Eqn 2. Our results below are independent of the observational frequency, since the profiles at different frequencies have similar shapes, with the amplitude of the SZ distortion scaled either up or down.
Once both profiles are known, one can determine the deviation in the value of the Hubble constant using Eqns(14) & (16). The deviation from the idealistic case can be parametrised as
$$f_H\frac{r_{core,true}}{r_{core,est,fit}}=\frac{H_{0,est}}{H_{0,true}}$$
(22)
The above formula has been used to determine the deviation of the estimated value of $`H_{}`$ from the actual value, for models listed in Table-1. The effect of cooling flow on the determination of the cosmological parameters are summarised in Table-2.
To begin with, one has to get best fitted values for $`r_{core}`$ (or $`\theta _c`$) and $`\beta `$ from different profiles. Since, the estimation of the Hubble constant depends on the determination of these parameters from the profiles, we look at this issue in more detail. We must keep in mind that the best fitted value of $`r_{core}`$ (or $`\theta _c`$) and $`\beta `$ depends on whether one decides to fit the X-ray or the SZ profiles, and the choice can lead to significant differences in the estimated value of $`H_{}`$. One of the reasons for the strong dependance on the nature of the profile can be the non-isothermality of the cluster gas. Recent observations indicate that intracluster gas has a temperature structure, see Markevitch et al 1998. This is because the y-parameter depends on the integral over $`T_e`$, while emissivity of thermal Bremsstrahlung depends on $`\sqrt{T_e}`$. The dependance of the Gaunt factor on $`T_e`$ isindirect and weak. Yoshikawa et al 1998 have shown that gas temperature drop in the central regions (their Fig. 3), should increase both $`r_{core}`$ and $`\beta `$ fitted to $`y(\theta )`$, and to alesser extent to $`S_X(\theta )`$, as compared to those compared to $`n_e(r)`$. This discrepancy increases at higher redshifts. However, in their case, there is little change in the gas density profile. Clumpiness can also give rise to different fits, resulting in an overestimation of the Hubble constant (Inagaki etal 1995).
There are two other important points that have to be kept in mind while fitting the profiles. First, we must remember that we are trying to fit a cluster having a finite profile with the formulae (Eqns 12 & 13) for isothermal $`\beta `$ profiles which is derived assuming the cluster to be of infinite extent. This can, by itself, lead to an overestimation of $`H_{}`$ (Inagaki et al 1995). Thus to have a good fit one must choose a segment of the profile such that, within that segment, the profiles (SZ or X-ray) for a finite cluster do not differ much from those of a hypothetical cluster of infinite extent. We found that SZ and X-ray profiles of clusters start differing from those of infinite size at radii greater than $`1.5`$ times the core radius. Hence, we have restricted our fitting to radii within $`1.5r_{core}`$.
Next, one must also be careful to exclude the region close to the sonic point, so that the X-ray spike is excluded from the fit. Also, the central portion in the SZ profile should be avoided as its inclusion can give an apparent central decrement less than its neighbouring points (see Schlickeiser 1991). We have fitted the SZ and X-ray profiles varying the inner cutoff radius and the results for a representative solution for each class of model are tabulated in Tables 3 & 4. Thus, all fittings were done for profiles extending from $`r=r_{min}`$ to $`r=1.5r_{core}`$.
As can be seen from Table-2, cooling flows can lead to an overestimation of the Hubble constant. However, we must emphasise, that it may not be possible to a priori estimate the amount of bias introduced in the measurement of the Hubble constant due to cooling flows. There is no simple correlation between the amount of cooling (i.e $`\dot{m}`$) and the change in the estimated $`H_{}`$ from the actual value. The total change depends not only on $`\dot{m}`$, but also on the position of cooling radius, sonic radius, temperature at the sonic point and the isothermal temperature characterising the hydrostatic cases, with which comparisons are made. Specifically, the fitted values of $`r_{core}(\theta _c)`$ and $`\beta `$ for cooling flow models differ from hydrostatic models according to shape of underlying profiles, which is marked by two important features, firstly the central excess of X-ray flux (or excess decrement of SZ flux), and secondly the the deviation from the smooth hydrostatic profile inside $`r_{cool}`$, the amount of overestimation mainly depends on these factors. For models with a central galaxy potential, there is always an over-estimation of $`H_{}`$, which is greater than the models without the central galaxy.
For the realistic cases of models C1 and C2, where we have a variable $`\dot{m}`$ with $`r`$ inside the cooling radius, the deviation of estimated Hubble constant from its actual value is almost the same. They are also greater than that of models A and B, having similar mass flow rates. This may be due to the fact that the maximum deviation in pressure from the hydrostatic cases is more in non steady cases, than in steady flows. Also non-steady cases are marked by the absence of the sonic radii and the subsequent drop in temperature.
We note that although the different choice of fitting may change the absolute determination of cosmological parameters, the trend i.e deviation from the correct values remains more or less unaffected. It is interesting to note that for B-type model (C1), which include mass deposition in cooling flows, the deviations decrease as one excludes a greater part of the cooling flow region (Table 4& 5). The other models show an increase instead. Here, we remind ourselves that models with mass deposition i.e C-type models are more realistic (Fabian 1994). It is possible that the unsually high value of deviation (Table 4) and the counter-intuitive trend of increasing deviation with decreasing portion of cooling flow region used for fitting (Table 4 & 5), arise because of the unrealistic modelling of cooling flows. If we take the model B1 as a realistic one, then Table 4 & 5 show that to obtain a value of the Hubble constant within an accuracy of $`10\%`$, one should have $`r_{min}0.8r_{cool}`$. In most cases, $`r_{cool}<r_{core}`$ (Fabian et al. 1984). However, as $`r_{cool}`$ cannot be determined without actually detecting a cooling flow in a cluster, we suggest that a significant portion of the profile within $`r_{core}`$ should be excluded as a precaution. The SZ and the X-ray profiles for the different models are shown in Figures 4 & 5.
## 5. Discussions and Conclusion
Our work on the effect of temperature structure of clusters and its effect on SZ decrement differs from other previous work of this nature in following way: this work takes into account the change in density profile as well as the temperature profile since both becomes important in the central region of the cluster. Also, previous authors have looked at the issue of non-isothermality of a cluster at radii greater than the core radius of the cluster, whereas we look at temperature change at regions inside the core radius. For them the density profile can still be well approximated by a $`\beta `$ profile, whereas for cooling flow solutions, density profile is vastly different. Further, they have neglected radiative cooling in their work. We for the first time look at SZ effect in presence of radiative cooling, by first solving the cooling flow equations for reasonable and physical solutions.
In summary, we find that the presence of a cooling flow in a cluster can lead to an overestimation of the Hubble constant determined from the Sunyaev-Zel’dovich decrement. We have used different models of cooling flows, with and without mass deposition, and found the deviation in the estimated value of the Hubble constant in the case of a cooling flow from that of hydrostatic equilibrium. We have used the usual procedure of fitting the SZ and X-ray profile with a $`\beta `$ profile to get an estimated value of $`r_{core}`$, and then compared with that for the case of gas in hydrostatic equilibrium in order to estimate the deviation in the Hubble constant. For the more realistic models with mass deposition (varying $`\dot{m}`$ with radius), we found that the deviation decreases with the exclusion of greater portions of the cooling flow region. Quantitatively, we found that for the deviation to be less than $`10\%`$, one should exclude a portion of the profile upto $`0.8r_{cool}`$. Since $`r_{cool}`$ is difficult to estimate without actually detecting a cooling flow, we have suggested that a significant portion of the profile inside $`r_{core}`$ should be excluded, to be safe.
There can be another important implication of the effect of cooling flows. With the upcoming satellite missions (MAP & Planck), we have come to the point where there are efforts to constrain $`\mathrm{\Omega }_0`$ with surveys of blank SZ fields (Bartlett et al 1998; Bartlett 2000), ultimately giving rise to SZ-selected catalog of clusters (Aghanim et al 1997). This method relies on estimating the number of SZ sources brighter than a given threshold flux (Barbosa et al 1996). The point to be noted is that since these surveys are essentially flux limited in nature, the validity of the analysis to determine $`\mathrm{\Omega }_0`$ depends crucially on the one-to-one association of flux-limits to corresponding mass limits of clusters. From our analysis above, it seems that it may not be possible to associate a unique cluster mass to a given SZ distortion, given the uncertainty due to the presence of cooling flows. This might lead to contamination in SZ cluster catalogs and the inference of $`\mathrm{\Omega }_0`$. Recently, attempts have been made to constrain $`\mathrm{\Omega }_{}`$ from variance measurement of brightness temperature in blank fields (Subrahmanyan et al 1998) and comparing them to simulated fields (Majumdar & Subrahmanyan, 2000), of cumulative SZ distortions from a cosmolgical distribution of clusters. These results may also be systematically affected due to the presence of clusters having cooling flows.
The estimations made in this paper strictly applies to cases where the image of the SZ effect is directly obtained by single dish observations. For interferometric observations, the interferometer samples the fourier transform of the sky brightness rather than the direct image of the sky. The fourier conjugate variables to the right ascension and declination form the $`uv`$ plane in the fourier domain. Due to spatial filtering by an interferometer, it is necessary that models be fitted directly to the data in the $`uv`$ plane, rather that to the image after deconvolution. We do not forsee drastic change from our inferences in such cases since the result mainly depends on the deviation of the SZ and X-ray profile in case of a cooling flow from those in hydrostatic equilibrium. This, however, should be looked in greater detail in future. We also note that with the growing number of high quality images of the SZ effect with interferometers, which have greater resolution than single dish antennas, the shape parameters of the clusters can be directly determined from the SZ dataset rather than from an X-ray image (Grego et al. 2000).
Finally, we would like to add, though the calculations presented in this paper were done using the dark matter profile (Eqn. 21), which is ”commonly used” for calculating cooling flow solutions, it is inconsistent with the dark matter profile (Navarro et al 1997) found in numerical simulations. (For a comparison of mass and gas distribution in clusters having cooling flows with different dark matter profiles, see Waxman & Miralda-Escude’, 1995). Moreover, we have neglected the self gravity of the gas. Suto et al (1998) have calculated the effect of including the self gravity of the gas in determining the gas density profile. To make strong conclusions about the effect of cooling flows in the determination of the Hubble constant, one should take both the points mentioned above into account.
We are grateful to Joseph Silk for his comments on the manuscript. We also thank the referee, Dr Asantha Cooray, for numerous suggestions which have helped us in improving the paper. One of the authors (SM) would like to thank the Raman Research Institute, Bangalore, for hospitality and for providing computational facilities. SM acknowledges help recieved from R.Sridharan on using IDL. Finally, he would also like to express his gratitude to Pijushpani Bhattacharjee for his constant encouragement.
Table 1 – Parameters for cooling flow solutions
| Solution | Mass Model | q | $`\dot{m}\left(r_{cool}\right)`$ | $`T_s`$ | $`r_s`$ | $`r_{cool}`$ | $`T_{iso}`$ |
| --- | --- | --- | --- | --- | --- | --- | --- |
| | | | (M yr<sup>-1</sup>) | (K) | (kpc) | (kpc) | (K) |
| A1 | Model A | 0 | 100 | $`6.5\times 10^6`$ | $`0.688`$ | $`127.5`$ | $`1.2\times 10^8`$ |
| A2 | Model A | 0 | 200 | $`6.5\times 10^6`$ | $`0.462`$ | $`96.1`$ | $`7.7\times 10^8`$ |
| A3 | Model A | 0 | 300 | $`6.5\times 10^6`$ | $`0.712`$ | $`132.2`$ | $`7.7\times 10^8`$ |
| B1 | Model B | 0 | 100 | $`4.0\times 10^6`$ | $`0.688`$ | $`85.7`$ | $`1.14\times 10^9`$ |
| B2 | Model B | 0 | 200 | $`6.5\times 10^6`$ | $`0.462`$ | $`89.6`$ | $`1.9\times 10^9`$ |
| B3 | Model B | 0 | 300 | $`6.5\times 10^6`$ | $`0.712`$ | $`110.3`$ | $`1.9\times 10^9`$ |
| C1 | Model A | 3 | 200 | | | $`111.6`$ | $`1.1\times 10^8`$ |
| C2 | Model A | 3 | 300 | | | $`132.2`$ | $`1.1\times 10^8`$ |
Table 2 – Effect on central decrement and $`H_{}`$ for $`r_{min}=0.1r_{core}`$
| Solution Type | $`\frac{H_{est}}{H_{true}}`$ | $`\mathrm{\Delta }y_0`$ |
| --- | --- | --- |
| | | (% change) |
| A1 | $`1.91`$ | $`35`$ |
| A2 | $`1.18`$ | $`11.5`$ |
| A3 | $`2.6`$ | $`25`$ |
| B1 | $`1.36`$ | $`12.0`$ |
| B2 | $`1.19`$ | $`9.0`$ |
| B3 | $`1.13`$ | $`8.5`$ |
| C1 | $`2.6`$ | $`14.0`$ |
| C2 | $`2.7`$ | $`11.3`$ |
Table 3 – Fitting of SZ profile and deviation of $`H_{}`$ for $`\dot{m}=200m_{}/yr`$
| Solution Type | $`r_{min}=0.2r_{cool}`$ | $`r_{min}=0.5r_{cool}`$ | $`r_{min}=0.8r_{cool}`$ |
| --- | --- | --- | --- |
| | $`H/H_{true}`$ | $`H/H_{true}`$ | $`H/H_{true}`$ |
| A2 | $`1.57`$ | $`1.74`$ | $`2.18`$ |
| B2 | $`1.44`$ | $`1.58`$ | $`2.06`$ |
| C1 | $`2.20`$ | $`1.18`$ | $`1.07`$ |
Table 4 – Fitting of X-ray profile and deviation of $`H_{}`$ for $`\dot{m}=200m_{}/yr`$
| Solution Type | $`r_{min}=0.5r_{cool}`$ | $`r_{min}=0.8r_{cool}`$ | $`r_{min}=0.9r_{cool}`$ | $`r_{min}=0.95r_{cool}`$ |
| --- | --- | --- | --- | --- |
| | $`H/H_{true}`$ | $`H/H_{true}`$ | $`H/H_{true}`$ | $`H/H_{true}`$ |
| A2 | | $`4.7`$ | $`2.7`$ | $`1.7`$ |
| B2 | | $`4.9`$ | $`2.24`$ | $`1.6`$ |
| C1 | $`1.69`$ | $`1.12`$ | $`1.02`$ | $`1.0`$ |
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# PHENOMENOLOGY OF NEUTRINO OSCILLATIONS AT A NEUTRINO FACTORY
## 1 Introduction
There have been several experiments $`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$ which suggest neutrino oscillations. It has been known in the two flavor framework that the solar neutrino deficit can be explained by neutrino oscillation with the set of parameters $`(\mathrm{\Delta }m_{}^2,\mathrm{sin}^22\theta _{})`$ $`(𝒪(10^5\mathrm{eV}^2),𝒪(10^2))`$ (small angle MSW solution), $`(𝒪(10^5\mathrm{eV}^2),𝒪(1))`$ (large angle MSW solution), or $`(𝒪(10^{10}\mathrm{eV}^2),𝒪(1))`$ (vacuum oscillation solution), and the atmospheric neutrino anomaly can be accounted for by $`(\mathrm{\Delta }m_{\text{atm}}^2,\mathrm{sin}^22\theta _{\text{atm}})(10^{2.5}\mathrm{eV}^2,1.0)`$. In the three flavor framework there are two independent mass squared differences and it is usually assumed that these two mass differences correspond to $`\mathrm{\Delta }m_{}^2`$ and $`\mathrm{\Delta }m_{\text{atm}}^2`$. Throughout this talk I will assume three neutrino species which can account for only the solar neutrino deficit and the atmospheric neutrino anomaly <sup>a</sup><sup>a</sup>aTo explain the LSND anomaly $`^\mathrm{?}`$ one needs at least four neutrino species.. Without loss of generality I assume $`|\mathrm{\Delta }m_{21}^2|<|\mathrm{\Delta }m_{32}^2|<|\mathrm{\Delta }m_{31}^2|`$ where $`\mathrm{\Delta }m_{ij}^2m_i^2m_j^2`$. The flavor eigenstates are related to the mass eigenstates by $`U_{\alpha j}`$ ($`\alpha =e,\mu ,\tau `$), where $`U_{\alpha j}`$ are the elements of the MNS mixing matrix U $`^\mathrm{?}`$:
$`\left(\begin{array}{c}\nu _e\\ \nu _\mu \\ \nu _\tau \end{array}\right)=U\left(\begin{array}{c}\nu _1\\ \nu _2\\ \nu _3\end{array}\right),`$ (7)
$`U`$ $``$ $`\left(\begin{array}{ccc}U_{e1}& U_{e2}& U_{e3}\\ U_{\mu 1}& U_{\mu 2}& U_{\mu 3}\\ U_{\tau 1}& U_{\tau 2}& U_{\tau 3}\end{array}\right)=\left(\begin{array}{ccc}c_{12}c_{13}\hfill & s_{12}c_{13}\hfill & s_{13}e^{i\delta }\hfill \\ s_{12}c_{23}c_{12}s_{23}s_{13}e^{i\delta }\hfill & c_{12}c_{23}s_{12}s_{23}s_{13}e^{i\delta }\hfill & s_{23}c_{13}\hfill \\ s_{12}s_{23}c_{12}c_{23}s_{13}e^{i\delta }\hfill & c_{12}s_{23}s_{12}c_{23}s_{13}e^{i\delta }\hfill & c_{23}c_{13}\hfill \end{array}\right).`$
With the mass hierarchy $`|\mathrm{\Delta }m_{21}^2||\mathrm{\Delta }m_{32}^2|`$ there are two possible mass patterns which are depicted in Fig. 1a and 1b, depending on whether $`\mathrm{\Delta }m_{32}^2`$ is positive or negative. If the matter effect is relevant then the sign of $`\mathrm{\Delta }m_{32}^2`$ can be determined by distinguishing neutrinos and anti-neutrinos. The Superkamiokande experiment uses water Cherenkov detectors and events of neutrinos and anti-neutrinos are unfortunately indistinguishable, so the sign of $`\mathrm{\Delta }m_{32}^2`$ is unknown to date.
It has been shown in the three flavor framework $`^\mathrm{?}`$ that combination of the CHOOZ reactor data $`^\mathrm{?}`$ and the atmospheric neutrino data of the Kamiokande and the Superkamiokande implies very small $`\theta _{13}`$, i.e., $`\mathrm{sin}^22\theta _{13}<0.1`$ which is essentially the result of the CHOOZ data. When $`|\theta _{13}|`$ is small, the MNS matrix looks like
$`U\left(\begin{array}{ccc}c_{}\hfill & s_{}\hfill & ϵ\hfill \\ s_{}c_{\text{atm}}\hfill & c_{}c_{\text{atm}}\hfill & s_{\text{atm}}\hfill \\ s_{}s_{\text{atm}}\hfill & c_{}s_{\text{atm}}\hfill & c_{\text{atm}}\hfill \end{array}\right),`$
where $`\theta _{12}`$, $`\theta _{23}`$ have been replaced by $`\theta _{}`$ and $`\theta _{\text{atm}}`$, respectively.
The measurement of $`\theta _{}\theta _{12}`$ and $`\theta _{\text{atm}}\theta _{23}`$ is expected to be greatly improved in the future experiments on solar and atmospheric neutrinos, so the remaining problems in the three flavor framework are to determine (1) the sign of $`\mathrm{\Delta }m_{32}^2`$, (2) the magnitude of $`\theta _{13}`$, (3) the magnitude of the CP phase $`\delta `$. Recently a lot of research have been done on neutrino factories, and the three problems mentioned above may be solved at neutrino factories. In this talk I would like to discus these three issues in some detail and show which set of parameters optimizes each signal. I will assume that the volume of the detector is 10 kt, the intensity of the beam is 10<sup>21</sup> muons/yr, and the data are taken for one year as the reference values in the following discussions. I will also assume $`E_\mu `$ 50 GeV.
## 2 Neutrino factories
Before discussing the three problems given at the end of the Introduction let me give a little background for neutrino factories. As has been shown in $`^{\mathrm{?},\mathrm{?}}`$, the information of neutrino oscillations can be obtained by looking at “wrong sign muons” which are produced in $`\nu _e\nu _\mu \mu ^{}`$ or $`\overline{\nu }_e\overline{\nu }_\mu \mu ^+`$ and the numbers $`N_{\text{wrong}}(\mu ^\pm )`$ of the wrong sign muons are given by $`^\mathrm{?}`$
$`N_{\text{wrong}}(\mu ^{})`$ $`=`$ $`n_T{\displaystyle \frac{12E_\mu ^2}{\pi L^2m_\mu ^2}}{\displaystyle d\left(\frac{E_\nu }{E_\mu }\right)\left(\frac{E_\nu }{E_\mu }\right)^2\left(1\frac{E_\nu }{E_\mu }\right)\sigma _{\nu N}(E_\nu )P(E_\nu )}`$
$`N_{\text{wrong}}(\mu ^+)`$ $`=`$ $`n_T{\displaystyle \frac{12E_\mu ^2}{\pi L^2m_\mu ^2}}{\displaystyle d\left(\frac{E_{\overline{\nu }}}{E_\mu }\right)\left(\frac{E_{\overline{\nu }}}{E_\mu }\right)^2\left(1\frac{E_{\overline{\nu }}}{E_\mu }\right)\sigma _{\overline{\nu }N}(E_{\overline{\nu }})P(E_{\overline{\nu }})},`$
where $`E_\mu `$ is the muon energy, $`L`$ is the length of the neutrino path, $`n_T`$ is the number of the target nucleons, $`\sigma _{\nu N}(E_\nu )`$ and $`\sigma _{\overline{\nu }N}(E_{\overline{\nu }})`$ are the (anti-)neutrino nucleon cross sections given by
$`\sigma _{\nu N}(E_\nu )`$ $`=`$ $`\left({\displaystyle \frac{E_\nu }{\text{GeV}}}\right)\times 0.67\times 10^{38}\text{cm}^2`$
$`\sigma _{\overline{\nu }N}(E_{\overline{\nu }})`$ $`=`$ $`\left({\displaystyle \frac{E_{\overline{\nu }}}{\text{GeV}}}\right)\times 0.33\times 10^{38}\text{cm}^2,`$
and $`P(E_\nu )`$ and $`P(E_{\overline{\nu }})`$ are the oscillation probabilities which are given by (on the assumption of constant density of the matter)
$`P(E_\nu )`$ $``$ $`P(\nu _e\nu _\mu )=s_{23}^2\mathrm{sin}^22\theta _{13}^{M^{()}}\mathrm{sin}^2\left({\displaystyle \frac{B^{()}L}{2}}\right)`$
$`P(E_{\overline{\nu }})`$ $``$ $`P(\overline{\nu }_e\overline{\nu }_\mu )=s_{23}^2\mathrm{sin}^22\theta _{13}^{M^{(+)}}\mathrm{sin}^2\left({\displaystyle \frac{B^{(+)}L}{2}}\right),`$ (13)
where $`A\sqrt{2}G_FN_e`$ stands for the matter effect $`^\mathrm{?}`$ of the Earth, $`\theta _{13}^{M^{(\pm )}}`$ is the effective mixing angle in matter given by
$`\mathrm{tan}2\theta _{13}^{M^{(\pm )}}{\displaystyle \frac{\mathrm{\Delta }E_{32}\mathrm{sin}2\theta _{13}}{\mathrm{\Delta }E_{32}\mathrm{cos}2\theta _{13}\pm A}},`$ (14)
and
$`B^{(\pm )}\sqrt{\left(\mathrm{\Delta }E_{32}\mathrm{cos}2\theta _{13}\pm A\right)^2+\left(\mathrm{\Delta }E_{32}\mathrm{sin}2\theta _{13}\right)^2}.`$ (15)
Using these formula, the ratio $`N_{\text{wrong}}(\mu ^+)/N_{\text{correct}}(\mu ^{})`$ is plotted in Fig. 2a and 2b as a function of $`E_\mu `$ and $`L`$ for typical values of $`\theta _{13}`$ with $`\mathrm{sin}^22\theta _{23}=1.0`$ and $`\mathrm{\Delta }m_{32}^2=3.5\times 10^3`$eV<sup>2</sup>. As was shown by Gomez-Cadenas $`^\mathrm{?}`$ the ratio of (background events)/$`N_{\text{correct}}(\mu )`$ is of order 10<sup>-5</sup>. In the case of smaller value of $`\theta _{13}`$ ($`\theta _{13}=1^{}`$ or $`\mathrm{sin}^22\theta _{13}=0.001`$), it should be possible to detect wrong sign events at neutrino factories if $`L>`$1000km for all the muon energies ($`10E_\mu 50`$ GeV), with our reference values of the beam and the detector (10$`{}_{}{}^{21}\mu `$/yr$``$10kt$``$1yr).
## 3 The sign of $`\mathrm{\Delta }m_{32}^2`$
As was mentioned in the Introduction, the mass pattern corresponds to either Fig. 1a or 1b, depending on whether $`\mathrm{\Delta }m_{32}^2`$ is positive or negative. Determination of this mass pattern is important, since Figs. 1a and 1b correspond to one and two mass states, assuming that the lowest mass is almost zero<sup>b</sup><sup>b</sup>bThe mixed dark scenario in which neutrinos have masses of order 1 eV seems to be disfavored $`^\mathrm{?}`$..
As we can see from (13), if $`\mathrm{\Delta }m_{32}^2>0`$ then the effective mixing angle $`\theta _{13}^{M()}`$ is enhanced and $`P(\nu _e\nu _\mu )`$ increases. On the other hand, if $`\mathrm{\Delta }m_{32}^2<0`$ then $`\theta _{13}^{M(+)}`$ is enhanced and $`P(\overline{\nu }_e\overline{\nu }_\mu )`$ increases. So, at neutrino factories where baseline is relatively large and therefore the matter effect plays an important role, the sign of $`\mathrm{\Delta }m_{32}^2`$ may be determined by looking at the difference between neutrino and anti-neutrino events which should reflect the difference between $`P(\nu _e\nu _\mu )`$ and $`P(\overline{\nu }_e\overline{\nu }_\mu )`$. In Figs. 3 and 4 the numbers of events per each neutrino energy $`E_\nu `$ ($`E_{\overline{\nu }}`$ and $`E_\nu `$ is identified here) at a neutrino factory with the parameter 10$`{}_{}{}^{21}\mu `$/yr$``$10kt$``$1yr are given for $`\mathrm{\Delta }m_{32}^2=3.5\times 10^3`$eV<sup>2</sup> (Fig. 3) and $`\mathrm{\Delta }m_{32}^2=3.5\times 10^3`$eV<sup>2</sup> (Fig. 4), respectively (I have assumed $`\mathrm{sin}^22\theta _{23}=1.0`$ and $`\mathrm{sin}^22\theta _{13}=0.095`$). The black and gray lines stand for wrong sign $`\mu ^{}`$ and $`\mu ^+`$ events, respectively, and five cases of the muon energy $`E_\mu =10,20,\mathrm{},50`$ GeV are plotted. For later purposes, behaviors with respect to the CP violating phase $`\delta `$ are also considered and the solid, dotted and dashed lines are the numbers of events with $`\delta =0`$, $`\delta =\pi /2`$, $`\delta =\pi /2`$, respectively. The deviation of the solid lines from the dotted or dashed lines is mainly due to the matter effects, and the deviation of the dotted lines from the dashed ones is in general smaller.
Since the cross section $`\sigma _{\nu N}`$ and $`\sigma _{\overline{\nu }N}`$ are different (the ratio is 2 to 1), it is useful to look at the quantity
$`{\displaystyle \frac{N_\nu 2N_{\overline{\nu }}}{\delta (N_\nu 2N_{\overline{\nu }})}}={\displaystyle \frac{N_\nu 2N_{\overline{\nu }}}{\sqrt{N_\nu +4N_{\overline{\nu }}}}}`$
whose absolute value should be much larger than one to demonstrate $`\mathrm{\Delta }m_{32}^2>0`$ or $`\mathrm{\Delta }m_{32}^2<0`$. Now let me introduce the quantity
$`R{\displaystyle \frac{\left[N_{\text{wrong}}(\mu ^{})2N_{\text{wrong}}(\mu ^+)\right]^2}{N_{\text{wrong}}(\mu ^{})+4N_{\text{wrong}}(\mu ^+)}}.`$
If $`R1`$ then we can deduce the sign of $`\mathrm{\Delta }m_{32}^2`$. The contour plot of $`R`$=const. is given in Fig. 5a and 5b for typical values of $`\theta _{13}`$ with $`\mathrm{\Delta }m_{32}^2=3.5\times 10^3`$eV$`{}_{}{}^{2}>`$0, $`\mathrm{sin}^22\theta _{23}=1.0`$. If $`\mathrm{sin}^22\theta _{13}`$ is not smaller than $`10^3`$, it is possible to determine the sign of $`\mathrm{\Delta }m_{32}^2`$. Irrespective of the value of $`\theta _{13}`$, $`L`$ 5000km, $`E_\mu =`$ 50 GeV seem to optimize the signal, as far as the quantity $`R`$ is concerned.
## 4 The magnitude of $`\theta _{13}`$
As we can see in (13), it is necessary to know the precise value of $`\theta _{23}`$ to determine $`\theta _{13}`$ accurately. For this purpose, it is useful to measure the independent quantity
$`1P(\nu _\mu \nu _\mu )`$ $`=`$ $`s_{23}^4\mathrm{sin}^22\theta _{13}^{M^{()}}\mathrm{sin}^2\left({\displaystyle \frac{B^{()}L}{2}}\right)`$ (16)
$`+`$ $`\mathrm{sin}^22\theta _{23}[\mathrm{sin}^22\theta _{13}^{M^{()}}\mathrm{sin}^2{\displaystyle \frac{L}{4}}(\mathrm{\Delta }E_{32}+AB^{()})`$
$`+`$ $`\mathrm{cos}^22\theta _{13}^{M^{()}}\mathrm{sin}^2{\displaystyle \frac{L}{4}}(\mathrm{\Delta }E_{32}+A+B^{()})].`$
Combining (13) and (16) and assuming that our knowledge on the density profile of the Earth is exact, we can determine $`\theta _{13}`$ and $`\theta _{23}`$. In practice, however, there is always uncertainty in the density in the Earth, particularly the density deep inside of the Earth (i.e., large $`L`$) is not very well known, so to determine $`\theta _{13}`$ precisely the neutrino path $`L`$ had better be small, say, $`L<1000`$km, as long as $`N_{\text{wrong}}(\mu )`$ exceeds the number of the background events.
## 5 The magnitude of $`\delta `$
In the hierarchical limit, to first order in $`\mathrm{\Delta }E_{21}/\mathrm{\Delta }E_{32}`$ and $`\mathrm{\Delta }E_{21}/A`$, the appearance probabilities are given by
$`\left\{\begin{array}{c}P(\nu _e\nu _\mu )\\ P(\overline{\nu }_e\overline{\nu }_\mu )\end{array}\right\}`$ $``$ $`s_{23}^2\mathrm{sin}^22\theta _{13}^{M^{()}}\mathrm{sin}^2\left({\displaystyle \frac{B^{()}L}{2}}\right)`$
$``$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{\Delta }E_{21}\mathrm{\Delta }E_{32}}{\lambda _+^{()}\lambda _{}^{()}}}\mathrm{sin}\delta \mathrm{sin}2\theta _{12}\mathrm{sin}2\theta _{23}\mathrm{sin}2\theta _{13}^{M^{()}}`$
$`\times `$ $`\mathrm{sin}\left({\displaystyle \frac{\lambda _+^{()}L}{2}}\right)\mathrm{sin}\left({\displaystyle \frac{\lambda _{}^{()}L}{2}}\right)\mathrm{sin}\left({\displaystyle \frac{B^{()}L}{2}}\right),`$
where
$`\lambda _\pm ^{()}{\displaystyle \frac{1}{2}}\left(A\mathrm{\Delta }E_{32}\pm B^{()}\right),\lambda _\pm ^{(+)}{\displaystyle \frac{1}{2}}\left(A\mathrm{\Delta }E_{32}\pm B^{(+)}\right).`$
Because of the matter effect, one of the two independent probabilities is enhanced while the other is suppressed. For the enhanced channel ($`\nu _e\nu _\mu `$ in the case of $`\mathrm{\Delta }m_{32}^2>0`$, $`\overline{\nu }_e\overline{\nu }_\mu `$ in the case of $`\mathrm{\Delta }m_{32}^2<0`$), the number of events is large and one might want to take advantage of this large number. Here I therefore would like to consider the following quantity
$`R_\delta {\displaystyle \frac{\left[N_{\text{wrong}}(\delta =\frac{\pi }{2})N_{\text{wrong}}(\delta =\frac{\pi }{2})\right]^2}{2\left[N_{\text{wrong}}(\delta =\frac{\pi }{2})+N_{\text{wrong}}(\delta =\frac{\pi }{2})\right]}}.`$ (19)
The denominator in (19) corresponds to square of statistical fluctuations of the number of events obtained from the T-invariant probability $`P(\nu _e\nu _\mu ;\delta )+P(\nu _e\nu _\mu ;\delta )`$ while the numerator does to squqare of the number of events obtained from T-violating probability $`P(\nu _e\nu _\mu ;\delta )P(\nu _e\nu _\mu ;\delta )`$. This $`R_\delta `$ is the quantity of T-violation instead of CP-violation. In fact $`R_\delta `$ cannot be determined by the experimental data only, but it requires knowledge on $`\mathrm{\Delta }m_{ij}^2`$, $`\theta _{ij}`$ and the density profile of the Earth to deduce $`R_\delta `$. Nevertheless, having large value of $`R_\delta `$ is a necessary condition to be able to measure $`\delta `$ and the experiments should be designed so that $`R_\delta `$ be maximized.
This suggestion is different from the one in $`^\mathrm{?}`$ in which it is proposed to subtract $`(N_\nu 2N_{\overline{\nu }})/(N_\nu +2N_{\overline{\nu }})`$ by the matter effect term. In either case, one needs the precise knowledge on $`\mathrm{\Delta }m_{ij}^2`$, $`\theta _{ij}`$ and the density of the Earth. To demonstrate $`\delta 0`$ it is necessary that $`R_\delta 1`$. The contour plot of the ratio $`R_\delta `$ is given in Figs. 6a and 6b for two sets of parameters. For the set of the oscillation parameters $`(\mathrm{\Delta }m_{21}^2,\mathrm{sin}^22\theta _{12})=(1.8\times 10^5\mathrm{eV}^2,0.76)`$, $`(\mathrm{\Delta }m_{32}^2,\mathrm{sin}^22\theta _{23})=(3.5\times 10^3\mathrm{eV}^2,1.0)`$ (Fig. 6a), which gives the best fit to the data of solar and atmospheric neutrinos, we have<sup>c</sup><sup>c</sup>c For solar neutrinos, there are three possible sets of parameters which give a very good fit to the data, but here I take the most optimistic one (large mixing angle MSW solution) to observe CP violation. For other solar neutrino solutions, observation of CP violation is either very difficult or impossible.
$`\underset{L,E_\mu }{\mathrm{max}}R_\delta 2.3`$ (20)
for $`\mathrm{sin}^22\theta _{13}=0.09`$, $`L3000`$km, $`\delta =\pi /2`$, $`E_\mu `$40 GeV with 10$`{}_{}{}^{21}\mu `$/yr$``$10kt$``$1yr.
If we take other set of parameters, e.g., $`(\mathrm{\Delta }m_{21}^2,\mathrm{sin}^22\theta _{12})=(1\times 10^4\mathrm{eV}^2,1.0)`$, $`(\mathrm{\Delta }m_{32}^2,\mathrm{sin}^22\theta _{23})=(6\times 10^3\mathrm{eV}^2,0.8)`$ (Fig. 6b), which are still in the allowed region of 99 %CL of data of solar and atmospheric neutrinos, we have
$`\underset{L,E_\mu }{\mathrm{max}}R_\delta 200`$
for $`\mathrm{sin}^22\theta _{13}=0.09`$, $`L3000`$km, $`\delta =\pi /2`$, $`E_\mu `$50 GeV with 10$`{}_{}{}^{21}\mu `$/yr$``$10kt$``$1yr. I have also calculated the ratio $`R_\delta `$ for smaller value of $`\theta _{13}`$, and found that the signal is optimized for almost the same set of the parameters ($`E_\mu `$50 GeV, $`L3500`$km), although the finite value of $`R_\delta `$ is obtained as a limit of 0/0 for very small value of $`\theta _{13}`$ and one has to make sure that we have a certain amount of events to be conclusive. If (20) happens to be the case, then after running the experiment for several years it may be possible to demonstrate $`\delta 0`$.
## 6 Summary
In this talk I have discussed some quantities (the sign of $`\mathrm{\Delta }m_{32}^2`$, the magnitude of $`\theta _{13}`$, the magnitude of $`\delta `$) which can be measured at neutrino factories, and made efforts to optimize the signals. Of course all the measurements would be impossible if $`\mathrm{sin}^22\theta _{13}10^3`$, but otherwise we may be able to determine the sign of $`\mathrm{\Delta }m_{32}^2`$, the values of $`\theta _{13}`$, $`\theta _{23}`$ and the value of $`\delta `$ with our reference values 10$`{}_{}{}^{21}\mu `$/yr$``$10kt$``$1yr. To measure $`\delta `$ it is necessary to know with great precision quantities such as the cross sections $`\sigma _{\nu N}`$, $`\sigma _{\overline{\nu }N}`$, the mixing angles $`\theta _{12}`$, $`\theta _{13}`$ $`\theta _{23}`$, the mass squared differences $`\mathrm{\Delta }m_{21}^2`$, $`\mathrm{\Delta }m_{32}^2`$ as well as density profile of the Earth. It is important to estimate how much the uncertainty on $`\mathrm{\Delta }m_{ij}^2`$, $`\theta _{ij}`$ and the density profile of the Earth affects the reach of the experiment, particularly in the measurement of CP violation, and this is a subject for future study.
## Acknowledgments
I would like to thank Yoshitaka Kuno and Yoshiharu Mori for suggesting this subject and for discussions. This research was supported in part by a Grant-in-Aid for Scientific Research of the Ministry of Education, Science and Culture, #12047222, #10640280.
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# Melting behavior of large disordered sodium clusters.
## I Introduction
Cluster melting is a topic of current theoretical and experimental interest, motivated by the observation of several size-dependent properties which have no analog in the bulk phase. Between those properties, we can mention the following: (i) the melting like transition does not occur at a well defined temperature as in the solid, but spreads over a finite temperature interval that widens as the cluster size decreases. The lower end of that interval defines the freezing temperature T<sub>f</sub> below which the cluster is completely solidlike, their constituent atoms just vibrating about their equilibrium positions. The upper part of the interval defines the melting temperature above which all the atoms can diffuse across the cluster and the “liquid” phase is completely established. Between those two temperatures the cluster is not fully solid nor fully liquid. It is in that transition region where the cluster-specific behavior emerges: (ia) Premelting effects like partial melting of the cluster (the most usual case is surface melting) or structural isomerizations upon heating, which lead to a melting in steps. (ib) The dynamic coexistence regime, where the cluster can fluctuate in time between being completely solid or liquid. (ii) Strong nonmonotonic variations of the melting temperature with size have been found in recent experiments on sodium clusters. The maxima in the melting temperature are not in exact correspondence with either electronic or atomic shell closing numbers, but bracketed by the two, suggesting that both effects are relevant to the melting process. It is important to note that the values of T<sub>f</sub> and T<sub>m</sub> as defined above are not yet amenable to the experiments, and that the experimental melting temperature is somewhere between these two values.
Previously we have reported density functional orbital-free molecular dynamics (OFMD) simulations of the melting process in sodium clusters Na<sub>N</sub>, with N=8,20,55,92, and 142. The OFMD technique is completely analogous to the method devised by Car and Parrinello (CP) to perform dynamic simulations at an ab initio level, but the electron density is taken as the dynamic variable, as opposed to the Kohn-Sham (KS) orbitals in the original CP method. This technique has been already used both in solid state and cluster physics. In contrast to simulations which use empirical interatomic potentials, the detailed electronic structure and the electronic contribution to the energy and the forces on the ions are recalculated efficiently every atomic time-step. The main advantage over KS-based methods is that the computational effort to update the electronic system increases linearly with the cluster size N, in contrast to the N<sup>3</sup> scaling of orbital-based methods. Indeed, these were the first molecular dynamics simulations of large clusters as Na<sub>92</sub> and Na<sub>142</sub> that included an explicit treatment of the electronic degrees of freedom.
A very important issue in the simulations of cluster melting is the election of the low-temperature isomer to be heated. A good knowledge of the ground state structure (global minimum) is required, as the details of the melting transition are known to be isomer-dependent. But the problem of performing a realistic global optimization search is exponentially difficult as size increases, so finding the global minima of Na<sub>92</sub> and Na<sub>142</sub> becomes impractical. In our previous work we directly started from icosahedral isomers, as there is some experimental and theoretical indications that suggest icosahedral packing in sodium clusters, and found a good agreement with the experimental results of Haberland’s group. However, we were unable to find those icosahedral structures by an unconstrained search method as simulated annealing, which always led to disordered isomers both for Na<sub>92</sub> and Na<sub>142</sub>. Although the icosahedral structures are more stable in all the cases, the energy difference between both isomers is approximately 0.02 eV/atom, which is very small. Amorphouslike structures have been found recently to be the ground state isomers of gold clusters for a number of sizes, and pair potential calculations performed by Doye and Wales predict that the amorphous state is favored by long potential ranges. The specific features of those structures are little or no spatial symmetry and a pair distribution function typical of glasses. Besides that, one usually finds a large number of amorphouslike isomers nearly degenerate in energy, which suggests that they occupy a substantial fraction of the phase space available to the cluster. Both the energy proximity to the more stable icosahedral isomers and the large entropy associated with the amorphous part of the phase space make plausible that amorphous isomers could be present in the experimental cluster beams, so their thermal properties deserve specific investigation. Apart from this, the study of melting in amorphouslike clusters is intrinsically interesting from a theoretical point of view. Thus, the goals of this work are to study the mechanisms by which the melting-like transition proceeds in two disordered isomers of Na<sub>92</sub> and Na<sub>142</sub>, to study the influence on the melting behavior of the disorder degree, and to make a meaningful comparison with the melting behavior of the corresponding icosahedral isomers. In the next section we briefly present some technical details of the method. The results are presented and discussed in section III and, finally, section IV summarizes our main conclusions.
## II Theory
The orbital-free molecular dynamics method is a Car-Parrinello total energy scheme which uses an explicit kinetic-energy functional of the electron density, and has the electron density as the dynamical variable, as opposed to the KS single particle wavefunctions. The main features of the energy functional and the calculation scheme have been described at length in previous work, and details of our method are as described by Aguado et al. In brief, the electronic kinetic energy functional of the electron density, $`n(\stackrel{}{r})`$, corresponds to the gradient expansion around the homogeneous limit through second order
$$T_s[n]=T^{TF}[n]+\frac{1}{9}T^W[n],$$
(1)
where the first term is the Thomas-Fermi functional (Hartree atomic units have been used)
$$T^{TF}[n]=\frac{3}{10}(3\pi ^2)^{2/3}n(\stackrel{}{r})^{5/3}𝑑\stackrel{}{r},$$
(2)
and the second is the lowest order gradient correction, where T<sup>W</sup>, the von Weizsäcker term, is given by
$$T^W[n]=\frac{1}{8}\frac{n(\stackrel{}{r})^2}{n(\stackrel{}{r})}𝑑\stackrel{}{r}.$$
(3)
The local density approximation is used for exchange and correlation. In the external field acting on the electrons, $`V_{ext}(\stackrel{}{r})=_nv(\stackrel{}{r}\stackrel{}{R}_n)`$, we take $`v`$ to be the local pseudopotential of Fiolhais et al., which reproduces well the properties of bulk sodium and has been shown to have good transferability to sodium clusters. The cluster is placed in a unit cell of a cubic superlattice, and the set of plane waves periodic in the superlattice is used as a basis set to expand the valence density. Following Car and Parrinello, the coefficients of that expansion are regarded as generalized coordinates of a set of fictitious classical particles, and the corresponding Lagrange equations of motion for the ions and the electron density distribution are solved as described in Ref. .
The calculations used a supercell of edge 71 a.u. and the energy cut-off in the plane wave expansion of the density was 8 Ryd. In all cases, a 64$`\times `$64$`\times `$64 grid was used. Previous tests indicate that the cut-offs used give good convergence of bond lengths and binding energies. The fictitious mass associated with the electron density coefficients ranged between 1.0$`\times `$10<sup>8</sup> and 3.3$`\times `$10<sup>8</sup> a.u., and the equations of motion were integrated using the Verlet algorithm for both electrons and ions with a time step ranging from $`\mathrm{\Delta }`$t = 0.73 $`\times `$ 10<sup>-15</sup> sec. for the simulations performed at the lowest temperatures, to $`\mathrm{\Delta }`$t = 0.34 $`\times `$ 10<sup>-15</sup> sec. for those at the highest ones. These choices resulted in a conservation of the total energy better than 0.1 %.
Several molecular dynamics simulation runs at different constant energies were performed in order to obtain the caloric curve for each cluster. The initial positions of the atoms for the first run were taken by slightly deforming the equilibrium low temperature geometry of the cluster. The final configuration of each run served as the starting geometry for the next run at a different energy. The initial velocities for every new run were obtained by scaling the final velocities of the preceding run. The total simulation times varied between 8 and 18 ps for each run at constant energy.
A number of indicators to locate the melting-like transition were employed. Namely, the specific heat defined by
$$C_v=[NN(1\frac{2}{3N6})<E_{kin}>_t<E_{kin}^1>_t]^1,$$
(4)
where N is the number of atoms and $`<>_t`$ indicates the average along a trajectory; the time evolution of the distance between each atom and the center of mass of the cluster
$$r_i(t)=\stackrel{}{R}_i(t)\stackrel{}{R}_{cm}(t);$$
(5)
the average over a whole dynamical trajectory of the radial atomic distribution function g(r), defined by
$$dN_{at}=g(r)dr$$
(6)
where $`dN_{at}(r)`$ is the number of atoms at distances from the center of mass between r and r + dr; and finally, the diffusion coefficient
$$D=\frac{1}{6}\frac{d}{dt}(<r^2(t)>),$$
(7)
which is obtained from the long time behavior of the mean square displacement $`<r^2(t)>=\frac{1}{Nn_t}_{j=1}^{n_t}_{i=1}^N[\stackrel{}{R}_i(t_{0j}+t)\stackrel{}{R}_i(t_{0j})]^2`$, where $`n_t`$ is the number of time origins, $`t_{0j}`$, considered along a trajectory.
## III Results
The most stable disordered structures of Na<sub>92</sub> and Na<sub>142</sub> that we have found with the simulated annealing technique are shown in Fig. 1, together with the relaxed icosahedral isomers obtained as explained in our previous work. The dynamical annealing runs were started from high-temperature liquid isomers thermalized during 30 ps at 600 K. The cooling strategy was to reduce the internal cluster temperature by a factor of 0.9999 of its actual value each twelve atomic time steps. With the chosen time step of 0.34 $`\times `$ 10<sup>-15</sup> sec. the temperature reduction is applied each four femtoseconds. A first important difference between the icosahedral and disordered isomers is that the former ones have a smoother surface. Besides that, no apparent spatial symmetry is observed in the disordered isomers. In Fig. 2 we show the short-time averages (sta) of the distances between each atom and the center of mass of the cluster for both isomers of Na<sub>142</sub>, obtained from an OFMD run at a low temperature. The values of $`<r_i(t)>_{sta}`$ are almost independent of time. The clusters are solid, the atoms merely vibrating around their equilibrium positions. Curve crossings are due to oscillatory motion and slight structural relaxations rather than diffusive motion. The difference between Figs. 2a and 2b is due to structure. For the icosahedral isomer in Fig. 2a one can distinguish quasidegenerate groups which are characteristic of the symmetry: one line near the center of mass of the cluster identifies the central atom, whose position does not exactly coincide with the center of mass because of the location of the five surface vacancies (147 atoms are needed to form a complete three-shell icosahedron); 12 lines correspond to the first icosahedral shell; another 42 complete the second shell, within which we can distinguish the 12 vertex atoms from the rest because their distances to the center of mass are slightly larger; finally, 82 lines describe the external shell, where again we can distinguish the 7 vertex atoms from the rest. In contrast, the lines for the disordered isomer in Fig. 2b are quite dispersed. Nevertheless, there is a narrow interval where the ionic density is very low, that serves to define separate surface and inner atomic shells. The case of Na<sub>92</sub> (not shown) is similar, but in that case the structure is more uniformously amorphous, and there is no way to distinguish between surface and inner shells. We will see below that this small difference is very important. Soler at al. have compared the structures of icosahedral and amorphous gold clusters, with similar results: while the atoms in the icosahedral isomers are clearly distributed in atomic shells with different radial distances to the center of mass of the cluster, in the case of amorphous clusters there are “atomic transfers” between shells that blur the atomic shell structure. In the case of gold clusters, the amorphous isomers turn out to be the minimum energy structures for a number of sizes. In the case of Na<sub>92</sub> and Na<sub>142</sub>, the icosahedral isomers are more stable than the lowest energy disordered isomers found (0.017 eV/atom and 0.020 eV/atom for Na<sub>92</sub> and Na<sub>142</sub> respectively).
For each cluster we have calculated the internal temperature, defined as the average of the ionic kinetic energy, as a function of the total energy– the so-called caloric curve. A thermal phase transition is indicated in the caloric curve by a change of slope, the slope being the specific heat; the height of the step gives an estimate of the latent heat of fusion. However, melting processes are more easily recognized as peaks in the specific heat as a function of temperature, that has been calculated directly from eq. (4). The fact that the specific heat peaks occur at the same temperatures as the slope changes of the caloric curve (see curves below) gives us confidence on the convergence of our results, as those two quantities have been obtained in different ways.
The specific heat curves for Na<sub>142</sub> (fig. 3) display two main peaks, suggestive of a two-step melting process. For the disordered cluster the two peaks are well separated in temperature, T$`{}_{s}{}^{dis}`$ 130 K and T$`{}_{m}{}^{dis}`$ 270 K, whereas they are much closer together for the icosahedral cluster, T$`{}_{s}{}^{ico}`$ 240 K and T$`{}_{m}{}^{ico}`$ 270 K, so close that only one slope change in the caloric curve can be distinguished in this case. The results suggest that the melting transition starts at a temperature T<sub>s</sub> and finishes at T<sub>m</sub> with the difference in the melting of the two isomers being the smaller T<sub>s</sub> value of the disordered isomer. In our previous work we showed that for the icosahedral cluster those two steps are due to surface melting and homogeneous melting, respectively. Here we show that a similar explanation is valid for the disordered Na<sub>142</sub> isomer. At T=160 K, a temperature between T$`{}_{}{}^{dis}{}_{s}{}^{}`$ and T$`{}_{}{}^{dis}{}_{m}{}^{}`$, the structure of the disordered cluster is more fluid than at low temperature (fig. 2). Fig. 4 indicates that the spherical atomic shells have separately melted, atoms undergoing diffusive motions mainly across a given shell, as seen in the bold lines, without an appreciable intershell diffusion (although some occasional intershell displacement has been observed). The larger spread of the upper bold line indicates that diffusion is appreciably faster in the surface shell. Thus the transition at 130 K can be identified with intrashell melting. Why it occurs at a rather low temperature is associated with the large spread in the radial distances of the atoms in each shell. The atomic shells are structurally disordered at low temperature, although kinetically frozen, like in a typical glass, so inducing diffusion in the shells of that cluster is rather easy and occurs at moderate temperatures. The surface melting stage does not develop in the icosahedral isomer until a temperature of T$`{}_{s}{}^{ico}`$ 240 K is reached. At this temperature, the time evolution of the surface atomic distances to the cluster center becomes very similar to those of the disordered Na<sub>142</sub> isomer at 160 K. Inducing diffusion in the surface of the icosahedral isomer requires a higher temperature because of the higher structural order of its surface. Once the surface of both isomers has melted, homogeneous melting occur at the same temperature, $``$ 270 K, in very good agreement with the experimental value of 280 K. From that temperature on, the liquid phase is completely established (all atoms can diffuse across the whole cluster volume) and differences between both isomers have no sense anymore.
The thermal behavior of both isomers is not so different. The radial atomic density distribution of the disordered isomer at 160 K (Fig. 5) shows a smoothed shape with respect to that found at low T, but a distribution of the atoms in separate surface and inner atomic shells can be clearly distinguished. In fact, the intermediate temperature distribution is similar to that found for the icosahedral isomer in the same temperature region, as the atomic shells are equally distinguished. The small gap in the ionic density (Figs. 2 and 5) of the disordered isomer at low temperature drives the cluster towards a well defined shell structure upon heating. We conclude that, despite the high orientational disorder in both surface and inner shells, the cluster can not be considered completely amorphous, as Fig. 5 at intermediate temperature shows some radial atomic ordering. There are still more similarities. The solidlike phase of the icosahedral isomer disappears as soon as a temperature of 130 K is reached, even though no peak in the specific heat is detected: there are isomerization transitions between different permutational isomers which preserve the icosahedral structure. These isomerizations involve the displacement of the five surface vacancies across the outer shell. Thus, both isomers leave the solidlike phase at $``$ 130 K, the only difference being that one has direct access to the intrashell melting stage while the other enters first an isomerization regime. Calvo and Spiegelmann have related the appareance of specific heat peaks to sudden openings of the available phase space. In the isomerization regime the icosahedral cluster has access just to a limited number of symmetry-equivalent isomers, while the phase structure of the disordered cluster opens suddenly to include a very large number of isomers, as all the possible position interchanges between two atoms of a given shell are allowed. Thus, a specific heat peak appears at T$``$130 K for the disordered isomer, but not for the icosahedral isomer. Any atomic shell distribution in the time average of g(r) disappears completely when homogeneous melting occur (Fig. 5).
The results for Na<sub>92</sub> are shown in fig. 6. Two-step melting is again observed in the icosahedral isomer, with a small prepeak in the specific heat at T$`{}_{s}{}^{ico}`$ 130 K and a large peak, corresponding to homogeneous melting, at T$`{}_{m}{}^{ico}`$ 240 K. In this case T$`{}_{}{}^{ico}{}_{s}{}^{}`$ and T$`{}_{}{}^{ico}{}_{m}{}^{}`$ are well separated. T$`{}_{}{}^{ico}{}_{s}{}^{}`$ is in the range where the isomerization processes in icosahedral Na<sub>142</sub> set in and the intrashell melting stage in disordered Na<sub>142</sub> develops. The larger number of vacancies in the surface shell of icosahedral Na<sub>92</sub>, as compared to icosahedral Na<sub>142</sub>, allows for true surface diffusion and these processes give rise to a distinct peak in the specific heat. The disordered isomer melts gradually over a wide temperature interval, and no prominent specific heat peaks nor important slope changes in the caloric curve are detected. Ercolessi et al. have also found a melting process without a latent heat of fusion for amorphous gold clusters with less than $``$ 90 atoms. In Fig. 7 we show the radial ionic density distribution of disordered Na<sub>92</sub> at several representative temperatures. At a temperature as low as 70 K there is no discernible atomic shell structure. The g(r) function for a higher temperature where the surface of the icosahedral isomer has already melted is not very different. At a temperature where the icosahedral isomer is liquid the only appreciable change in g(r) is due to the thermal expansion of the cluster. The structure of cold disordered Na<sub>92</sub> is both radially and orientationally disordered. The structure is kinetically frozen, but there seems to be no barriers impeding the exploration of the liquid part of the phase space. In fact, the cluster is already in that region at low temperature, as Fig. 7 shows. This is seen most clearly in the evolution of the diffusion behavior with temperature. In Fig. 8 we show the temperature variation of the square root of the diffusion coefficient. While the two steps in the melting of the icosahedral isomer are clearly detected as slope changes at the corresponding temperatures, the value of $`\sqrt{D}`$ for the disordered isomer increases with temperature in a smooth way, without any abrupt change. Thus, the opening of the available phase space proceeds in a gradual way and specific heat peaks are not detected.
We have found a very different thermal behavior for two clusters that were classified in principle as disordered. Although just two examples are not enough to draw general conclusions, we believe that the thermal behavior typical of amorphouslike sodium clusters is that found for Na<sub>92</sub>, and that what is lacking is a clear definition of what an amorphous cluster is. In line with the discussion of “atomic transfers” between shells advanced by Soler et al, we propose that a large orientational disorder is not enough for a cluster to be classified as amorphouslike. Only when those atomic transfers are maximal, in such a way that no local regions with low atomic density exist, the cluster can be considered completely amorphous. The existence of those regions, however small they may be (as is the case of Na<sub>142</sub>), promote the creation of appreciable free energy barriers in the potential energy surface, so sudden access to a substantial region of the available phase space is expected above a certain temperature, and peaks will appear in the temperature evolution of the specific heat. On the contrary, the absence of such low atomic density regions facilitates the diffusion of the atoms across the whole cluster volume right from the start of the heating process. As the liquidlike phase is established precisely when all the atoms in the cluster can diffuse across the cluster volume, we expect that no appreciable free energy barriers will be found in these cases, and no specific heat peaks will be detected.
A few comments regarding the quality of the simulations and of the annealing runs are perhaps in order here. The orbital-free representation of the atomic interactions is much more efficient than the more accurate KS treatments, but is still substantially more expensive computationally than a simulation using phenomenological many-body potentials. Such potentials contain a number of parameters that are usually chosen by fitting some bulk and/or molecular properties. In contrast our model is free of external parameters, although there are approximations in the kinetic and exchange-correlation functionals. The orbital-free scheme accounts, albeit approximately, for the effects of the detailed electronic distribution on the total energy and the forces on the ions. We feel that this is particularly important in metallic clusters for which a large proportion of atoms are on the surface and experience a very different electronic environment than an atom in the interior. Furthermore, the adjustment of the electronic structure and consequently the energy and forces to rearrangements of the ions is also taken into account. But the price to be paid for the more accurate description of the interactions is a less complete statistical sampling of the phase space. The simulation times are substantially shorter than those that can be achieved in phenomenological simulations. The cooling rate employed in the annealing runs is also faster than those that can be achieved by using empirical potentials. Nevertheless, we expect that the locations of the several transitions are reliable, because all the indicators we have used, both thermal and structural ones, are in essential agreement on the temperature at which the transitions start. On the other hand, the disordered isomers found in different annealing runs did not show substantial structural or energetic differences with respect to those studied here.
## IV Summary
The melting-like transition in disordered Na<sub>142</sub> and Na<sub>92</sub> has been investigated by applying an orbital-free, density-functional molecular dynamics method. The computational effort which is required is modest in comparison with the traditional Car-Parrinello Molecular Dynamics technique based on Kohn-Sham orbitals, that would be very costly for clusters of this size. This saving allows the study of large clusters.
A disordered isomer of Na<sub>142</sub> melts in two steps as evidenced by the thermal indicators. The transition at T$`{}_{s}{}^{dis}`$ 130 K is best described as intrashell melting. This is followed at T$`{}_{m}{}^{dis}`$ 270 K by homogeneous melting. Melting is found to depend on the starting low-temperature isomer. Specifically, for an icosahedral Na<sub>142</sub> isomer, the analysis of thermal, macroscopic properties places those two stages much closer in temperature, 240 K and 270 K respectively. Nevertheless, isomerization transitions are observed in the icosahedral isomer at a temperature as low as T$`{}_{s}{}^{dis}`$130 K. These isomerizations involve the motion of the five atomic vacancies in the cluster surface, preserve the icosahedral structure and do not give rise to any pronounced feature in the caloric curve. In the disordered isomer there is not a separate isomerization regime (something rather evident because there is not an underlying ordered structure in each shell), but there is a melting-in-steps process, due to the distribution of the atoms in different shells. Thus, the melting of both isomers is not as different as suggested by the thermal indicators. An icosahedral isomer of Na<sub>92</sub> melts also in a similar way: there is a minor peak in C<sub>v</sub> at T$`{}_{s}{}^{ico}`$130K that indicates surface melting, and a main, homogeneous melting peak at T$`{}_{m}{}^{ico}`$240 K. The lower value of T$`{}_{}{}^{ico}{}_{s}{}^{}`$, as compared to Na<sub>142</sub>, is due to the larger number of surface vacancies, and the melting-in-step process is due to the atomic shell structure. The melting of disordered Na<sub>92</sub> proceeds instead gradually, and spreads over a very wide temperature interval. The thermal indicators as the caloric curve or the specific heat do not show any indications of an abrupt transition, which suggests that the phase space available to the cluster does not increases suddenly at any given temperature. The square root of the diffusion coefficient increases with temperature in a smooth way, in contrast to the step diffusive behavior of icosahedral Na<sub>92</sub>. It has been suggested that the absence of any abrupt transition is closely related to the absence of any shell structure in the radial atomic density distribution, which should be considered a necessary condition for a cluster to be considered completely amorphous. In this, sense, only the disordered isomer of Na<sub>92</sub> can be considered rigorously amorphous, while the disordered Na<sub>142</sub> isomer should be considered just partially amorphous.
In summary, we have found a number of structural properties that have an important effect on the melting properties of sodium clusters. A melting in steps process is to be expected in almost all clusters where a clear distribution of the atoms in radial shells exists, as is the case of both isomers of Na<sub>142</sub> and of the icosahedral isomer of Na<sub>92</sub>. In those cases, intrashell diffusive motion starts at a temperature T<sub>intra</sub>, lower than the temperature T<sub>inter</sub> at which intershell diffusive motion begins to be important. The difference T<sub>inter</sub>-T<sub>intra</sub> is small if the orientational order of the shells is large (for example, icosahedral order) and the number of vacancies in each shell is small: this occurs for icosahedral Na<sub>142</sub>, with just five surface vacancies; A limit case is icosahedral Na<sub>55</sub>. With two complete atomic shells, intrashell motions are as difficult as intershell motions and the two transitions merge into one. When one shell contains a large number of vacancies, the two temperatures are well separated no matter how high the orientational order is: this is exemplified by the case of icosahedral Na<sub>92</sub>. Also, if the shells have a high orientational disorder, the two transitions are well separated in temperature no matter how close we are from an icosahedral shell closing: an example is the disordered Na<sub>142</sub> isomer. Finally, a gradual melting process without any abrupt transition is to be expected for all those clusters which have both orientational and radial disorder, that is for amorphouslike clusters: this is the case of amorphous Na<sub>92</sub>.
ACKNOWLEDGMENTS: This work has been supported by DGES (Grant PB98-0368) and Junta de Castilla y León (VA70/99). The author acknowledges useful discussions with J. M. López, J. A. Alonso, and M. J. Stott.
Captions of Figures.
Figure 1 Structures of the low temperature (a) amorphous Na<sub>142</sub>, (b) icosahedral Na<sub>142</sub>, (c) amorphous Na<sub>92</sub> and (d) icosahedral Na<sub>92</sub> isomers.
Figure 2 Short-time averaged distances $`<r_i(t)>_{sta}`$ between each atom and the center of mass in Na<sub>142</sub>, as functions of time for (a) the icosahedral isomer at T= 30 K and (b) the amorphous isomer at T= 47 K.
Figure 3 Caloric and specific heat curves of Na<sub>142</sub>, taking the internal cluster temperature as the independent variable. The deviation around the mean temperature is smaller than the size of the circles.
Figure 4 Short-time averaged distances $`<r_i(t)>_{sta}`$ between each atom and the center of mass in amorphous Na<sub>142</sub> as functions of time at T= 160 K. The bold lines follow the evolution of a particular atom in the surface shell and another in the outermost core shell.
Figure 5 Time averaged radial atomic density distribution of the amorphous isomer of Na<sub>142</sub>, at some representative temperatures.
Figure 6 Caloric and specific heat curves of Na<sub>92</sub>, taking the internal cluster temperature as the independent variable. The deviation around the mean temperature is smaller than the size of the circles.
Figure 7 Time averaged radial atomic density distribution of the amorphous isomer of Na<sub>92</sub>, at some representative temperatures.
Figure 8 Square root of the diffusion coefficient as a function of temperature for the icosahedral and amorphous isomers of Na<sub>92</sub>.
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# Generalized Jaynes-Cummings Model with Intensity-Dependent and Non-resonant Coupling
## I Introduction
In the preceding paper we extended our earlier work describing interactions between a two-level system and a shape-invariant system. Here we present another generalization.
The model studied in Ref. is a generalization of the Jaynes-Cummings model . In the standard Jaynes-Cummings model the “field” is described by a harmonic oscillator. In our generalization it can be described by any shape-invariant system. In developing this model we made extensive use of the algebraic approach to the supersymmetric quantum mechanics . In this paper we further generalize the model to one with intensity-dependent interactions.
The standard Jaynes-Cummings model is an idealized model describing the interaction of matter with electromagnetic radiation. A variant of the Jaynes-Cummings model takes the coupling between matter and the radiation to depend on the intensity of the electromagnetic field . This model has great relevance since this kind of interaction means that the coupling is proportional to the amplitude of the field which is a very simple case of a nonlinear interaction corresponding to a more realistic physical situation. The results of this model can also give insight into the behavior of other quantum systems with strong nonlinear interactions.
## II The Generalized Intensity-Dependent and Non-resonant Jaynes-Cummings Hamiltonian
The expression of the intensity-dependent and non-resonant Jaynes-Cummings Hamiltonian can be written as
$$\widehat{𝐇}=\widehat{A}^{}\widehat{A}+\frac{1}{2}[\widehat{A},\widehat{A}^{}]\left(\widehat{\sigma }_3+1\right)+\alpha \left(\widehat{\sigma }_+\widehat{A}\sqrt{\widehat{A}^{}\widehat{A}}+\widehat{\sigma }_{}\sqrt{\widehat{A}^{}\widehat{A}}\widehat{A}^{}\right)+\mathrm{}\mathrm{\Delta }\widehat{\sigma }_3,$$
(1)
where $`\alpha `$ is a constant related with the coupling strength, $`\mathrm{\Delta }`$ is a constant related with the detuning of the system and $`\widehat{\sigma }_i`$, with $`i=1,\mathrm{\hspace{0.17em}2},\mathrm{and}\mathrm{\hspace{0.17em}\hspace{0.17em}3}`$, are the Pauli matrices.
However, the harmonic oscillator systems, used in this context, is only the simplest example of supersymmetric and shape-invariant potential. Our goal at this point is to generalize that Hamiltonian for all supersymmetric and shape-invariant systems. With this purpose we introduce the operators
$`\widehat{𝐒}=\widehat{\sigma }_+\widehat{A}+\widehat{\sigma }_{}\widehat{A}^{}`$ (3)
$`\widehat{𝐒}_i=\widehat{\sigma }_+\widehat{A}\sqrt{\widehat{A}^{}\widehat{A}}+\widehat{\sigma }_{}\sqrt{\widehat{A}^{}\widehat{A}}\widehat{A}^{},`$ (4)
where
$$\widehat{\sigma }_\pm =\frac{1}{2}\left(\widehat{\sigma }_1\pm i\widehat{\sigma }_2\right),$$
(5)
and, now, the operators $`\widehat{A}`$ and $`\widehat{A}^{}`$ satisfy the shape invariance condition . Using this definition we can decompose the Jaynes-Cummings Hamiltonian in the form
$$\widehat{𝐇}=\widehat{𝐇}_o+\widehat{𝐇}_{int},$$
(6)
where
$`\widehat{𝐇}_o=\widehat{𝐒}^2,`$ (8)
$`\widehat{𝐇}_{int}=\alpha \widehat{𝐒}_i+\mathrm{}\mathrm{\Delta }\widehat{\sigma }_3.`$ (9)
We search for the eigenstates of $`\widehat{𝐇}`$ and, in this case, it is more convenient to work with its $`B`$-operator expressions, which can be written as
$`\widehat{𝐒}^2=\left[\begin{array}{cc}\widehat{T}\widehat{B}_{}\widehat{B}_+\widehat{T}^{}& 0\\ 0& \widehat{B}_+\widehat{B}_{}\end{array}\right]\left[\begin{array}{cc}\widehat{H}_2& 0\\ 0& \widehat{H}_1\end{array}\right]`$ (11)
$`\widehat{𝐇}_{int}=\alpha \left[\begin{array}{cc}\beta & \widehat{T}\widehat{B}_{}\sqrt{\widehat{B}_+\widehat{B}_{}}\\ \sqrt{\widehat{B}_+\widehat{B}_{}}\widehat{B}_+\widehat{T}^{}& \beta \end{array}\right]=\alpha \left[\begin{array}{cc}\beta & \widehat{T}\widehat{B}_{}\sqrt{\widehat{H}_1}\\ \sqrt{\widehat{H}_1}\widehat{B}_+\widehat{T}^{}& \beta \end{array}\right],`$ (12)
where $`\beta =\mathrm{}\mathrm{\Delta }/\alpha `$. We use the same notation as the preceding paper . There we show that the states
$$\mathrm{\Psi }_m^{(\pm )}=\left[\begin{array}{cc}\widehat{T}& 0\\ 0& \pm 1\end{array}\right]\left[\begin{array}{c}C_m^{(\pm )}m\\ C_{m+1}^{(\pm )}m+1\end{array}\right],m=0,1,2,\mathrm{}$$
(13)
are the eigenstates of the operator $`\widehat{𝐒}^2`$
$$\widehat{𝐒}^2\mathrm{\Psi }_m^{(\pm )}=_{m+1}\mathrm{\Psi }_m^{(\pm )}.$$
(14)
In this case $`C_{m,m+1}^{(\pm )}C_{m,m+1}^{(\pm )}[R(a_1),R(a_2),R(a_3),\mathrm{}]`$ are auxiliary coefficients and, $`m`$ and $`m+1`$ are the abbreviated notation for the states $`\psi _m`$ and $`\psi _{m+1}`$ .
At this point, we observe that the wave-state orthonormalization conditions imply in the following relations among the $`C`$’s real coefficients
$`\mathrm{\Psi }_m^{(\pm )}\mathrm{\Psi }_m^{(\pm )}`$ $`=`$ $`\left[C_m^{(\pm )}\right]^2+\left[C_{m+1}^{(\pm )}\right]^2=1`$ (16)
$`\mathrm{\Psi }_m^{()}\mathrm{\Psi }_m^{(\pm )}`$ $`=`$ $`C_m^{(\pm )}C_m^{()}C_{m+1}^{(\pm )}C_{m+1}^{()}=0.`$ (17)
Now, considering that $`\widehat{𝐒}^2`$ and $`\widehat{𝐇}_{int}`$ commute then it is possible to find a common set of eigenstates. We can use this fact to determine the eigenvalues of $`\widehat{𝐇}_{int}`$ and the relations among the $`C`$’s coefficients. For that we need to calculate
$$\widehat{𝐇}_{int}\mathrm{\Psi }_m^{(\pm )}=\lambda _m^{(\pm )}\mathrm{\Psi }_m^{(\pm )},$$
(18)
where $`\lambda _m^{(\pm )}`$ are the eigenvalues to be determined. Using the Eqs. (II), (6) and (13), the last eigenvalue equation can be rewritten in a matrix form as
$$\alpha \left[\begin{array}{cc}\beta & \widehat{T}\widehat{B}_{}\sqrt{\widehat{H}_1}\\ \sqrt{\widehat{H}_1}\widehat{B}_+\widehat{T}^{}& \beta \end{array}\right]\left[\begin{array}{cc}\widehat{T}& 0\\ 0& \pm 1\end{array}\right]\left[\begin{array}{c}C_m^{(\pm )}m\\ C_{m+1}^{(\pm )}m+1\end{array}\right]=\lambda _m^{(\pm )}\left[\begin{array}{c}C_m^{(\pm )}m\\ C_{m+1}^{(\pm )}m+1\end{array}\right],$$
(19)
Since the $`C`$’s coefficients commute with the $`\widehat{A}`$ or $`\widehat{A}^{}`$ operators, then the last matrix equation permits to obtain the following equations
$`\left[\alpha \beta \lambda _m^{(\pm )}\right]\left(\widehat{T}C_m^{(\pm )}\widehat{T}^{}\right)\widehat{T}m\pm \alpha C_{m+1}^{(\pm )}\widehat{T}\widehat{B}_{}\sqrt{\widehat{H}_1}m+1=0`$ (21)
$`\alpha \left(\widehat{T}C_m^{(\pm )}\widehat{T}^{}\right)\sqrt{\widehat{H}_1}\widehat{B}_+m\left[\alpha \beta +\lambda _m^{(\pm )}\right]C_{m+1}^{(\pm )}m+1=0.`$ (22)
Introducing the operator
$$\widehat{Q}^{}=\left(\widehat{B}_+\widehat{B}_{}\right)^{1/2}\widehat{B}_+=\left(\widehat{H}_1\right)^{1/2}\widehat{B}_+$$
(23)
one can write the normalized eigenstate of $`\widehat{H}_1`$ as
$$m=\left(\widehat{Q}^{}\right)^m0,$$
(24)
and, with Eqs. (23) and (24) we can show that
$`\widehat{B}_+m=\sqrt{_{m+1}}m+1,`$ (26)
$`\widehat{T}\widehat{B}_{}m+1=\sqrt{_{m+1}}\widehat{T}m.`$ (27)
Therefore, we have that
$`\widehat{T}\widehat{B}_{}\sqrt{\widehat{H}_1}m+1`$ $`=`$ $`\widehat{T}\widehat{B}_{}\sqrt{_{m+1}}m+1`$ (28)
$`=`$ $`\sqrt{_{m+1}}\widehat{T}\widehat{B}_{}m+1`$ (29)
$`=`$ $`_{m+1}\widehat{T}m,`$ (30)
and
$`\sqrt{\widehat{H}_1}\widehat{B}_+m`$ $`=`$ $`\sqrt{\widehat{H}_1}\sqrt{_{m+1}}m+1`$ (31)
$`=`$ $`\sqrt{_{m+1}}\sqrt{\widehat{H}_1}m+1`$ (32)
$`=`$ $`_{m+1}m+1,`$ (33)
Using Eqs. (30) and (33), then Eqs. (19) take the form
$`\left\{\left[\alpha \beta \lambda _m^{(\pm )}\right]\left(\widehat{T}C_m^{(\pm )}\widehat{T}^{}\right)\pm \alpha _{m+1}C_{m+1}^{(\pm )}\right\}\widehat{T}m=0`$ (35)
$`\left\{\alpha _{m+1}\left(\widehat{T}C_m^{(\pm )}\widehat{T}^{}\right)\left[\alpha \beta +\lambda _m^{(\pm )}\right]C_{m+1}^{(\pm )}\right\}m+1=0.`$ (36)
From Eqs. (II) it follows that
$$\lambda _m^{(\pm )}=\pm \alpha \sqrt{_{m+1}^2+\beta ^2},$$
(37)
and
$$C_{m+1}^{(\pm )}=\left(\frac{\sqrt{_{m+1}^2+\beta ^2}\beta }{_{m+1}}\right)\left(\widehat{T}C_m^{(\pm )}\widehat{T}^{}\right).$$
(38)
Eqs. (II) and (38) imply that
$$C_{m+1}^{(\pm )}=C_m^{()},$$
(39)
and the eigenstates and eigenvalues of the generalized intensity-dependent and non-resonant Jaynes-Cummings Hamiltonian can be written as
$$E_m^{(\pm )}=_{m+1}\pm \sqrt{\alpha ^2_{m+1}^2+\mathrm{}^2\mathrm{\Delta }^2},$$
(40)
and
$$\mathrm{\Psi }_m^{(\pm )}=\left[\begin{array}{cc}\widehat{T}& 0\\ 0& \pm 1\end{array}\right]\left[\begin{array}{c}C_m^{(\pm )}m\\ C_m^{()}m+1\end{array}\right],m=0,1,2,\mathrm{}$$
(41)
a) The Intensity-Dependent Resonant Limit
From these general results we can verify two simple limiting cases. The first one corresponds to the resonant situation, which is for $`\mathrm{\Delta }=0`$ $`(\beta =0)`$. Using these conditions into Eqs. (38) and (40) and Eqs. (II) we can promptly conclude that
$$E_m^{(\pm )}=\left(1\pm \alpha \right)_{m+1},$$
(42)
and
$$C_{m+1}^{(\pm )}=\widehat{T}C_m^{(\pm )}\widehat{T}^{}=C_m^{(\pm )}=\frac{1}{\sqrt{2}}.$$
(43)
Therefore the intensity-dependent Jaynes-Cummings resonant eigenstate is given by
$$\mathrm{\Psi }_m^{(\pm )}=\frac{1}{\sqrt{2}}\left[\begin{array}{cc}\widehat{T}& 0\\ 0& \pm 1\end{array}\right]\left[\begin{array}{c}m\\ m+1\end{array}\right],m=0,1,2,\mathrm{}$$
(44)
If we compare this last particular result with that one found in the reference , we conclude that the intensity-dependent and intensity-independent generalized Jaynes-Cummings Hamiltonians have the same eigenstates in the resonant situation.
b) The Standard Intensity-Dependent Jaynes-Cummings Limit
The second important limit corresponds to the standard intensity-dependent Jaynes-Cummings case, related with the harmonic oscillator system. In this limit we have that $`\widehat{T}=\widehat{T}^{}1`$, $`\widehat{B}_{}\widehat{a}`$, $`\widehat{B}_+\widehat{a}^{}`$, $`\mathrm{\Delta }=\omega \omega _o`$ and $`_{m+1}=(m+1)\mathrm{}\omega `$. Using these conditions in Eqs. (38), (40) and Eqs. (II) we can promptly conclude that
$$E_m^{(\pm )}=(m+1)\mathrm{}\omega \pm \mathrm{}\sqrt{\alpha ^2\omega ^2(m+1)^2+(\omega \omega _o)^2},$$
(45)
and
$$C_{m+1}^{(\pm )}=\gamma _m^{(\pm )}C_m^{(\pm )}=C_m^{()}=\frac{1}{\sqrt{1+\left(\gamma _m^{()}\right)^2}},$$
(46)
where
$`\gamma _m^{(\pm )}=\sqrt{1+\delta _m^2}\delta _m,`$ (48)
$`\delta _m={\displaystyle \frac{\omega \omega _o}{\alpha \omega (m+1)}}.`$ (49)
Therefore the standard intensity-dependent Jaynes-Cummings eigenstate, written in a matrix form, is given by
$$\mathrm{\Psi }_m^{(\pm )}=\frac{1}{\sqrt{1+\left(\gamma _m^{(\pm )}\right)^2}}\left[\begin{array}{cc}1& 0\\ 0& \pm \gamma _m^{(\pm )}\end{array}\right]\left[\begin{array}{c}m\\ m+1\end{array}\right],m=0,1,2,\mathrm{}.$$
(50)
## III Time Evolution of the System
To resolve the the time-dependent Schrödinger equation for intensity-dependent and non-resonant Jaynes-Cummings systems
$$i\mathrm{}\frac{}{t}\mathrm{\Psi }(t)=\left(\widehat{𝐇}_o+\widehat{𝐇}_{int}\right)\mathrm{\Psi }(t)$$
(51)
we can write the state as
$$\mathrm{\Psi }(t)=\mathrm{exp}\left(i\widehat{𝐇}_ot/\mathrm{}\right)\mathrm{\Psi }_i(t),$$
(52)
and, by substituting this into Schrödinger equation and taking into account the commutation property between $`\widehat{𝐇}_o`$ and $`\widehat{𝐇}_{int}`$, we obtain
$$i\mathrm{}\frac{}{t}\mathrm{\Psi }_i(t)=\widehat{𝐇}_{int}\mathrm{\Psi }_i(t),$$
(53)
Now, we can introduce the evolution matrix $`\widehat{𝐔}_i(t,0)`$, related with the interaction Hamiltonian, by
$$\mathrm{\Psi }_i(t)=\widehat{𝐔}_i(t,0)\mathrm{\Psi }_i(0).$$
(54)
with
$$i\mathrm{}\frac{}{t}\widehat{𝐔}_i(t,0)=\widehat{𝐇}_{int}\widehat{𝐔}_i(t,0),$$
(55)
that is, in matrix form, written as
$$i\mathrm{}\left[\begin{array}{cc}\widehat{U}_{11}^{}& \widehat{U}_{12}^{}\\ \widehat{U}_{21}^{}& \widehat{U}_{22}^{}\end{array}\right]=\alpha \left[\begin{array}{cc}\beta & \widehat{T}\widehat{B}_{}\sqrt{\widehat{H}_1}\\ \sqrt{\widehat{H}_1}\widehat{B}_+\widehat{T}^{}& \beta \end{array}\right]\left[\begin{array}{cc}\widehat{U}_{11}& \widehat{U}_{12}\\ \widehat{U}_{21}& \widehat{U}_{22}\end{array}\right],$$
(56)
where the primes denote the time derivative. One fast way to diagonalize the evolution matrix differential equation is by differentiating Eq. (55) with respect to time. After that, if we use again the same Eq. (55), we find
$$i\mathrm{}\frac{^2}{t^2}\widehat{𝐔}_i(t,0)=\widehat{𝐇}_{int}\frac{}{t}\widehat{𝐔}_i(t,0)=\frac{1}{i\mathrm{}}\widehat{𝐇}_{int}^2\widehat{𝐔}_i(t,0),$$
(57)
which can be written as
$$\left[\begin{array}{cc}\widehat{U}_{11}^{\prime \prime }& \widehat{U}_{12}^{\prime \prime }\\ \widehat{U}_{21}^{\prime \prime }& \widehat{U}_{22}^{\prime \prime }\end{array}\right]=\left[\begin{array}{cc}\widehat{\omega }_1& 0\\ 0& \widehat{\omega }_2\end{array}\right]\left[\begin{array}{cc}\widehat{U}_{11}& \widehat{U}_{12}\\ \widehat{U}_{21}& \widehat{U}_{22}\end{array}\right],$$
(58)
where
$`\mathrm{}\widehat{\omega }_1=\alpha \sqrt{(\widehat{T}\widehat{B}_{}\widehat{B}_+\widehat{T}^{})^2+\beta ^2}=\sqrt{\alpha ^2\widehat{H}_2^2+(\mathrm{}\mathrm{\Delta })^2},`$ (60)
$`\mathrm{}\widehat{\omega }_2=\alpha \sqrt{(\widehat{B}_+\widehat{B}_{})^2+\beta ^2}=\sqrt{\alpha ^2\widehat{H}_1^2+(\mathrm{}\mathrm{\Delta })^2}.`$ (61)
Now, since by initial conditions $`\widehat{𝐔}_i(0,0)=\widehat{𝐈}`$, then we can write the solution of the evolution matrix differential equation (57) as
$$\widehat{𝐔}_i(t,0)=\left[\begin{array}{cc}\mathrm{cos}(\widehat{\omega }_1t)& \mathrm{sin}(\widehat{\omega }_1t)\widehat{C}\\ \mathrm{sin}(\widehat{\omega }_2t)\widehat{D}& \mathrm{cos}(\widehat{\omega }_2t)\end{array}\right],$$
(62)
and the $`\widehat{C}`$ and $`\widehat{D}`$ operators can be determined by the unitary transformation conditions
$$\widehat{𝐔}_i^{}(t,0)\widehat{𝐔}_i(t,0)=\widehat{𝐔}_i(t,0)\widehat{𝐔}_i^{}(t,0)=\widehat{𝐈}.$$
(63)
Following the same steps used in the appendix A of the reference we can conclude that to satisfy the unitary conditions (63) these operators must have the form
$`\widehat{C}=\widehat{D}^{}={\displaystyle \frac{i}{\widehat{H}_2^{1/4}}}\sqrt{\widehat{T}\widehat{B}_{}}`$ (65)
$`\widehat{D}=\widehat{C}^{}=\sqrt{\widehat{B}_+\widehat{T}^{}}{\displaystyle \frac{i}{\widehat{H}_2^{1/4}}}.`$ (66)
Therefore, we can write the final expression of the time evolution matrix of the system as
$$\widehat{𝐔}_i(t,0)=\left[\begin{array}{cc}\mathrm{cos}(\widehat{\omega }_1t)& \mathrm{sin}(\widehat{\omega }_1t)\widehat{C}\\ \mathrm{sin}(\widehat{\omega }_2t)\widehat{C}^{}& \mathrm{cos}(\widehat{\omega }_2t)\end{array}\right].$$
(67)
For Jaynes-Cummings systems an important physical quantity to see how the system under consideration evolves in time is the population inversion factor , defined by
$$\widehat{𝐖}(t)\widehat{\sigma }_+(t)\widehat{\sigma }_{}(t)\widehat{\sigma }_{}(t)\widehat{\sigma }_+(t)=\widehat{\sigma }_3(t),$$
(68)
where the time dependence of the operators is related with the Heisenberg picture. In this case, the time evolution of the population inversion factor will be given by
$$\frac{d\widehat{\sigma }_3(t)}{dt}=\frac{1}{i\mathrm{}}\widehat{𝐔}_i^{}(t,0)[\widehat{\sigma }_3,\widehat{𝐇}]\widehat{𝐔}_i(t,0),$$
(69)
and since we have
$$[\widehat{\sigma }_3,\widehat{𝐇}]=\alpha [\widehat{\sigma }_3,\widehat{𝐒}_i]=2\alpha \widehat{𝐒}_i\widehat{\sigma }_3,$$
(70)
then Eq. (69) can be written as
$$\frac{d\widehat{\sigma }_3(t)}{dt}=\frac{2i\alpha }{\mathrm{}}\widehat{𝐒}_i(t)\widehat{\sigma }_3(t).$$
(71)
We can obtain a differential equation with constant coefficients for $`\widehat{\sigma }_3(t)`$ by taking the time derivative of Eq. (71)
$$\frac{d^2\widehat{\sigma }_3(t)}{dt^2}=\frac{2i\alpha }{\mathrm{}}\left\{\frac{d\widehat{𝐒}_i(t)}{dt}\widehat{\sigma }_3(t)+\widehat{𝐒}_i(t)\frac{d\widehat{\sigma }_3(t)}{dt}\right\}.$$
(72)
Having in mind that
$$\frac{d\widehat{𝐒}_i(t)}{dt}=\frac{1}{i\mathrm{}}\widehat{𝐔}_i^{}(t,0)[\widehat{𝐒}_i,\widehat{𝐇}]\widehat{𝐔}_i(t,0),$$
(73)
and,
$$[\widehat{𝐒}_i,\widehat{𝐇}]=\alpha \beta [\widehat{𝐒}_i,\widehat{\sigma }_3]=2\alpha \beta \widehat{𝐒}_i\widehat{\sigma }_3,$$
(74)
we can conclude that
$$\frac{d\widehat{𝐒}_i(t)}{dt}=\frac{2i\alpha \beta }{\mathrm{}}\widehat{𝐒}_i(t)\widehat{\sigma }_3(t).$$
(75)
Now using Eqs. (71) and (75) into Eq. (72) we obtain
$$\frac{d^2\widehat{\sigma }_3(t)}{dt^2}+\widehat{𝚯}^2\widehat{\sigma }_3(t)=\widehat{𝐅}(t)$$
(76)
where
$`\widehat{𝚯}^2={\displaystyle \frac{4\alpha ^2}{\mathrm{}^2}}\widehat{𝐒}_i^2`$ (78)
$`\widehat{𝐅}(t)={\displaystyle \frac{4\alpha ^2\beta }{\mathrm{}^2}}\widehat{𝐔}_i^{}(t,0)\widehat{𝐒}_i\widehat{𝐔}_i(t,0).`$ (79)
The Eq. (75) corresponds to a non-homogeneous linear differential equation for $`\widehat{\sigma }_3(t)`$ with constant coefficients since $`\widehat{𝐒}_i^2`$ and $`\widehat{𝐇}`$ commute and, therefore, $`\widehat{𝚯}`$ is a constant of the motion. The general solution of this differential equation can be written as
$$\widehat{\sigma }_3(t)=\widehat{\sigma }^H(t)+\widehat{\sigma }^P(t),$$
(80)
and each matrix element of the homogeneous solution, satisfies the differential equation
$$\frac{d^2\widehat{\sigma }_{jk}^H(t)}{dt^2}+\widehat{\nu }_j^2\widehat{\sigma }_{jk}^H(t)=0,j,k=1,\mathrm{or}\mathrm{\hspace{0.33em}2},$$
(81)
with
$`\mathrm{}\widehat{\nu }_1=2\alpha \widehat{T}\widehat{B}_{}\widehat{B}_+\widehat{T}^{}=2\alpha \widehat{H}_2,`$ (83)
$`\mathrm{}\widehat{\nu }_2=2\alpha \widehat{B}_+\widehat{B}_{}=2\alpha \widehat{H}_1.`$ (84)
The solution of Eq. (81) is given by
$$\widehat{\sigma }_{jk}^H(t)=\widehat{y}_j(t)\widehat{c}_{jk}+\widehat{z}_j(t)\widehat{d}_{jk},$$
(85)
where
$`\widehat{y}_j(t)=\mathrm{cos}(\widehat{\nu }_jt)`$ (87)
$`\widehat{z}_j(t)=\mathrm{sin}(\widehat{\nu }_jt),`$ (88)
and the coefficients $`\widehat{c}_{jk}`$ and $`\widehat{d}_{jk}`$ can be determined by the initial conditions.
The matrix elements of the particular solution of the $`\widehat{\sigma }_3(t)`$ differential equation needs to satisfy
$$\frac{d^2\widehat{\sigma }_{jk}^P(t)}{dt^2}+\widehat{\nu }_j^2\widehat{\sigma }_{jk}^P(t)=\widehat{F}_{jk}(t),j,k=1,\mathrm{or}\mathrm{\hspace{0.33em}2},$$
(89)
and can be obtained by the variation of parameter or by Green function methods, giving
$$\widehat{\sigma }_{jk}^P(t)=\widehat{\nu }_j^1\left\{\widehat{z}_j(t)_0^t\xi \widehat{y}_j(\xi )\widehat{F}_{jk}(\xi )\widehat{y}_j(t)_0^t𝑑\xi \widehat{z}_j(\xi )\widehat{F}_{jk}(\xi )\right\},$$
(90)
where we used that the Wronskian of the system of solutions $`\widehat{y}_j(t)`$ and $`\widehat{z}_j(t)`$ is given by $`\widehat{\nu }_j`$.
After we determine the elements of the $`\widehat{𝐅}(t)`$-matrix, it is necessary to resolve the integrals in Eq. (90) to obtain the explicit expression of the particular solution. In the appendix we show that, using Eqs. (II), (67), and (76), it is possible to conclude that these matrix elements can be written as
$`\widehat{\sigma }_{11}^P(t)`$ $`=`$ $`i{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_1^1\sqrt{\widehat{T}\widehat{B}_{}}\left\{\widehat{z}_2(t)𝒢_{CS}^{(+)}(t;\widehat{\nu }_2,\widehat{\omega }_2,\widehat{\omega }_1)\widehat{y}_2(t)𝒢_{SS}^{(+)}(t;\widehat{\nu }_2,\widehat{\omega }_2,\widehat{\omega }_1)\right\}\widehat{H}_2^{3/4}`$ (92)
$`+`$ $`i{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_1^1\widehat{H}_2^{3/4}\left\{\widehat{z}_1(t)𝒢_{SC}^{()}(t;\widehat{\nu }_1,\widehat{\omega }_1,\widehat{\omega }_2)\widehat{y}_1(t)𝒢_{CC}^{()}(t;\widehat{\nu }_1,\widehat{\omega }_1,\widehat{\omega }_2)\right\}\sqrt{\widehat{B}_+\widehat{T}^{}},`$ (93)
$`\widehat{\sigma }_{12}^P(t)`$ $`=`$ $`{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_1^1\sqrt{\widehat{T}\widehat{B}_{}}\left\{\widehat{z}_2(t)𝒢_{CC}^{(+)}(t;\widehat{\nu }_2,\widehat{\omega }_2,\widehat{\omega }_1)\widehat{y}_2(t)𝒢_{SC}^{(+)}(t;\widehat{\nu }_2,\widehat{\omega }_2,\widehat{\omega }_1)\right\}\sqrt{\widehat{H}_2\widehat{T}\widehat{B}_{}}`$ (94)
$`+`$ $`{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_1^1\widehat{H}_2^{3/4}\left\{\widehat{z}_1(t)𝒢_{SS}^{()}(t;\widehat{\nu }_1,\widehat{\omega }_1,\widehat{\omega }_2)+\widehat{y}_1(t)𝒢_{CS}^{()}(t;\widehat{\nu }_1,\widehat{\omega }_1,\widehat{\omega }_2)\right\}\widehat{H}_1^{1/4},`$ (95)
$`\widehat{\sigma }_{21}^P(t)`$ $`=`$ $`{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_2^1\sqrt{\widehat{B}_+\widehat{T}^{}\widehat{H}_2}\left\{\widehat{z}_1(t)𝒢_{CC}^{(+)}(t;\widehat{\nu }_1,\widehat{\omega }_1,\widehat{\omega }_2)\widehat{y}_1(t)𝒢_{SC}^{(+)}(t;\widehat{\nu }_1,\widehat{\omega }_1,\widehat{\omega }_2)\right\}\sqrt{\widehat{B}_+\widehat{T}^{}}`$ (96)
$`+`$ $`{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_2^1\widehat{H}_1^{1/4}\left\{\widehat{z}_2(t)𝒢_{SS}^{()}(t;\widehat{\nu }_2,\widehat{\omega }_2,\widehat{\omega }_1)\widehat{y}_2(t)𝒢_{CS}^{()}(t;\widehat{\nu }_2,\widehat{\omega }_2,\widehat{\omega }_1)\right\}\widehat{H}_2^{3/4},`$ (97)
$`\widehat{\sigma }_{22}^P(t)`$ $`=`$ $`i{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_2^1\sqrt{\widehat{B}_+\widehat{T}^{}\widehat{H}_2}\left\{\widehat{z}_1(t)𝒢_{CS}^{(+)}(t;\widehat{\nu }_1,\widehat{\omega }_1,\widehat{\omega }_2)\widehat{y}_1(t)𝒢_{SS}^{(+)}(t;\widehat{\nu }_1,\widehat{\omega }_1,\widehat{\omega }_2)\right\}\widehat{H}_1^{1/4}`$ (98)
$`+`$ $`i{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_2^1\widehat{H}_1^{1/4}\left\{\widehat{z}_2(t)𝒢_{SC}^{()}(t;\widehat{\nu }_2,\widehat{\omega }_2,\widehat{\omega }_1)+\widehat{y}_2(t)𝒢_{CC}^{()}(t;\widehat{\nu }_2,\widehat{\omega }_2,\widehat{\omega }_1)\right\}\sqrt{\widehat{H}_2\widehat{T}\widehat{B}_{}},`$ (99)
where $`\gamma =4\alpha ^2\beta /\mathrm{}^2`$, and the auxiliary functions are given by
$$𝒢_{XY}^{(\pm )}(t;\widehat{p},\widehat{q},\widehat{r})=_{XY}(t;\widehat{p}\widehat{q},\widehat{r})\pm _{XY}(t;\widehat{p}+\widehat{q},\widehat{r}),X,Y=C\mathrm{or}S,$$
(100)
with
$`_{CC}(t;\widehat{x},\widehat{w})`$ $``$ $`{\displaystyle _0^t}𝑑\xi \mathrm{cos}(\widehat{x}\xi )\mathrm{cos}(\widehat{w}\xi )`$ (102)
$`=`$ $`{\displaystyle \underset{m,n=0}{\overset{\mathrm{}}{}}}(1)^{m+n}{\displaystyle \frac{\widehat{x}^{2m}\widehat{w}^{2n}}{(2m)!(2n)!}}{\displaystyle \frac{t^{2m+2n+1}}{(2m+2n+1)}}`$ (103)
$`_{CS}(t;\widehat{x},\widehat{w})`$ $``$ $`{\displaystyle _0^t}𝑑\xi \mathrm{cos}(\widehat{x}\xi )\mathrm{sin}(\widehat{w}\xi )`$ (104)
$`=`$ $`{\displaystyle \underset{m,n=0}{\overset{\mathrm{}}{}}}(1)^{m+n}{\displaystyle \frac{\widehat{x}^{2m}\widehat{w}^{2n+1}}{(2m)!(2n+1)!}}{\displaystyle \frac{t^{2m+2n+2}}{(2m+2n+2)}}`$ (105)
$`_{SC}(t;\widehat{x},\widehat{w})`$ $``$ $`{\displaystyle _0^t}𝑑\xi \mathrm{sin}(\widehat{x}\xi )\mathrm{cos}(\widehat{w}\xi )`$ (106)
$`=`$ $`{\displaystyle \underset{m,n=0}{\overset{\mathrm{}}{}}}(1)^{m+n}{\displaystyle \frac{\widehat{x}^{2m+1}\widehat{w}^{2n}}{(2m+1)!(2n)!}}{\displaystyle \frac{t^{2m+2n+2}}{(2m+2n+2)}}`$ (107)
$`_{SS}(t;\widehat{x},\widehat{w})`$ $``$ $`{\displaystyle _0^t}𝑑\xi \mathrm{sin}(\widehat{x}\xi )\mathrm{sin}(\widehat{w}\xi )`$ (108)
$`=`$ $`{\displaystyle \underset{m,n=0}{\overset{\mathrm{}}{}}}(1)^{m+n}{\displaystyle \frac{\widehat{x}^{2m+1}\widehat{w}^{2n+1}}{(2m+1)!(2n+1)!}}{\displaystyle \frac{t^{2m+2n+3}}{(2m+2n+3)}}.`$ (109)
With these results for the particular solution we can conclude that
$$\widehat{\sigma }_{ij}^P(0)=0=\frac{d\widehat{\sigma }_{ij}^P(0)}{dt}.$$
(110)
Now, using Eqs. (71), (80), (85) and the initial conditions, we have
$`\left[\widehat{\sigma }_3(0)\right]_{ij}=\widehat{c}_{ij}`$ (112)
$`\left[{\displaystyle \frac{d\widehat{\sigma }_3(0)}{dt}}\right]_{ij}={\displaystyle \frac{2i\alpha }{\mathrm{}}}\left[\widehat{𝐒}_i(0)\widehat{\sigma }_3(0)\right]_{ij}=\widehat{\nu }_i\widehat{d}_{ij}.`$ (113)
Therefore, the final expression for the elements of the population inversion matrix of the system can be written as
$$[\widehat{\sigma }_3(t)]_{ij}=\mathrm{cos}(\widehat{\nu }_it)\left[\widehat{\sigma }_3(0)\right]_{ij}+\frac{2i\alpha }{\mathrm{}}\mathrm{sin}(\widehat{\nu }_it)\widehat{\nu }_i^1\left[\widehat{𝐒}(0)_i\widehat{\sigma }_3(0)\right]_{ij}+\widehat{\sigma }_{ij}^P(t).$$
(114)
Again, using these final results we can verify two limiting cases.
a) The Intensity-Dependent Resonant Limit
The first one corresponds to the intensity-dependent resonant $`(\mathrm{\Delta }=0)`$. Eqs. (58), (67), (81) and (III) allow us to conclude that, in this case, the evolution matrix of the system is given by
$$\widehat{𝐔}_i(t,0)=\left[\begin{array}{cc}\mathrm{cos}\left(\frac{1}{2}\widehat{\nu }_1t\right)& \mathrm{sin}\left(\frac{1}{2}\widehat{\nu }_1t\right)\widehat{C}\\ \mathrm{sin}\left(\frac{1}{2}\widehat{\nu }_2t\right)\widehat{C}^{}& \mathrm{cos}\left(\frac{1}{2}\widehat{\nu }_2t\right)\end{array}\right].$$
(115)
and the elements of the population inversion of the system is
$$[\widehat{\sigma }_3(t)]_{ij}=\mathrm{cos}(\widehat{\nu }_it)\left[\widehat{\sigma }_3(0)\right]_{ij}+\frac{2i\alpha }{\mathrm{}}\mathrm{sin}(\widehat{\nu }_it)\widehat{\nu }_i^1\left[\widehat{𝐒}_i(0)\widehat{\sigma }_3(0)\right]_{ij}.$$
(116)
b) The Standard Intensity-Dependent Jaynes-Cummings Limit
This second important limit corresponds to the case of the harmonic oscillator system, and in this limit we have that $`\widehat{T}=\widehat{T}^{}1`$, $`\widehat{B}_{}\widehat{a}`$, $`\widehat{B}_+\widehat{a}^{}`$ and $`[\widehat{a},\widehat{a}^{}]=\mathrm{}\omega `$. With these conditions the operators $`\widehat{\omega }_1`$ and $`\widehat{\omega }_2`$ commute, and this fact permits to evaluate the integrals related with the particular solution of the population inversion elements using trigonometric product relations. Using that and the expressions obtained in the appendix, after a considerable amount of algebra and trigonometric product relations we can show that is possible to write the expressions for the $`\widehat{\sigma }_{ij}^P(t)`$-matrix elements as
$`\widehat{\sigma }_{11}^P(t)`$ $`=`$ $`i{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_1^1\sqrt{\widehat{a}}\left\{𝒦_S(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_2)𝒦_S(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_2)\right\}\left(\widehat{a}\widehat{a}^{}\right)^{3/4}`$ (118)
$``$ $`i{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_1^1\left(\widehat{a}\widehat{a}^{}\right)^{3/4}\left\{𝒦_S(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_1)𝒦_S(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_1)\right\}\sqrt{\widehat{a}^{}}`$ (119)
$`\widehat{\sigma }_{12}^P(t)`$ $`=`$ $`{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_1^1\sqrt{\widehat{a}}\left\{𝒦_C(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_2)𝒦_C(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_2)\right\}\sqrt{\widehat{a}\widehat{a}^{}\widehat{a}}`$ (120)
$``$ $`{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_1^1\left(\widehat{a}\widehat{a}^{}\right)^{3/4}\left\{𝒦_C(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_1)𝒦_C(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_1)\right\}\left(\widehat{a}^{}\widehat{a}\right)^{1/4}`$ (121)
$`\widehat{\sigma }_{21}^P(t)`$ $`=`$ $`{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_2^1\sqrt{\widehat{a}^{}\widehat{a}\widehat{a}^{}}\left\{𝒦_C(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_1)+𝒦_C(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_1)\right\}\sqrt{\widehat{a}^{}}`$ (122)
$``$ $`{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_2^1\left(\widehat{a}^{}\widehat{a}\right)^{1/4}\left\{𝒦_C(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_2)𝒦_C(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_2)\right\}\left(\widehat{a}\widehat{a}^{}\right)^{3/4}`$ (123)
$`\widehat{\sigma }_{22}^P(t)`$ $`=`$ $`i{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_2^1\sqrt{\widehat{a}^{}\widehat{a}\widehat{a}^{}}\left\{𝒦_S(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_1)+𝒦_S(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_1)\right\}\left(\widehat{a}^{}\widehat{a}\right)^{1/4}`$ (124)
$``$ $`i{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_2^1\left(\widehat{a}^{}\widehat{a}\right)^{1/4}\left\{𝒦_S(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_2)+𝒦_S(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_2)\right\}\sqrt{\widehat{a}\widehat{a}^{}\widehat{a}},`$ (125)
where, now, the auxiliary functions are given by
$`𝒦_S(t;\widehat{p},\widehat{q},\widehat{r})`$ $`=`$ $`{\displaystyle \frac{\widehat{r}\mathrm{sin}\left[\left(\widehat{p}+\widehat{q}\right)t\right]\left(\widehat{p}+\widehat{q}\right)\mathrm{sin}\left(\widehat{r}t\right)}{\widehat{r}^2\left(\widehat{p}+\widehat{q}\right)^2}}`$ (127)
$`𝒦_C(t;\widehat{p},\widehat{q},\widehat{r})`$ $`=`$ $`{\displaystyle \frac{\widehat{r}\mathrm{cos}\left[\left(\widehat{p}+\widehat{q}\right)t\right]\widehat{r}\mathrm{cos}\left(\widehat{r}t\right)}{\widehat{r}^2\left(\widehat{p}+\widehat{q}\right)^2}}.`$ (128)
Considering the expressions above we may easily verify that the particular solution for the population inversion factor must still satisfy the initial conditions (110). Therefore, in this case the final expression of the population inversion factor has the same form given by Eq. (114), with
$`\mathrm{}\widehat{\nu }_1=2\alpha \widehat{a}\widehat{a}^{},\mathrm{}\widehat{\nu }_2=2\alpha \widehat{a}^{}\widehat{a},`$ (130)
$`\mathrm{}\widehat{\omega }_1=\alpha \sqrt{(\widehat{a}\widehat{a}^{})^2+\beta ^2},\mathrm{}\widehat{\omega }_2=\alpha \sqrt{(\widehat{a}^{}\widehat{a})^2+\beta ^2}.`$ (131)
## IV Conclusions
In this article we introduced a class of shape-invariant bound-state problems which represent two-level systems. The corresponding coupled-channel Hamiltonians generalize the intensity-dependent and non-resonant Jaynes-Cummings Hamiltonian. These models are not only interesting on their own account. Being exactly solvable coupled-channels problems they may help to assess the validity and accuracy of various approximate approaches to the coupled-channel problems .
## ACKNOWLEDGMENTS
This work was supported in part by the U.S. National Science Foundation Grants No. PHY-9605140 and PHY-0070161 at the University of Wisconsin, and in part by the University of Wisconsin Research Committee with funds granted by the Wisconsin Alumni Research Foundation. A.B.B. acknowledges the support of the Alexander von Humboldt-Stiftung. M.A.C.R. acknowledges the support of Fundação de Amparo à Pesquisa do Estado de São Paulo (Contract No. 98/13722-2). A.N.F.A. acknowledges the support of Fundação Coordenação de Aperfeiçoamento de Pessoal de Nível Superior (Contract No. BEX0610/96-8). A.B.B. is grateful to the Max-Planck-Institut für Kernphysik and H.A. Weidenmüller for the very kind hospitality.
Appendix
In this appendix we show the steps necessary to obtain the explicit expression of the elements of the population inversion particular solution. To resolve the integrals in the Eq. (90), first we need to determine the elements of the $`\widehat{𝐅}(t)`$-matrix. To do that we can use Eqs. (II), (76), and (67) to conclude that
$`\widehat{F}_{11}(t)`$ $`=`$ $`\gamma \left\{\mathrm{cos}(\widehat{\omega }_1t)\widehat{T}\widehat{B}_{}\mathrm{sin}(\widehat{\omega }_2t)\widehat{C}^{}+\widehat{C}\mathrm{sin}(\widehat{\omega }_2t)\widehat{B}_+\widehat{T}^{}\mathrm{cos}(\widehat{\omega }_1t)\right\}`$ (133)
$`=`$ $`i\gamma \left\{\sqrt{\widehat{T}\widehat{B}_{}}\mathrm{cos}(\widehat{\omega }_2t)\mathrm{sin}(\widehat{\omega }_1t)\widehat{H}_2^{1/4}\widehat{H}_2^{1/4}\mathrm{sin}(\widehat{\omega }_1t)\mathrm{cos}(\widehat{\omega }_2t)\sqrt{\widehat{B}_+\widehat{T}^{}}\right\}`$ (134)
$`\widehat{F}_{12}(t)`$ $`=`$ $`\gamma \left\{\mathrm{cos}(\widehat{\omega }_1t)\widehat{T}\widehat{B}_{}\mathrm{cos}(\widehat{\omega }_2t)\widehat{C}\mathrm{sin}(\widehat{\omega }_2t)\widehat{B}_+\widehat{T}^{}\mathrm{sin}(\widehat{\omega }_1t)\widehat{C}\right\}`$ (135)
$`=`$ $`\gamma \left\{\sqrt{\widehat{T}\widehat{B}_{}}\mathrm{cos}(\widehat{\omega }_2t)\mathrm{cos}(\widehat{\omega }_1t)\sqrt{\widehat{T}\widehat{B}_{}}+\widehat{H}_2^{1/4}\mathrm{sin}(\widehat{\omega }_1t)\mathrm{sin}(\widehat{\omega }_2t)\widehat{H}_1^{1/4}\right\}`$ (136)
$`\widehat{F}_{21}(t)`$ $`=`$ $`\gamma \left\{\mathrm{cos}(\widehat{\omega }_2t)\widehat{B}_+\widehat{T}^{}\mathrm{cos}(\widehat{\omega }_1t)\widehat{C}^{}\mathrm{sin}(\widehat{\omega }_1t)\widehat{T}\widehat{B}_{}\mathrm{sin}(\widehat{\omega }_2t)\widehat{C}^{}\right\}`$ (137)
$`=`$ $`\gamma \left\{\sqrt{\widehat{B}_+\widehat{T}^{}}\mathrm{cos}(\widehat{\omega }_1t)\mathrm{cos}(\widehat{\omega }_2t)\sqrt{\widehat{B}_+\widehat{T}^{}}+\widehat{H}_1^{1/4}\mathrm{sin}(\widehat{\omega }_2t)\mathrm{sin}(\widehat{\omega }_1t)\widehat{H}_2^{1/4}\right\}`$ (138)
$`\widehat{F}_{22}(t)`$ $`=`$ $`\gamma \left\{\widehat{C}^{}\mathrm{sin}(\widehat{\omega }_1t)\widehat{T}\widehat{B}_{}\mathrm{cos}(\widehat{\omega }_2t)+\mathrm{cos}(\widehat{\omega }_2t)\widehat{B}_+\widehat{T}^{}\mathrm{sin}(\widehat{\omega }_1t)\widehat{C}\right\}`$ (139)
$`=`$ $`i\gamma \left\{\sqrt{\widehat{B}_+\widehat{T}^{}}\mathrm{cos}(\widehat{\omega }_1t)\mathrm{sin}(\widehat{\omega }_2t)\widehat{H}_1^{1/4}\widehat{H}_1^{1/4}\mathrm{sin}(\widehat{\omega }_2t)\mathrm{cos}(\widehat{\omega }_1t)\sqrt{\widehat{T}\widehat{B}_{}}\right\},`$ (140)
where $`\gamma =4\alpha ^2\beta /\mathrm{}^2`$, and we used the properties
$`\widehat{C}\widehat{C}^{}=\widehat{C}^{}\widehat{C}=1`$ (142)
$`\widehat{C}\mathrm{sin}(\widehat{\omega }_2t)=\mathrm{sin}(\widehat{\omega }_1t)\widehat{C}`$ (143)
$`\widehat{C}^{}\mathrm{cos}(\widehat{\omega }_1t)=\mathrm{cos}(\widehat{\omega }_2t)\widehat{C}^{}`$ (144)
$`\sqrt{\widehat{T}\widehat{B}_{}}\widehat{\omega }_2^n=\widehat{\omega }_1^n\sqrt{\widehat{T}\widehat{B}_{}}`$ (145)
$`\sqrt{\widehat{B}_+\widehat{T}^{}}\widehat{\omega }_1^n=\widehat{\omega }_2^n\sqrt{\widehat{B}_+\widehat{T}^{}},`$ (146)
proved in the appendix A of the Ref. , together with the operators relations
$`\widehat{C}\sqrt{\widehat{B}_+\widehat{T}^{}}=\sqrt{\widehat{T}\widehat{B}_{}}\widehat{C}^{}=i\widehat{H}_2^{1/4}`$ (148)
$`\sqrt{\widehat{B}_+\widehat{T}^{}}\widehat{C}=\widehat{C}^{}\sqrt{\widehat{T}\widehat{B}_{}}=i\widehat{H}_1^{1/4}.`$ (149)
At this point, if we remember that $`[\widehat{\nu }_j,\widehat{\omega }_j]=0`$, $`(j=1,\mathrm{or}\mathrm{\hspace{0.33em}2})`$, then we conclude that we can use the trigonometric relations involving the product of trigonometric function with arguments $`\widehat{\nu }_jt`$ and $`\widehat{\omega }_jt`$ because, in this case, we know that $`\mathrm{exp}(\widehat{\nu }_jt)\mathrm{exp}(\pm \widehat{\omega }_jt)=\mathrm{exp}[(\widehat{\nu }_j\pm \widehat{\omega }_j)t]`$. Now, using this fact, the commutators
$$[\widehat{\nu }_1,\widehat{H}_2]=[\widehat{\omega }_1,\widehat{H}_2]=[\widehat{\nu }_2,\widehat{H}_1]=[\widehat{\omega }_2,\widehat{H}_1]=0,$$
(150)
and the properties (ACKNOWLEDGMENTS), we can show that
$`\widehat{y}_1(t)\widehat{F}_{11}(t)`$ $`=`$ $`i{\displaystyle \frac{\gamma }{2}}\sqrt{\widehat{T}\widehat{B}_{}}\left\{\mathrm{cos}\left[(\widehat{\nu }_2\widehat{\omega }_2)t\right]\mathrm{sin}(\widehat{\omega }_1t)+\mathrm{cos}\left[(\widehat{\nu }_2+\widehat{\omega }_2)t\right]\mathrm{sin}(\widehat{\omega }_1t)\right\}\widehat{H}_2^{1/4}`$ (152)
$`+`$ $`i{\displaystyle \frac{\gamma }{2}}\widehat{H}_2^{1/4}\left\{\mathrm{sin}\left[(\widehat{\nu }_1\widehat{\omega }_1)t\right]\mathrm{cos}(\widehat{\omega }_2t)\mathrm{sin}\left[(\widehat{\nu }_1+\widehat{\omega }_1)t\right]\mathrm{cos}(\widehat{\omega }_2t)\right\}\sqrt{\widehat{B}_+\widehat{T}^{}}`$ (153)
$`\widehat{y}_1(t)\widehat{F}_{12}(t)`$ $`=`$ $`{\displaystyle \frac{\gamma }{2}}\sqrt{\widehat{T}\widehat{B}_{}}\left\{\mathrm{cos}\left[(\widehat{\nu }_2\widehat{\omega }_2)t\right]\mathrm{cos}(\widehat{\omega }_1t)+\mathrm{cos}\left[(\widehat{\nu }_2+\widehat{\omega }_2)t\right]\mathrm{cos}(\widehat{\omega }_1t)\right\}\sqrt{\widehat{T}\widehat{B}_{}}`$ (154)
$`+`$ $`{\displaystyle \frac{\gamma }{2}}\widehat{H}_2^{1/4}\left\{\mathrm{sin}\left[(\widehat{\nu }_1+\widehat{\omega }_1)t\right]\mathrm{sin}(\widehat{\omega }_2t)\mathrm{sin}\left[(\widehat{\nu }_1\widehat{\omega }_1)t\right]\mathrm{sin}(\widehat{\omega }_2t)\right\}\widehat{H}_1^{1/4}`$ (155)
$`\widehat{y}_2(t)\widehat{F}_{21}(t)`$ $`=`$ $`{\displaystyle \frac{\gamma }{2}}\sqrt{\widehat{B}_+\widehat{T}^{}}\left\{\mathrm{cos}\left[(\widehat{\nu }_1\widehat{\omega }_1)t\right]\mathrm{cos}(\widehat{\omega }_2t)+\mathrm{cos}\left[(\widehat{\nu }_1+\widehat{\omega }_1)t\right]\mathrm{cos}(\widehat{\omega }_2t)\right\}\sqrt{\widehat{B}_+\widehat{T}^{}}`$ (156)
$`+`$ $`{\displaystyle \frac{\gamma }{2}}\widehat{H}_1^{1/4}\left\{\mathrm{sin}\left[(\widehat{\nu }_2+\widehat{\omega }_2)t\right]\mathrm{sin}(\widehat{\omega }_1t)\mathrm{sin}\left[(\widehat{\nu }_2\widehat{\omega }_2)t\right]\mathrm{sin}(\widehat{\omega }_1t)\right\}\widehat{H}_2^{1/4}`$ (157)
$`\widehat{y}_2(t)\widehat{F}_{22}(t)`$ $`=`$ $`i{\displaystyle \frac{\gamma }{2}}\sqrt{\widehat{B}_+\widehat{T}^{}}\left\{\mathrm{cos}\left[(\widehat{\nu }_1\widehat{\omega }_1)t\right]\mathrm{sin}(\widehat{\omega }_2t)+\mathrm{cos}\left[(\widehat{\nu }_1+\widehat{\omega }_1)t\right]\mathrm{sin}(\widehat{\omega }_2t)\right\}\widehat{H}_1^{1/4}`$ (158)
$`+`$ $`i{\displaystyle \frac{\gamma }{2}}\widehat{H}_1^{1/4}\left\{\mathrm{sin}\left[(\widehat{\nu }_2\widehat{\omega }_2)t\right]\mathrm{cos}(\widehat{\omega }_1t)\mathrm{sin}\left[(\widehat{\nu }_2+\widehat{\omega }_2)t\right]\mathrm{cos}(\widehat{\omega }_1t)\right\}\sqrt{\widehat{T}\widehat{B}_{}}.`$ (159)
In the same way, we can show that
$`\widehat{z}_1(t)\widehat{F}_{11}(t)`$ $`=`$ $`i{\displaystyle \frac{\gamma }{2}}\sqrt{\widehat{T}\widehat{B}_{}}\left\{\mathrm{sin}\left[(\widehat{\nu }_2\widehat{\omega }_2)t\right]\mathrm{sin}(\widehat{\omega }_1t)+\mathrm{sin}\left[(\widehat{\nu }_2+\widehat{\omega }_2)t\right]\mathrm{sin}(\widehat{\omega }_1t)\right\}\widehat{H}_2^{1/4}`$ (161)
$``$ $`i{\displaystyle \frac{\gamma }{2}}\widehat{H}_2^{1/4}\left\{\mathrm{cos}\left[(\widehat{\nu }_1\widehat{\omega }_1)t\right]\mathrm{cos}(\widehat{\omega }_2t)\mathrm{cos}\left[(\widehat{\nu }_1+\widehat{\omega }_1)t\right]\mathrm{cos}(\widehat{\omega }_2t)\right\}\sqrt{\widehat{B}_+\widehat{T}^{}}`$ (162)
$`\widehat{z}_1(t)\widehat{F}_{12}(t)`$ $`=`$ $`{\displaystyle \frac{\gamma }{2}}\sqrt{\widehat{T}\widehat{B}_{}}\left\{\mathrm{sin}\left[(\widehat{\nu }_2\widehat{\omega }_2)t\right]\mathrm{cos}(\widehat{\omega }_1t)+\mathrm{sin}\left[(\widehat{\nu }_2+\widehat{\omega }_2)t\right]\mathrm{cos}(\widehat{\omega }_1t)\right\}\sqrt{\widehat{T}\widehat{B}_{}}`$ (163)
$``$ $`{\displaystyle \frac{\gamma }{2}}\widehat{H}_2^{1/4}\left\{\mathrm{cos}\left[(\widehat{\nu }_1+\widehat{\omega }_1)t\right]\mathrm{sin}(\widehat{\omega }_2t)\mathrm{cos}\left[(\widehat{\nu }_1\widehat{\omega }_1)t\right]\mathrm{sin}(\widehat{\omega }_2t)\right\}\widehat{H}_1^{1/4}`$ (164)
$`\widehat{z}_2(t)\widehat{F}_{21}(t)`$ $`=`$ $`{\displaystyle \frac{\gamma }{2}}\sqrt{\widehat{B}_+\widehat{T}^{}}\left\{\mathrm{sin}\left[(\widehat{\nu }_1\widehat{\omega }_1)t\right]\mathrm{cos}(\widehat{\omega }_2t)+\mathrm{sin}\left[(\widehat{\nu }_1+\widehat{\omega }_1)t\right]\mathrm{cos}(\widehat{\omega }_2t)\right\}\sqrt{\widehat{B}_+\widehat{T}^{}}`$ (165)
$``$ $`{\displaystyle \frac{\gamma }{2}}\widehat{H}_1^{1/4}\left\{\mathrm{cos}\left[(\widehat{\nu }_2+\widehat{\omega }_2)t\right]\mathrm{sin}(\widehat{\omega }_1t)\mathrm{cos}\left[(\widehat{\nu }_2\widehat{\omega }_2)t\right]\mathrm{sin}(\widehat{\omega }_1t)\right\}\widehat{H}_2^{1/4}`$ (166)
$`\widehat{z}_2(t)\widehat{F}_{22}(t)`$ $`=`$ $`i{\displaystyle \frac{\gamma }{2}}\sqrt{\widehat{B}_+\widehat{T}^{}}\left\{\mathrm{sin}\left[(\widehat{\nu }_1\widehat{\omega }_1)t\right]\mathrm{sin}(\widehat{\omega }_2t)+\mathrm{sin}\left[(\widehat{\nu }_1+\widehat{\omega }_1)t\right]\mathrm{sin}(\widehat{\omega }_2t)\right\}\widehat{H}_1^{1/4}`$ (167)
$``$ $`i{\displaystyle \frac{\gamma }{2}}\widehat{H}_1^{1/4}\left\{\mathrm{cos}\left[(\widehat{\nu }_2\widehat{\omega }_2)t\right]\mathrm{cos}(\widehat{\omega }_1t)\mathrm{cos}\left[(\widehat{\nu }_2+\widehat{\omega }_2)t\right]\mathrm{cos}(\widehat{\omega }_1t)\right\}\sqrt{\widehat{T}\widehat{B}_{}}.`$ (168)
Now, the non-commutativity between the operators $`\widehat{\omega }_1`$ and $`\widehat{\omega }_2`$ imply that to calculate the integrals involving the terms given by the Eqs. (150) and (ACKNOWLEDGMENTS) we need to use the series expansion of the trigonometric functions. In this case the integrals can be easily done because the time variable can be considered as a parameter factor. Finally, using these results into Eq. (90) is trivial to find the expression (III) for the matrix elements of the particular solution.
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# Untitled Document
SPATIOTEMPORAL CHAOS OF SOLITON IN A GENERALIZED SKYRME MODEL WITH THE MODIFIED SYMMETRY - BREAKING TERM .
Nguyen Vien Tho <sup>1</sup><sup>1</sup>1E-mail: ngvtho@dng.vnn.vn
Hue University, Hue, Vietnam
Phu Chi Hoa<sup>2</sup><sup>2</sup>2E-mail: pchihoa@hcm.vnn.vn
Dalat University, Dalat, Vietnam
## Abstract
The chiral symmetry-breaking term of the Skyrme model with massive pion is modified to obtain the hedgehog profile function which is in best coincidence with the kink-like profile function. For the modified Lagrangian, the minimum of the energy of the B=2 twisty skyrmion configuration is lower than the values for both the cases of the Skyrme Lagrangian with and without the non-modified symmetry-breaking term. The equations of motion for the time-dependent hedgehog of this model and for a generalizated Skyrme model including sixth-order stabilizing term are derived and integrated nummerically. The time evolution of soliton is obtained. We have observed the seft-exitation of soliton because of the fast developement of fluctuation.
I. INTRODUCTION
It is known that one can give a description of low energy hadron physics on the base of a semiclassical quatization of the soliton solution of the Skyrme model (skyrmion) $`\left[13\right].`$The skyrmion configuration, denoted by $`U(\stackrel{}{x})`$, is a map: $`R^3SU(2)`$, with the condition $`U1`$ as $`\left|\stackrel{}{x}\right|\mathrm{}`$, required for finite energy . One could analyze in details the case of the spherically symmetry configuration, when the matrix field $`U(\stackrel{}{x})`$ has the form (the hedgehog):
$$U(\stackrel{}{x})=\mathrm{exp}[iF(r)\stackrel{}{x}.\stackrel{}{\tau }]$$
$`(1.1)`$
where $`r=\left|\stackrel{}{x}\right|`$ and $`\stackrel{}{\tau }`$ is the Pauli matrices. Here $`F(r)`$ is a profile function with the boundary conditions $`F(0)=n\pi `$ (n is an integer) and F(r) $`0`$as r$`\mathrm{}`$. F(r) obeys a nonlinear differential equation which could be solve numerically.
An analytic form of the skyrmion profile function seems to be very useful for many purposes. Attemps were made to obtain a such analytic approximation: (i) by computing the holonomy of $`SU(2)`$ Yang- Mills instantons in $`R^4`$ along lines paralell to the time -axis $`\left[4\right]`$; (ii) by identifying the skyrmion profile function with the sin-Gordon kink field (if we replace $`x`$ by $`r`$) $`\left[5\right]`$. The approach (ii) is very attractive because of its simplicity. Moreover, unlike the instanton approach, the sin-Gordon kink has a fixed scale, so there are no abitrary scale parameters which have to be fixed by hand in order to minimize the energy. With the explicit expression of the kink-like profile function $`F(r)`$, in Sec.2, the symmetry-breaking term of the Lagrangian is modified to obtain the hedgehog profile function for B=1 skyrmion which is in best coincidence with the kink-like profile function. Then we have calculated the energy of four different B=2 skyrmion configurations and lead to the result: the lowest value of the energy is obtained in the case of twisty configuration for the modified Lagrangian .
The results on numerical integration of the equation of motion for the time-dependent hedgehog indicate a dynamical chaostic character of fluctuations around the static soliton solution . Such a study is interesting from the viewpoint of considering the Skyrme model as a nonlinear dynamical system. In Sec.3, we find the equation of motion for the time-dependent hedgehog of the considered model. The results of numerical integration of the obtained equation are presented. The fluctuation of the profile function $`\delta F(x,t)`$ as well as the time dependence of the amplitudes of different modes of the fluctuations are plotted.
The Skyrme’s Lagrangian is of fourth order in field derivatives. Various alternative models have been proposed which preserved the form of the original Lagrangian while extending it to higher orders \[8-12\]. The incorporation of the higher order terms, on the one hand, improves the fit of observables, and, on the other hand, gives a reasonable physical interpretation for stabilizing terms in Lagrngian. For example, one could introduce a sixth-order term \[8-9\]
$$^{\left(6\right)}=\frac{\epsilon _6^2}{2}B^\mu B_\mu ,$$
$`(1.2)`$
where $`B^\mu `$ is the baryon current
$$B^\mu =\frac{\epsilon ^{\mu \nu \alpha \beta }}{24\pi ^2}Tr[(U^+_\nu U)(U^+_\alpha U)(U^+_\beta U)].$$
$`(1.3)`$
This term may be understood an representing effects of $`\varpi `$ mesons, while the Skyrme’s fourth-order term should be viewed as representing effects of $`\rho `$ mesons. In Sec. 4, we consider the dynamical behavior of the model which the Skyrme Lagrangian added by the sixth-order term (1.2 ) and the modified symmetry-breaking term. The discussion of the results is given in Conclusion.
II. THE SIN-GORDON KINK FIELD AND THE CHOICE OF THE CHIRAL SYMMETRY BREAKING TERM
The Skyrme’s Lagrangian density takes the form
$$=\frac{F_\pi ^2}{16}Tr\left(_\mu U^\mu U^+\right)+\frac{1}{32e^2}Tr[(_\mu U)U^+,(_\nu U)U^+]^2.$$
$`(2.1)`$
One considers also the Skyrme’s Lagrangian with a pion mass term
$$^{}=+_{SB},$$
$`(2.2)`$
$$_{SB}=\frac{1}{8}m_\pi ^2F_\pi ^2Tr\left(U+U^+2\right).$$
$`(2.3)`$
From (2.1), based on equations of motion, one found the nonlinear differential equation for the B=1 hedgehog profile function
$$\frac{xF^{}}{2}+\left[\frac{x^2}{4}+2Sin^2F\right]F^{^{\prime \prime }}+Sin2F\left(F^{^{}}\right)^2\frac{Sin\left(2F\right)}{4}\frac{Sin^2FSin\left(2F\right)}{x^2}=0,$$
$`(2.4)`$
where $`x=eF_\pi r`$ is the dimensionless radial distance.
The kink-like function has the form
$$F\left(x\right)=4arctan\left(e^x\right),$$
$`(2.5)`$
which satisfies the same boundary conditions. The kink-like profile function (2.5) has an exponential decay for large x . The same of an asymptotic behaviour is also obtained when the Skyrme Lagrangian (2.1) is added by a chiral symmetry-breaking term as (2.3) . Accordingly, the right-hand side of (2.4) is not zero, but equals to
$$\frac{\beta ^2}{4}x^2SinF,$$
$`(2.6)`$
where $`\beta =\frac{m_\pi }{eF_\pi }`$, $`m_\pi `$=140 MeV, e=4.84, $`F_\pi `$ =108 MeV.
However, if one substitutes the kink-like profile function (2.5) in the left-hand side of (2.4) and compare with (2.6) (in Fig.1a and Fig.1b), it is seen that they are different. In order to make the equality to be satisfied approximately, we must modify the pion mass term as following
$$3.5\times 10^7\frac{\beta ^2}{4}x^2SinF.$$
$`(2.7)`$
After modification the pion mass term as (2.7), we have plotted it in Fig.1c. We see that the approximate equality may be acceptable. So, one should be able to choice the symmetry-breaking term as follow
$$_{SB}^{(mod.)}=\frac{\epsilon }{8}m_\pi ^2F_\pi ^2Tr\left(U+U^+2\right),$$
$`(2.8)`$
where $`\epsilon =3.5\times 10^7`$.
Now, we consider an ansazt has the form
$$N=\{Cosk\varphi Sin\vartheta ,Sink\varphi Sin\vartheta ,Cos\vartheta \},$$
$`(2.9)`$
with this ansazt the soliton mass is
$$M=M_2+M_4,$$
$`(2.10)`$
$$M_2=\frac{\gamma }{4}_0^{\mathrm{}}𝑑xx^2_0^\pi 𝑑\vartheta Sin\vartheta \{\left(F^{}\right)^2+[k^2+1]\frac{Sin^2F}{x^2}\},$$
$`(2.11)`$
$$M_4=\gamma _0^{\mathrm{}}𝑑xx^2_0^\pi 𝑑\vartheta Sin\vartheta \{[\frac{Sin^2\vartheta }{Sin^2\vartheta }k^2+1]\left(F^{}\right)^2+\frac{Sin^2F}{x^2}k^2\}\frac{Sin^2F}{x^2},$$
$`(2.12)`$
where $`\gamma =\frac{\pi F_\pi }{e},`$ k is integer- k=1 corresponds to the case of the sherically symmetry hedgehod, k$`2`$ to the case of twisty skyrmion configuration. The variation of (2.10) in F(x) give the following equation
$$[x^2+2a\mathrm{sin}^2F]F^{^{\prime \prime }}+2xF^{^{}}+[a(F^{^{}})^2\frac{a}{4}2b\frac{\mathrm{sin}^2F}{x^2}]\mathrm{sin}(2F)=0,$$
$`(2.13)`$
where
$$a=_0^\pi [k^2+1]\mathrm{sin}\vartheta d\vartheta ,$$
$`(2.14)`$
$$b=k^2_0^\pi \mathrm{sin}\vartheta d\vartheta ,$$
$`(2.15)`$
which for $`k=2`$ we have $`a=10`$ and $`b=8`$.
From above-mentioned results and the formulas (2.11, 2.12), we consider the different cases of B=2 skyrmion configurations .
The first way of obtaining B=2 Skyrme field is to alter the boundary condition on the hedgehog profile function so that $`F(0)=2\pi `$ . One can generate an approximation to this Skyrme field by using the kink-like profile function
$$F\left(x\right)=8arctan\left(e^x\right),$$
$`(2.16)`$
substitute it to (2.11) and (2.12) with k=1, we get
$$M_a=M_2+M_446.7\gamma .$$
$`(2.17)`$
The second way of obtaining B=2 skyrme field is to leave the boundary condition on the profile unchanged but to have the skyrme field rotate twice as rapidly as the radial vector x under a change in the azimuthal angle around an axis . From(2.1), (2.11) and (2.12) with k=2, we get
$$M_b=M_2+M_438\gamma .$$
$`(2.18)`$
In the third case, based on the Lagrangian with $`_{SB}`$ (2.3) and $`k=2`$, we obtained the equation of motion which is the equation (2.13) added by the term (2.6) in the right hand side. Solve this equation, we have nummerical dat file. Substituting it to (2.11) and (2.12) we get
$$M_c=M_2+M_438.4\gamma .$$
$`(2.19)`$
Analogously, in the last case, based on the Lagrangian with $`_{SB}^{(mod.)}`$ and $`k=2`$ we obtained the equation of motion which is the equation (2.13) added by the term (2.7) in the right hand side.We find numerically F(x) (dat file) and substitute it to (2.11) and (2.12) we get
$$M_d=M_2+M_427.5\gamma .$$
$`(2.20)`$
III. THE EQUATION FOR THE TIME -DEPENDENT HEDGEHOG
AND THE TIME EVOLUTION OF THE SOLITON SOLUTION .
We consider the Skyrme Lagrangian (2.1) is added by a chiral symmetry-breaking term (2.8). It is convenient to parametrize the SU(2) matrix field U(x) by the pion field isovector $`\stackrel{}{\pi }\left(x\right)`$
$$U\left(x\right)=\frac{1+i\stackrel{}{\eta }\left(x\right)\stackrel{}{\tau }}{1i\stackrel{}{\eta }\left(x\right)\stackrel{}{\tau }},$$
$`(3.1)`$
where
$$\stackrel{}{\eta }\left(x\right)=\frac{\stackrel{}{\pi }\left(x\right)}{F_\pi },$$
$`(3.2)`$
and $`\stackrel{}{\tau }=(\tau _1,\tau _2,\tau _3)`$ are the Pauli matrices. In the parametrization (3.1) the expressions of the Cartan forms $`(_\mu U)U^+`$ for the SU(2) group have been calculated , and on the base of these expressions we get the expressions of the Lagrangian (2.1) an (2.8) in the form
$$=\frac{F_\pi ^2}{2}\frac{\left(_\mu \stackrel{}{\eta }^\mu \stackrel{}{\eta }\right)}{\left(1+\stackrel{}{\eta }^2\right)^2}\frac{4}{e^2}\frac{\left[_\mu \stackrel{}{\eta }\times ^\mu \stackrel{}{\eta }\right]^2}{\left(1+\stackrel{}{\eta }^2\right)^4},$$
$`(3.3)`$
$$_{SB}^{(mod.)}=\frac{\epsilon }{2}m_\pi ^2F_\pi ^2\frac{\stackrel{}{\eta }^2}{\left(1+\stackrel{}{\eta }^2\right)}.$$
$`(3.4)`$
The time-dependent hedgehog corresponds to the following ansatz
$$\stackrel{}{\eta }(r,t)=tan\left[\frac{F(r,t)}{2}\right]\frac{\stackrel{}{r}}{r},$$
$`(3.5)`$
where F(r,t) is the profile function. Corresponding to this ansatz, $``$ and $`_{SB}^{(mod.)}`$ are given by
$$=\frac{F_\pi ^2}{8}\{(\dot{F})^2(F^{^{}})^2\frac{2\mathrm{sin}^2F}{r^2}\}+\frac{1}{2e^2}\frac{\mathrm{sin}^2F}{r^2}\{2(\dot{F})^22(F^{^{}})^2\frac{\mathrm{sin}^2F}{r^2}\},$$
$`(3.6)`$
$$_{SB}^{(mod.)}=\frac{\epsilon m_\pi ^2F_\pi ^2}{4}\left(\mathrm{cos}F1\right).$$
$`(3.7)`$
The Hamiltonian is given by
$$H=4\pi r^2dr\{\frac{F_\pi ^2}{8}\{(\dot{F})^2+(F^{^{}})^2+\frac{2\mathrm{sin}^2F}{r^2}\}+\frac{1}{2e^2}\frac{\mathrm{sin}^2F}{r^2}\{2(\dot{F})^2+2(F^{^{}})^2+\frac{\mathrm{sin}^2F}{r^2}\}+$$
$$\frac{\epsilon m_\pi ^2F_\pi ^2}{4}(1\mathrm{cos}F)\}.$$
$`(3.8)`$
The variational equation for the profile function is
$$\frac{xF^{}}{2}+\left[\frac{x^2}{4}+2Sin^2F\right](F^{^{\prime \prime }}\ddot{F})+Sin2F(F^{}_{}{}^{}2\dot{F}^2)\frac{Sin\left(2F\right)}{4}\frac{Sin^2FSin\left(2F\right)}{x^2}$$
$$\frac{\epsilon \beta ^2}{4}x^2SinF=0,$$
$`(3.9)`$
where $`\beta =\frac{m_\pi }{eF_\pi }`$, $`m_\pi `$=140 MeV, e=4.84, $`F_\pi `$ =108 MeV \[ 14 \] , and $`x=eF_\pi r`$, $`\tau =eF_\pi t`$ are the dimensionless distance and time. The primes and the dots mean the derivatives with respect to x and $`\tau `$, respectively. Hereafter t is always understood as the dimensionless time.
It is convenient to write F(x,t) in the form
$$F(x,t)=F\left(x\right)+\delta F(x,t),$$
$`(3.10)`$
where F(x) is the profile function of the static hedgehog. To find $`\delta F(x,t)`$ we use the following boundary conditions
$$\delta F(0,t)=\delta F(L,t)=0.$$
$`(3.11)`$
This condition is automatically satisfied by the harmonic expansion
$$\delta F(x,t)=\underset{j=1}{\overset{N1}{}}A_j(t)\mathrm{sin}\left(\frac{j\pi x}{L}\right),$$
$`(3.12)`$
where L is the size of the spatial volume, N is the number of the points of the discretized spatial variable x, $`A_j(t)`$ are the amplitudes of jth fluctuation modes. These amplitudes are obtained by inverting the series given in (3.12)
$$A_j(t)=\frac{2}{L}_0^L\delta F(x,t)\mathrm{sin}(\frac{j\pi x}{L})𝑑x.$$
$`(3.13)`$
In our calculation we choose L=16, N=128 and the initial excitation mode is j =16. We denote $`AA_{16}(0)`$. We have studied the time evolution of the system for the pertubation parameter $`A=0.1`$. The equation (3.9) should be reduced to a system coupled second order differential equations for the time variable t. We solve this system by using the Runge-Kutta procedure and obtain the solutions for $`\delta F(x,t)`$ and $`A_j(t)`$. We have plotted the fluctuation $`\delta F(x,t)`$ at t=0, 100, 200, 300 and 500 in Fig.2a to Fig.2e, respectively. We see that apart from large fluctuations near $`x=0`$, $`\left|\delta F(x,t)\right|0.1`$ . This is understandable as the Skyrmion dynamics is dominated by the small x-region and the deviation from the Skyrmion is small. The dynamical behavior of the system is understood better by observing the time evolution of the amplitudes of various harmonic modes. In Fig.3a to Fig.3c we have plotted $`A_j(t)`$ for j=8, 16, 32 , respectively. It is seen that the amplitude of the initial mode decreases gradually and it has some kind of periodicity in the variation, while amplitudes of other modes increase on the average. The time evolution of the mode of j=8 and j=16 has a periodic behavior when it was consider in a small interval of the time. It is to be expected that after a longer interval of time, all the modes would be of the same magnitude leading to ”thermalization” and ”spatio-temporal chaos”.
IV. DYNAMICAL BEHAVIOR OF SKYRMION IN A GENERALIZED SKYRME MODEL .
We consider the Lagrangian
$$=^{(2)}+^{(4)}+^{(6)}+_{SB}^{(mod.)},$$
$`(4.1)`$
where $`^{(2)}+^{(4)}`$ is the Lagrangian of the Skyrme model (2.1), $`^{(6)}`$ is given by (1.2 ), $`_{SB}^{(mod.)}`$ is given by (2.8). In the parametrization (3.1), the expressions of the Cartan forms $`(_\mu U)U^+`$ for the SU(2) group have been calculated . We get the Lagrangian (1.2) in the form
$$^{(6)}=\frac{4\epsilon _6^2}{3\pi ^4}\frac{1}{\left(1+\stackrel{}{\eta }^2\right)^6}\{(\stackrel{}{\eta })^63(\stackrel{}{\eta })^2(_\alpha \stackrel{}{\eta }^\beta \stackrel{}{\eta })(_\beta \stackrel{}{\eta }^\alpha \stackrel{}{\eta })+$$
$$2(_\nu \stackrel{}{\eta }^\beta \stackrel{}{\eta })(_\beta \stackrel{}{\eta }^\alpha \stackrel{}{\eta })(_\alpha \stackrel{}{\eta }^\nu \stackrel{}{\eta })\}.$$
$`(4.2)`$
By the same caculation in Sec. III, we have obtained the variational equation for the profile function
$$(\frac{x^2}{4}+2Sin^2F+\frac{\gamma }{4}\frac{Sin^4F}{x^2})(\ddot{F}F^\mathrm{"})+(Sin\left(2F\right)+\frac{\gamma }{4}\frac{Sin^2FSin\left(2F\right)}{x^2})(\dot{F}^2F^2)$$
$$(\frac{x}{2}\frac{\gamma }{4}\frac{Sin^4F}{x^3})F^{}+\frac{Sin\left(2F\right)}{4}+\frac{Sin^2FSin\left(2F\right)}{x^2}+\frac{\epsilon \beta ^2}{4}x^2SinF=0,$$
$`(4.3)`$
where $`x=eF_\pi r`$, $`\tau =eF_\pi t`$ are the dimensionless distance and time. The primes and the dots mean the derivatives with respect to x and $`\tau `$, respectively, $`\gamma =\frac{F_\pi ^2\epsilon _6^2e^4}{\pi ^4}`$, $`\beta =\frac{m_\pi }{eF_\pi }`$. We choose $`\epsilon _6^2=5fm^2`$, and t is understood as the dimensionless time.
Based on the expressions (3.10) to (3.13) and the choice of L=16, N=128, $`AA_{16}(0)=0.1`$, we consider the dynamical behavior of the system with Lagragian ( 4.1) by observing the time evolution of the amplitudes of various harmonic modes. We have plotted the profile function of the static hedgehog and the fluctuation $`\delta F(x,t)`$ at t = 0, 100, 200, 300 in Fig. 4a to Fig. 4d. It is a clear indication that $`\delta F(x,t)`$ has violent fluctuations for substantially longer intervalls in x.
The plots of $`A_j(t)`$ for the cases of j= 8, 16, 64, 127 ( Fig. 5a to Fig. 5d ) indicate that the fluctuation amplitudes in the generalized Skyrme model which Lagrangian added by the sixth-order term develop much faster than in the Skyrme model. For this model, the amplitudes of ”spontaneous” fluctuations appear after $`t150`$. One can say that a self-exitation of soliton takes place after $`t150`$.
V. CONCLUSION
We have calculated the energy of four B=2 skyrmion configurations. The comparison of the values of the energy in four cases shows that in the case of modified Lagrangian one could obtain B=2 configuration with lowest energy. That is, this configuration is more close to the real energy minimum of B=2 skyrmion.
The B=1 hedgehog kink-like profile function (2.5) is a convenient analytic approximation to the numerical solution. But how can modify the Lagrangian to obtain the profile function (2.5), the modification made in this paper is one of the answers to the question.
By integrating the equation (3.1) for the time-dependent hedgehog we have obtained the information about the dynamical behavior of the sotiton in the model (3.3, 3.4). We have plotted $`\delta F(x,t)`$ at various moments and the development of amplitudes of fluctuation modes $`A_j(t)`$. The plots of $`A_j(t)`$ for the case A=0.1 indicate that the process of thermalization takes place sooner, and the fluctuation amplitudes in the considered model develop much faster than in the Skyrme model (see ). Besides, in , for A=0.1 the amplitude of the initial (32nd) mode has extremely regular periodic behavior while in our initial (16nd) mode the amplitudes decreases gradually on the average.
When higher order terms are included, the theory becomes much more nonlinear, so it is clearly that the phase space around the soliton solution in the generalized Skyrme model was shown to be more stochastic than in the original Skyrme model.
We are grateful to To Ba Ha for valuable help. This work was supported in part by the National Basic Research Program in Natural Sciences under the grant number KT-04-1.2/99.
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# Infrared observations of serendipitous hard Chandra X-ray sources
## 1 Introduction
The X-ray Background (XRB) above 2 keV has at last been mostly resolved into point sources by the Chandra X-ray Observatory (Mushotzky et al 2000; Brandt et al 2000). Chandra’s superb sub-arcsecond imaging provided the high-sensitivity confusion-less images for this breakthrough.
The generally accepted model for the XRB is that it is dominated by absorbed active galactic nuclei (AGN: Setti & Woltjer 1989; Madau, Ghisellini & Fabian 1994; Comastri et al 1995), a significant fraction of which are Compton-thick (i.e. have an absorbing column $`>10^{24}\mathrm{cm}^2`$): collectively such objects are Type II AGN. So far, only the 2-7 keV XRB has been (mostly) resolved with Chandra. This regime is not expected to be sensitive to large numbers of Compton-thick objects, where the X-ray emission emerges only above 5 keV (Wilman & Fabian 1999; Wilman, Fabian & Nulsen 2000); it should reveal sources with column densities more typically of $`10^{2223}`$$`\mathrm{cm}^2`$. The Compton-thick objects provide much of the power where the XRB spectrum peaks in $`\nu I_\nu `$. Correction of the XRB for absorption (Fabian & Iwasawa 1999) shows that accretion at the standard 10 per cent efficiency onto massive black holes can approximately account for the local mass density in black holes (Magorrian et al 1998) and that about 85 per cent of that accretion power is absorbed and re-radiated in the mid- to far-IR.
Despite the recent progress resolving the hard XRB into discrete sources with the correct collective spectrum, the actual identification of many of these objects is, however, not straightforward. Roughly one-third are blue broad-line quasars, another third are identified with faint, optically-normal galaxies and the final third have only extremely faint optical counterparts or no detectable counterpart at all (Mushotzky et al 2000; Brandt et al 2000; Maiolino et al 2000). Mushotzky et al (2000) find infrared HK band counterparts for most of their X-ray detected sources.
We have used SCUBA maps of the cores of the lensing clusters A2390 and A1835 to place deep submillimetre limits on the three serendipitous Chandra sources which lie in the field (Fabian et al 2000). Only one (marginal) source is detected in both the X-ray and submillimetre bands. Of three X-ray sources in the HST field of A2390 we find that one plausibly has a photometric redshift of 0.9 and conclude from its hard X-ray spectrum that it is a Type II quasar; the other two have $`V>26`$ (Fabian et al 2000). Such objects are therefore difficult to follow up in the optical band alone. Although two SCUBA sources in the A370 field that optically resemble AGN are detected in X-rays (Bautz et al, in preparation), none of the 10 SCUBA sources in fields flanking the Hubble Deep Field is detected in a deep Chandra observation (Hornschemeier et al 2000).
In this paper we present observations which are part of a programme to determine the origin of the optically-faint, hard X-ray sources that are most likely responsible for the hard X-ray background. Given the possibility that the sources are highly redshifted and/or obscured, we have sought infrared counterparts at the X-ray target position. Our work bridges other surveys in this field which are either very deep, or shallow with a wider area coverage. By targeting the serendipitous sources from several 10-20 ksec Chandra observations, we are able to select the very brightest absorbed sources in each field for follow-up.
## 2 Selection and properties of X-ray targets
The X-ray sources were selected from those found serendipitously in the field of Chandra observations of galaxy clusters obtained during the Guaranteed Time of one of us (ACF). These observations were typically of 10-20 ksec duration (see Table 1) and all were taken so that the on-axis pointing position fell on the ACIS-S3 detector, apart from that with sequence number 800010. We shall present a detailed analysis of the X-ray properties of the entire sample of serendipitous sources elsewhere; here we detail follow-up observations of a preliminary small sample of sources of a very specific type.
We examined each of the cluster fields for serendipitous sources using the Chandra Interactive Analysis of Observations (CIAO) detection software. All three available detect algorithms (celldetect, wavdetect and vtpdetect) were used in order to search for sources under the different assumptions and methods that each employ. In the first instance, we ran the detect software on unbinned data in three energy bands: 0.5-7 keV, ‘soft’ 0.5-2 keV and ‘hard’ 2-7 keV. Those sources with detectable hard counts were then selected for more detailed follow-up. As detailed later (section 4.5), the higher sensitivity of Chandra below 2 keV means that even genuinely hard sources can still show plenty of counts in our soft band. Most of the sources in this paper were detected close to the pointing position, either in the S3 or S2 chips. The exceptions were the two sources in the field of A2199 (CXOU J162850.9+392434 and CXOU J162827.8+392343), which were both in I3, and CXOU J031946.4+413734 and CXOU J031946.2+413737 in the Perseus field which were in the I2 chip. The sources were 2-9 arcminutes off-axis, with a reduced effective area down to 80 per cent of the on-axis value.
We calculated the number of counts in each of these three energy bands using a box typically of side 16 pixels (8 arcsec; although the box size was increased for sources well off-axis, ie for CXOU J162827.8+392343, CXOU J091340.9+410314 and CXOU J140048.4+024954) from which to extract the source counts. The background countrate was estimated from a box around the source (excluding the source box) of sidelength typically 55 pixels (27.5 arcsec; or larger for those objects with countrate extracted from a larger source box). The only exception was in the case of CXOU J031946.4+413734 and CXOU J031946.2+413737 in the Perseus field, where the sources were so close that the background rate for both was estimated from a close region of sky. The counts in all three energy bands, as well as the ratio of soft/hard counts are given in Table 1.
As the Chandra observations were relatively early on in the operation of the satellite, all were affected by inaccuracies in the aspect solution during the initial pipeline processing of the data. This led to offsets of up to 8 arcseconds between the coordinate positions given by Chandra and those of the sky. This offset was easy to correct for in the observations of IRAS 09104+4109 (800017) and Perseus (800010), where the position of the AGN can be determined with half-arcsecond precision from the Chandra hard-band image, to be compared to sub-arcsecond radio core position. A hard nuclear peak was not, however, found in the Chandra images of A1795, A1835 or A2199, so the X-ray optical registration was done by a systematic cross-correlation between all the X-ray sources against optical sources from the Digitized Sky Survey (DSS). Through this process we improved the accuracy of the attitude solution to within about 2 arcseconds. Note we assume that only a simple translation in sky coordinates is required, with no rotation or stretching. The (registered) X-ray source positions were then compared to optical sources on the DSS. Although we have been using multi-colour (archival) optical imaging data of our fields, many of the sources lie a few arcmin out from the centre of the cluster, and are not always covered by previous observations.
This inaccurate aspect solution to the X-ray data and the uncertainties inherent in the correction of the Chandra coordinates to the sky presents some confusion when applied to conventional naming of the source. The coordinates used to name the source are derived from a definite ID of the source from IR imaging presented later in this paper. Where no source is detected in the IR imaging we use the IR-detected sources within this field to apply the correct offset between the Chandra frame of reference and the sky, and thus derive the source name.
We selected our targets from those sources that had a soft-to-hard (S/H) ratio of less than 3.5 (see Table 1; also section 4.5 for implications of the S/H ratio), and either a very faint, or no, optical identification on the DSS. The exceptions were two intriguing sources in the ACIS-I observation of the Perseus cluster. Their unusually close proximity to each other (5 arcsec) and exact association with two close optical sources marked them out as particularly interesting and worthy of detailed follow-up. The DSS B- and R-band images around each source observed in this paper are shown in Fig 1, along with matching Chandra images in the 0.5-7 keV, 0.5-2 keV and 2-7 keV energy bands. Faint optical identifications are visible not only for CXOU J031946.4+413734 and CXOU J031946.2+413737, but also CXOU J091357.5+405938, CXOU J162827.8+392343 and (perhaps) CXOU J091360.0+405548. The faint optical source seen in the CXOU J134905.8+263752 box is several arcseconds away from the X-ray position, and thus we assume it is not associated.
## 3 Observations
### 3.1 Optical spectra
Optical spectra of the close X-ray sources CXOU J031946.4+413734 and CXOU J031946.2+413737 were obtained in service time with the ISIS double beam spectrograph on the William Herschel Telescope (WHT) on La Palma, during the night of 1999 Dec 15. The 1 arcsec-wide slit was oriented at a position angle of 138 in order to detect the emission from both sources. The total exposure was 3000 s, and the R158R and R158B gratings were used to produce a wavelength range of 3500-5500Å on the EEV chip in the blue arm, and 5150-8050Å on the Tek chip in the red arm of the instrument. The data were bias-subtracted, flat-fielded, wavelength-calibrated from exposures of an arclamp, and corrected for the Galactic extinction of E(B-V)=0.31 in this direction.
### 3.2 Infrared observations
Near-infrared spectra and images were taken during the nights of 2000 Feb 24-25 at the United Kingdom Infrared Telescope (UKIRT) in Hawaii. The night of Feb 24 was of very good seeing and transparency, both of which, however, deteriorated by the time of our observations on Feb 25. We used two instruments: the 2D grating spectrometer CGS4, and the cooled infrared camera IRCAM3/TUFTI. A full log of observations is shown in Table 2.
CGS4 was used with the 40 l/mm grating on the long camera, yielding a pixel scale of 0.618 arcsec per pixel. The slit width was set to 4 pixels (ie 2.47 arcsec), and observations were taken using the standard ‘quad-slide’ nodding pattern of a-b-b-a along the slit. Spectra were taken in the first order in each of the H and K bands on some objects, as well as of corresponding spectrophotometric and atmospheric absorption standards in the same band. The data were reduced using the standard reduction package CGS4DR V1.3-0. The object spectrum was divided by the spectrum of a standard star observed at approximately the air mass, assumed to approximate to a black body, in order to remove any common atmospheric features. The object spectra were flux-calibrated using a standard star; the flux calibration was estimated to be accurate to within ten per cent. The final, calibrated spectrum of the object was extracted from three rows centred on the object peakup row, and three rows around the sky peakup row 19 pixels away.
IRCAM3/TUFTI is an imaging camera with a scale of 0.0814 arcsec/pixel, and a total field of view of 20.8$`\times `$20.8 arcsec. TUFTI was used in ND STARE mode, with standard read-out. Objects were observed using a jitter grid of 9 points, each separated by 6 arcsec. We observed for 60 s at each of the grid points, except for CXOU J140100.2+025720, leading to total exposure time of 540 s in each of the J,H and K bands. The imaging observations for CXOU J140100.2+025720 were curtailed because we were approaching the end of that night’s observing. The data were reduced using the standard ORACDR reduction package available from UKIRT, and flux-calibrated using observations of several standard stars in the same bands.
## 4 Results and Discussion
### 4.1 X-ray luminosities
We derive the bolometric luminosity of the X-ray sources from the observed 0.5-7 keV count rate, assuming the emission originates in a non-thermal power-law with photon index of $`\mathrm{\Gamma }=2`$, subject only to Galactic absorption. The observed 0.5-7 keV fluxes range over $`4.363.3\times 10^{15}`$$`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$(Table 1), implying a bolometric luminosity range of $`1.724.8\times 10^{44}`$ and $`9.0130.8\times 10^{45}`$$`\mathrm{erg}\mathrm{s}^1`$for redshifts of 0.5 and 3 respectively.
### 4.2 Optical spectra
The optical spectra of CXOU J031946.4+413734 and CXOU J031946.2+413737 are shown in Fig 2. CXOU J031946.4+413734 (the optically-fainter, X-ray-harder source to the south-east) has a solitary broad line observed at 6454Å, with a FWHM of 5460$`\mathrm{km}\mathrm{s}^1`$and an intensity of 1.6$`\times 10^{15}`$$`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$. The most likely identification for this is MgII, implying a redshift for the source of $`z=1.307`$; however we do not obviously see CII\]$`\lambda `$2336 or CIII\]$`\lambda `$1909 emission line at 5365Å and 4404Å, or H$`\alpha `$ at 1.514$`\mu `$m in the infrared spectrum (Fig 3; the 2-$`\sigma `$ emission line nearby is at 1.504$`\mu `$m). The line can only be CIII\] if the expected Ly$`\alpha `$ at 4110Å is completely extinguished by dust absorption; it is unlikely to be Ly$`\alpha `$ itself, as the continuum does not cut off due to a Ly$`\alpha `$-forest as expected for a source at $`z=4.3`$. CXOU J031946.2+413737 shows a relatively featureless continuum, with too little flux for a reliable redshift to be obtained from cross-correlation with a template galaxy spectrum.
### 4.3 Infrared spectra
We detected continuum emission from four of the objects observed with CGS4, only failing to detect any signal at the position of CXOU J091360.0+405548 at an upper limit of $`5.1\times 10^{19}`$$`\mathrm{erg}\mathrm{s}^1\mathrm{cm}^2`$Å$`^1`$. (CXOU J091360.0+405548 is the source we also failed to detect from the infrared imaging, see next section). All the detected objects have flat infrared spectra with no significant emission-line features (Fig 3). Where magnitudes from the infrared imaging are available (see next section), they are compatible with the spectra (see Fig 3 for a direct comparison). The flux inferred from the infrared imaging tends to be higher than that determined from the spectra, which is to be expected if our slit placement does not quite cover all the light from the object.
We searched for emission lines in those parts of the infrared spectrum least affected by any OH atmospheric absorption (if the sky is variable then we may not be able to correct for these features fully in the data reduction and spurious emission lines would be introduced; these regions of the spectrum are shown by the dotted lines in Fig 3). We find no significant emission lines, with 3-$`\sigma `$ upper limits to the equivalent width of 185-720Å in H, and 145-190Å in K (depending on the length of the exposure). We would expect some Balmer or Paschen line emission to be visible for targets in most of the redshift range $`0<z<4`$, and MgII to be visible for sources $`4.3<z<5.4`$. It is thus surprising that none of our targets show strong emission lines in the CGS4 spectra. This suggests that any active nucleus – and any ionized gas that surrounds it – is heavily obscured.
### 4.4 Infrared imaging
Of the nine sources imaged in good seeing conditions, we detect infrared sources in the 21$`\times `$21 arcsec TUFTI field of view for eight of them. A further three objects attempted during conditions of poor seeing and transparency were not detected. Due to the variability of the conditions we are unable to quantify the detection limits. CXOU J091360.0+405548 (our least significant X-ray source) was not detected to a 3-$`\sigma `$ limit of 22.66 in J, 22.27 in H and 21.62 in K (assuming a point source). The J, H and K-band TUFTI images are shown in Fig 4 for all sources where infrared sources were found in the field. In Table 3 we list the infrared magnitudes (and estimated DSS B and R magnitudes or limits) of each source near the centre of the TUFTI image, along with its offset from the boresight of TUFTI (and thus from the X-ray position). The position of this aim-point and its constancy was established from the observations of standard stars throughout the night.
There is a very clear single identification of a source near the centre of the TUFTI field of view for seven of the eight sources shown in Fig 4. The source for CXOU J140106.9+024934 is resolved into three separate components in the K band. CXOU J091352.8+405829 is the only source (out of those closest to the centre) to show any significant ellipticity. The registration of the X-ray coordinates to the sky (as indicated by the offset of the source from the aimpoint of TUFTI) is, not surprisingly, best in the field of IRAS 09104+4109, the source with a clear detection of the point-like AGN at hard energies to cross-correlate with the radio position. Offsets are larger for the other fields, but all are good to within 2.1 arcsec. Only CXOU J140100.2+025720 has no clear identification, as there are two sources that are both at a larger offset from the centre than expected: it is not clear that either of these are the identification of the X-ray source. For clarity we will, however, refer to them as CXOU J140100.2+025720 (SSE) and (ESE), according to the direction of the offset from the position of the X-ray source CXOU J140100.2+025720. There is no correlation between the J-K colours of our detected objects and the ratio of their soft to hard X-ray flux, but this is not unexpected as we have deliberately chosen a very small range of X-ray flux ratio to investigate in the first place.
### 4.5 Absorbed AGN?
We consider it unlikely that our detected objects are optically-faint stars, as any star with a sufficient level of X-ray activity is unlikely to show such a featureless infrared continuum. The infrared colours of our detected targets are very different to those of main sequence, giant and supergiant stars of all spectral types (Fig 5), although the infrared colours alone are consistent with some IR-selected stars (eg the samples shown in Fig 5 of Dickinson et al 2000). We note also that stars comprised only 6 per cent of the sources found in the ROSAT Deep Survey (Schmidt et al 1998). The infrared colours of our sources are better matched to those of X-ray-selected quasars in the compilation of Elvis et al (1994).
We performed a simple X-ray spectral analysis of the two brightest sources, CXOU J134905.8+263752 and CXOU J134849.0+263716 using XSPEC. The spectra were fitted with power-law models of fixed photon index $`\mathrm{\Gamma }=2`$, absorbed (in our frame) by a column density N<sub>H</sub>. The X-ray spectra require excess column densities of N$`{}_{\mathrm{H}}{}^{}=2.6\pm 0.7`$ and $`3.5_{1.1}^{+1.3}\times 10^{21}`$$`\mathrm{cm}^2`$, for CXOU J134905.8+263752 and CXOU J134849.0+263716 (errors are 1-$`\sigma `$). For example, the fit to CXOU J134905.8+263752 for a power-law model (of fixed $`\mathrm{\Gamma }`$=2) with only the Galactic absorption gives $`\chi ^2`$=24.5 (for 12 degrees of freedom), improving to $`\chi ^2`$=6.1 (for 11 d.o.f.) when the absorption is allowed to vary freely, giving N$`{}_{\mathrm{H}}{}^{}=2.6\pm 0.7\times 10^{21}`$$`\mathrm{cm}^2`$(Fig 6). A fit with the absorption fixed at Galactic, and now with the power-law slope varying freely yields $`\mathrm{\Gamma }=1.4\pm 0.14`$, with $`\chi ^2`$=9.6 (11 d.o.f.), and a completely free fit gives $`\mathrm{\Gamma }=1.9\pm 0.2`$ and N$`{}_{\mathrm{H}}{}^{}=2.3\pm 0.9\times 10^{21}`$$`\mathrm{cm}^2`$, with $`\chi ^2`$=6.06 (for 10 d.o.f.). Thus there is a greater than 95 per cent probability that the source requires excess absorption over the Galactic column. The models with $`\mathrm{\Gamma }=2`$ and free N<sub>H</sub> predict de-absorbed fluxes at 1keV of $`2.0\pm 0.3\times 10^5`$ and $`8.8_{1.7}^{+1.8}\times 10^6`$ photons$`\mathrm{cm}^2\mathrm{keV}^1\mathrm{s}^1`$for CXOU J134905.8+263752 and CXOU J134849.0+263716 respectively (equivalent to 13 and 6 nJy).
We converted these de-absorbed fluxes at 1 keV to intrinsic K-band magnitudes, assuming a spectral index of $`\alpha _{\mathrm{KX}}=1.16`$ (such that $`S\nu ^\alpha `$). This index is the median for the sample of 41 quasars in Elvis et al (1994), who chose optically-bright objects with high signal-to-noise ratio Einstein X-ray spectra. They are predominantly low redshift objects, so K-corrections should be negligible. The predicted K-band magnitude of 17.3 for CXOU J134905.8+263752 is slightly fainter than the observed value of 17.01$`\pm `$0.06, whereas the predicted magnitude of 18.2 for CXOU J134849.0+263716 is a lot brighter than the observed value of 19.99$`\pm `$0.25. This suggests that for the total K-band light of CXOU J134849.0+263716 to be due wholly to the AGN continuum, it must be reddened by $`2`$ magnitudes at K. For a Galactic dust-to-gas ratio, this equates to N$`{}_{\mathrm{H}}{}^{}=3.5\times 10^{22}\mathrm{cm}^2`$, substantially more than the value derived from the X-ray fitting.
We extend this calculation to predict optical and infrared magnitudes for all our TUFTI-detected sources, assuming that all the emission is due to a typical quasar at a range of redshifts. We use PIMMS to obtain the normalization of the quasar spectral energy distribution (SED) at 1 keV from the observed 0.5-7 keV count rate, assuming $`\mathrm{\Gamma }=2`$ using the Galactic column. We then extrapolate from the inferred X-ray flux at 2 keV using the (rest-frame) SED of a typical radio-quiet quasar, approximated from that given in Elvis et al (1994); $`\alpha _{OX}=1.4`$ (bridging the monochromatic luminosities at 2500Å and 2keV), $`\alpha =0.15`$ (over 0.25-1.28$`\mu `$m), $`\alpha =2.38`$ (1.28-1.85$`\mu `$m) and $`\alpha =1.38`$ ($`>`$1.85$`\mu `$m).
For all the sources, the observed SEDs are redder in slope than that expected from the quasar predictions (Fig 7). The majority of the objects have unabsorbed quasar predictions that are insufficient (by up to 1.5 magnitudes) to account for the observed magnitudes, suggesting that the AGN continuum is not the main contributor to the infrared magnitudes. Where the (low-redshift) quasar SEDs are more compatible with the observed infrared magnitudes (CXOU J091340.9+410314, CXOU J134905.8+263752) the quasar SED still overestimates the optical band limits. There remains only one source (CXOU J134849.0+263716) where the extrapolated spectrum vastly overestimates the observed magnitudes; here any AGN component must be reddened by at least 1.5 magnitudes in the infrared. These results are, however, sensitive to the $`\alpha _{OX}`$ used: varying $`\alpha _{OX}`$ by only $`\pm 0.1`$ leads to a corresponding $`\pm 0.6`$ mag variation in the optical-infrared magnitudes shown in Fig 7. Even if $`\alpha _{OX}`$ is adjusted so that the quasar can account for all of the K-band light, its continuum must be highly reddened to also fit the observed B and R limits or magnitudes (we assume that the lack of detection in the optical is not due to variability). The results indicate both that a host galaxy contributes much of the infrared light of many of our sources, and that any contribution from the quasar continuum requires significant reddening.
Further constraints on the redshift and N<sub>H</sub> may be obtained by comparing the soft to hard (S/H) count ratios (Table 1) with those predicted by XSPEC models, as shown in Table 4 and Fig 8. For a given model (ie $`z`$ and N<sub>H</sub>), the differences between the predictions for the two brightest sources reflect the different response matrices of the front- (CXOU J134905.8+263752 on S2) and back- (CXOU J134849.0+263716 on S3) illuminated CCDs. The observed S/H ratios are $`2.4\pm 0.4`$ (for CXOU J134905.8+263752) and $`2.5\pm 0.5`$ (for CXOU J134849.0+263716). Very roughly, the values in Table 4 suggest that if our two sources had power-law spectra, they have N$`{}_{\mathrm{H}}{}^{}10^{22}`$$`\mathrm{cm}^2`$if at low redshift, or more like N$`{}_{\mathrm{H}}{}^{}10^{23}`$$`\mathrm{cm}^2`$if at $`z1`$, As a first approximation, the Tables may also be used to interpret the S/H count ratios for the other sources, for which there are insufficient counts to justify spectral fitting \[using the CXOU J134905.8+263752 (and CXOU J134849.0+263716) predicted values as guides for sources on the front- (back-)illuminated chips\]. We have attempted simple spectral fitting of the sources with around 50 counts, and although the error bars are larger, the results confirm the conclusions of Table 4. The X-ray colours of all our sources are consistent with the X-ray emission originating in an AGN at a range of redshift, but only if that emission is absorbed by an intrinsic column density of at least N$`{}_{\mathrm{H}}{}^{}10^{22}`$$`\mathrm{cm}^2`$at $`z>0.5`$. We conclude that the bulk of our objects have X-ray column densities which classify them as Compton-thin Type II objects.
### 4.6 Photometric redshifts
Given that at least one of the detected infrared objects is resolved to have significant ellipticity, and that the host galaxies of AGN usually make the dominant contribution to the near-infrared light of AGN (Rix et al 1999), we also compare the infrared colours and magnitudes to those of galaxies at a range of redshift. The lack of strong emission lines in the infrared spectra argues against strong starburst behaviour, so we consider only elliptical galaxies. We used the programme HYPERZ (Bolzonella, Pello & Miralles 2000), which fits a grid of template galaxy spectra (generated from the GISSEL library; Bruzual & Charlot 1993) to the observed magnitudes, with variations permitted in the age, redshift and intrinsic reddening. We attempted a variety of fits: a 3 Gyr-old elliptical galaxy (with e-folding time $`\tau `$ of 1 Gyr) fitted without and with the possibility of intrinsic reddening; and a choice between an elliptical galaxy at ages 1.5, 3 and 5 Gyr again without and then with reddening. In all cases we assumed the reddening law of Calzetti et al (2000) and constrained the age so that it did not exceed the age of the Universe at the redshift under consideration (assuming $`H_0=50`$$`\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$). The maximum redshift considered was $`z=6`$.
Our results for the basic model (a 3 Gyr-old elliptical with no intrinsic reddening) are shown for three of our detected sources in Table 5. We also plot the best-fit solutions for this model to CXOU J134905.8+263752, CXOU J134849.0+263716 and CXOU J140106.9+024934 in Fig 9. Increasing the range of possible elliptical galaxy ages, and/or allowing the possibility of intrinsic reddening to this model decreased the $`\chi ^2`$, but broadened the range of possible redshifts, and sometimes introduced a secondary $`\chi ^2`$ minimum at lower redshift. Without deeper optical limits it is hard to constrain the redshift in any detail; we can only say that the infrared colours are not inconsistent with an origin from elliptical galaxies above a redshift of one.
Host galaxies to radio-quiet quasars generally have a modest luminosity of just less than, or around that of an L galaxy, with radio-loud quasars lying in galaxies a factor of 2-5 times brighter (Rix et al 1999). We thus also compare our observed infrared magnitudes to those expected from a passively-evolving, unabsorbed L elliptical galaxy (R. McMahon & A. Aragon-Salamanca, private communication) in Fig 10. This crude comparison of magnitudes suggests that CXOU J140106.9+024934, CXOU J134849.0+263716 and CXOU J140100.2+025720 (ESE) are consistent with L galaxies at a redshift of $`z=2`$, with other sources such as CXOU J140100.2+025720 (SSE), CXOU J091340.9+410314, CXOU J091352.8+405829 and CXOU J162827.8+392343 being nearer $`z=1`$. The inferred redshifts may be underestimated, however, if the infrared colours and magnitudes also include a contribution from a central quasar continuum.
## 5 Summary and conclusions
We have carried out follow-up observations of optically-faint, X-ray hard, serendipitous Chandra sources, and find that they are readily detected in the near-infrared. Spectra in the infrared of some of them appear flat and featureless, suggesting that strong emission-line activity is either absent or heavily obscured. Only one source – which is the brightest optically – shows a strong emission line in an optical spectrum, which we cannot identify unambiguously. The 0.5-7 keV fluxes of our sources are in the range $`0.34.6\times 10^{14}`$$`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$which, if due to an unabsorbed non-thermal quasar spectrum, imply bolometric luminosities of a few times $`10^{44}10^{46}`$$`\mathrm{erg}\mathrm{s}^1`$, for redshifts between $`0.1<z<3`$. The X-ray colours of all our sources are consistent with an origin in quasars only if absorbed by an intrinsic column density of at least $`10^{22}`$$`\mathrm{cm}^2`$. If we extrapolate a typical radio-quiet quasar spectrum at $`0<z<3`$ from the observed X-ray flux into the optical and infrared wavebands, we find that a host galaxy probably contributes much of the infrared light, and that the optical quasar continuum requires significant reddening. The infrared magnitudes and colours are consistent with a (host) L$``$Lgalaxy at moderate redshifts $`0.5<z<2.5`$.
Although we do not have deep optical images for most of our objects, the limits imply that the optical and infrared properties may resemble so-called extremely red objects (EROs: see Scodeggio & Silva 2000 and references therein). These have proven difficult to follow up even with large telescopes, and good spectra and redshifts are only available for a small fraction. Whether there is a deeper, physical connection between optically faint Chandra sources and EROs must await larger samples of Chandra sources with deep optical and infrared coverage.
The properties of our small sample of objects is consistent with the existence of a population of moderately absorbed (i.e. Compton thin Type II) quasars at $`z=12`$, as predicted by recent models for the hard XRB (Madau et al 1994; Celotti et al 1995; Comastri et al 1995; Wilman & Fabian 1999).
Finally, we note that the objects we have likely found, namely Type II quasars, are qualitatively different from Type II Seyferts. They do not obviously have any strong narrow-line region (although deeper spectra covering a wider band on more objects are needed to be definite on this issue), nor any obvious scattered blue continuum or blue light from star formation. They would not even be classified as active galaxies on the basis of what we have seen so far of their optical and infrared properties. It is only on the basis of the X-ray emission that we classify them as Type II objects, where Type II means strongly obscured (see also the discussion in Matt et al 2000). An appropriate name is X-ray Type II quasar. Of course, if 85 per cent of accretion power is absorbed (Fabian & Iwasawa 1999) then the objects dominating that power cannot have the broad torus opening-angles commonly ascribed to Seyfert galaxies. The simple and successful geometrical unification scheme for Seyfert galaxies (e.g. Antonucci 1993) cannot extend to quasars.
## 6 Acknowledgements
We are grateful to the Chandra project for the X-ray data, and to M. Bolzonella, R. Pello and J.-M. Miralles for making HYPERZ available. CSC and ACF thank the Royal Society, PG thanks the Isaac Newton Trust and the Overseas Research Trust, and RJW and RMJ thank the PPARC for financial support. We thank the service queue scheme on the William Herschel Telescope, which is operated on the island of La Palma by the Isaac Newton Group in the Spanish Observatorio del Roque de los Muchachos of the Instituto de Astrofisica de Canarias. The United Kingdom Infrared Telescope is operated by the Joint Astronomy Centre on behalf of the U.K. Particle Physics and Astronomy Research Council. This research has made use of the NASA/IPAC Extragalactic Database (NED), and the Digitized Sky Surveys which were produced at the Space Telescope Science Institute under U.S. Government grant NAG W-2166.
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# Uniform Bounds and Symbolic Powers on Smooth Varieties
## Introduction
The purpose of this note is to show how one can use multiplier ideals to establish effective uniform bounds on the multiplicative behavior of certain families of ideal sheaves on a smooth algebraic variety. In particular, we prove a quick but rather surprising result concerning the symbolic powers of radical ideals on such a variety.
Let $`X`$ be a non-singular quasi-projective variety defined over the complex numbers $`𝐂`$, and let $`ZX`$ be a reduced subscheme of $`X`$.<sup>1</sup><sup>1</sup>1All of our results are local in nature, so there is no loss in taking $`X`$ to be an affine variety. In this case one can work with the coordinate ring $`𝐂[X]`$ of $`X`$ in place of its structure sheaf $`𝒪_X`$. Denote by
$$𝔮=_Z𝒪_X$$
the ideal sheaf of $`Z`$, so that $`𝔮`$ is a sheaf of radical ideals. We shall be concerned with the symbolic powers $`𝔮^{(m)}`$ of $`𝔮`$. According to a well-known theorem of Zariski and Nagata (see , Chapter 3.9) $`𝔮^{(m)}`$ can be described as the sheaf of all function germs vanishing to order $`m`$ at a general point of each irreducible component of $`Z`$ (or equivalently at every point of $`Z`$):
$$𝔮^{(m)}=\{f𝒪_X\text{ord}_x(f)m\text{for all }xZ\}.$$
It is evident that $`𝔮^m𝔮^{(m)}`$, but in general of course the inclusion is strict. However Swanson established (in a much less restrictive setting<sup>2</sup><sup>2</sup>2Swanson’s theorem holds in particular on any normal variety over a field of any characteristic.) that there exists an integer $`k=k(Z)`$ depending on $`Z`$ such that
$$𝔮^{(km)}𝔮^m\text{ for all }m𝐍.$$
On geometric grounds this already seems rather striking since membership in the symbolic power on the left is tested at general smooth points of $`Z`$, whereas the actual power on the right reflects also its singular points. So one’s first guess might be that the worse the singularities of $`Z`$, the larger one will have to take the coefficient $`k(Z)`$ to be. Surprisingly enough this is not the case, and in fact one has a uniform statement depending only on the codimension of $`Z`$:
###### Theorem A.
Assume that every component of $`Z`$ has codimension $`e`$ in $`X`$. Then
$$𝔮^{(me)}𝔮^m\text{ for all }m𝐍.$$
In particular, if $`dimX=n`$ then $`𝔮^{(mn)}𝔮^m`$ for every radical ideal $`𝔮𝒪_X`$ and every natural number $`m1`$. One can see the Theorem as providing further confirmation of Huneke’s philosophy that there are unexpected uniform bounds lurking in commutative algebra.
Theorem A is a very simple application of the theory of multiplier ideals. In commutative algebra these were introduced by Lipman in connection with the Briançon-Skoda theorem.<sup>3</sup><sup>3</sup>3Lipman called them “adjoint ideals”, but “multiplier ideal” has become standard in higher dimensional geometry. The name derives from their analytic construction, where they arise as sheaves of multipliers (see ). More general constructions, which we use here, have in the meantime become extremely important in the study of higher dimensional algebraic varieties (cf. , , , , , ). In brief, we consider families $`𝔞_{}=\{𝔞_k\}`$ of ideals $`𝔞_k𝒪_X`$ — such as the symbolic powers $`𝔮^{()}=\{𝔮^{(k)}\}`$ — satisfying the relations
$$𝔞_{\mathrm{}}𝔞_m𝔞_{\mathrm{}+m}\text{ for all }\mathrm{},m1.$$
For each $`\mathrm{}1`$ we associate to such a family an asymptotic multiplier ideal $`𝒥(𝔞_{\mathrm{}})𝒪_X`$ which reflects the asymptotic behavior of all the ideals $`𝔞_p\mathrm{}`$ for $`p0`$. Using the subadditivity theorem of , we prove
###### Theorem B.
If $`𝒥(𝔞_{\mathrm{}})𝔟`$ for some index $`\mathrm{}`$ and some ideal $`𝔟`$, then
$$𝔞_m\mathrm{}𝔟^m$$
for every $`m1`$.
In the case of symbolic powers it is elementary to check that $`𝒥(𝔮^{(e)})𝔮`$, so Theorem A follows from the “abstract” Theorem B. As another application, we establish a result (Theorem 2.6) rendering effective and extending in certain directions a theorem of Izumi , dealing with ideals arising from a valuation.
We have been guided by the viewpoint that the families $`\{𝔞_k\}`$ share some of the behavior of the linear series $`|kD|`$ associated to multiples of a divisor $`D`$ on a projective variety, and that one can try to adapt geometric tools to the present setting. We hope that these and other ideas from higher dimensional complex geometry will find further algebraic applications in the future. Going in the other direction, Hochster and Huneke have used the theory of tight closure to reprove and generalize the results of the present paper dealing with symbolic powers: in paticular they show that Theorem A holds for any regular local ring containing a field, and they remove the hypothesis that $`𝔮`$ be radical (see §3 for further discussion). This illustrates once again the close but somewhat mysterious connections between tight closure methods and the more geometric outlook appearing here.
Our exposition is organized into two sections. In the first, we construct the multiplier ideals we use and establish their basic properties. The applications are given in §2.
We are grateful to Mel Hochster and Craig Huneke for valuable discussions and encouragement, and to Jessica Sidman for some Macaulay scripts related to Example 2.3. We also wish to record our debt to the work of Irena Swanson and her collaborators, through which we learned of many of the questions discussed here.
## 1. Graded Families and Multiplier Ideals
In this section we construct the multiplier ideals we require, and give their basic properties. Quick overviews of the general theory of multiplier ideals appear in and , §1, and a survey of some of the applications in algebraic geometry is given in . The forthcoming book will contain a detailed exposition, which in the meantime can be found in the lecture notes . In particular, contains full proofs of all the facts about multiplier ideals quoted in the following paragraphs.
Let $`X`$ be a non-singular complex quasi-projective variety, and $`𝔞𝒪_X`$ a non-zero ideal sheaf on $`X`$. A log resolution of $`𝔞`$ is a projective birational map $`\mu :X^{}X`$, with $`X^{}`$ non-singular, such that $`𝔞𝒪_X^{}=𝒪_X^{}(F)`$ for an effective Cartier divisor $`F`$ on $`X`$ with the property that the sum of $`F`$ and the exceptional divisor of $`\mu `$ has simple normal crossing support. Such resolutions can be construced by resolving the singularities of the blow-up of $`𝔞`$. We write $`K_{X^{}/X}=K_X^{}\mu ^{}K_X`$ for the relative canonical divisor of $`X^{}`$ over $`X`$.
Given a rational number $`c>0`$, the multiplier ideal associated to $`c`$ and $`𝔞`$ is defined by fixing a log resolution as above, and setting
$$𝒥(X,c𝔞)=𝒥(c𝔞)=\mu _{}𝒪_X^{}\left(K_{X^{}/X}[cF]\right).$$
Here $`cF`$ is viewed as an effective $`𝐐`$-divisor, and its integer part $`[cF]`$ is defined by taking the integral part of the coefficient of each of its components. The fact that $`𝒥(c𝔞)`$ is indeed an ideal follows from the observation that $`𝒥(c𝔞)\mu _{}𝒪_X^{}(K_{X^{}/X})=𝒪_X`$. An important point is that this definition is independent of the log resolution $`\mu `$.
It follows immediately from the definition that if $`c𝐍`$, then $`𝒥(c𝔞)=𝒥(𝔞^c)`$. This being so, we sometimes prefer to use “exponential notation” $`𝒥(𝔞^c)`$ for the multiplier ideal $`𝒥(c𝔞)`$ for an arbitrary rational number $`c>0`$. Note that we are not trying to attach any actual meaning to the expression $`c𝔞`$ or $`𝔞^c`$ when $`c`$ is non-integral. Nonetheless, the possibility of being able to work with rational coefficients is critical in applications.
As a variant, given ideals $`𝔞,𝔟𝒪_X`$, and rational numbers $`c,d>0`$, we define $`𝒥((c𝔞)(d𝔟))`$ (or $`𝒥(𝔞^c𝔟^d)`$ in exponential notation) by taking a common log resolution $`\mu :X^{}X`$ of $`𝔞`$ and $`𝔟`$, with $`𝔞𝒪_X^{}=𝒪_X^{}(F_1),𝔟𝒪_X^{}=𝒪_X^{}(F_2)`$, and putting $`𝒥(𝔞^c𝔟^d)=\mu _{}𝒪_X^{}\left(K_{X^{}/X}[cF_1+dF_2]\right)`$. It is sometimes useful also to adopt the convention that if $`𝔞=(0)`$, then $`𝒥(c𝔞)=(0)`$ for all $`c>0`$.
The most important local property of multiplier ideals is the Restriction Theorem, due in the algebro-geometric setting to Esnault and Viehweg. Specifically, let $`YX`$ be a smooth subvariety, and let $`𝔞𝒪_X`$ be an ideal sheaf whose zeroes do not contain $`Y`$. Then $`𝔞_Y=𝔞𝒪_Y`$ is an ideal sheaf on $`Y`$, and the result in question states that one has an inclusion:
(1)
$$𝒥(Y,c𝔞_Y)𝒥(X,c𝔞)𝒪_Y$$
of ideal sheaves on $`Y`$. This is established by reducing to the case where $`Y`$ has codimension one, and applying vanishing theorems. The restriction theorem leads in turn to the Subaddivity Theorem of , which states (in exponential notation) that given ideals $`𝔞,𝔟𝒪_X`$ and rational numbers $`c,d>0`$, one has the inclusion:
(2)
$$𝒥(𝔞^c𝔟^d)𝒥(𝔞^c)𝒥(𝔟^d).$$
To prove this, one first of all applies the Künneth formula to check that
$$𝒥(X\times X,(p_1^1𝔞)^c(p_2^1𝔟)^d)=p_1^1𝒥(X,𝔞^c)p_2^1𝒥(X,𝔟^d),$$
where $`p_1,p_2:X\times XX`$ are the projections, where we are somewhat abusively writing $`f^1𝔬𝒪_V`$ for inverse image $`𝔬𝒪_V`$ of an ideal $`𝔬𝒪_W`$ under a morphism $`f:VW`$. Then one restricts to the diagonal. Note for later reference that in “additive notation” (2) implies
(3)
$$𝒥(cm𝔞)𝒥(c𝔞)^m$$
for every integer $`m1`$.
Some of the most interesting applications of multiplier ideals (for instance , ) depend on the fact that one can make asymptotic constructions. A natural algebraic setting for these is described in the following
###### Definition 1.1.
A graded family or graded system of ideals $`𝔞_{}=\{𝔞_k\}`$ is a collection of ideal sheaves $`𝔞_k𝒪_X`$ $`(k1)`$ satisfying
(4)
$$𝔞_k𝔞_{\mathrm{}}𝔞_{k+\mathrm{}}\text{ for all }k,\mathrm{}1.$$
To avoid unnecessary complications, we assume also that $`𝔞_k(0)`$ for $`k0`$.
Note that if we set $`A_0=𝒪_X`$, then condition (4) is equivalent to the statement that $`_{k0}𝔞_k`$ is a graded $`𝒪_X`$-algebra. The asymptotic constructions that follow are particularly useful in case this algebra is not finitely generated (or at least not known to be so).
###### Example 1.2.
* Let $`(0)𝔞𝒪_X`$ be a fixed ideal, and take $`𝔞_k=𝔞^k`$ to be the $`k^{\text{th}}`$ power of $`𝔞`$. Then the $`\{𝔞_k\}`$ form a graded family. One should view this as a trivial example.
* Let $`D`$ be a divisor on a projective variety $`X`$. When $`H^0(X,𝒪_X(kD))0`$ let $`𝔟_k=𝔟\left(|kD|\right)`$ be the base-ideal of the complete linear series $`|kD|`$, and put $`𝔟_k=(0)`$ otherwise. Then $`𝔟_{}=\{𝔟_k\}`$ forms a graded family of ideals.
* Let $`(0)𝔮𝒪_X`$ be a radical ideal. Then the symbolic powers $`\{𝔮^{(k)}\}`$ form a graded family of ideals that we denote by $`𝔮^{()}`$.
* Let $`\nu :YX`$ be a proper birational map, and let $`D`$ be a non-zero effective Cartier divisor on $`Y`$. Then we get a graded family of ideals $`𝔬_{}=\{𝔬_k\}`$ on $`X`$ by putting $`𝔬_k=\nu _{}𝒪_Y(kD)`$. Note that this includes the symbolic powers $`𝔮^{(k)}`$ in (iii) as a special case, as well as the graded family of ideals associated to an $`𝔪`$-valuation on $`X`$ in the sense of .
* Let $`p(t)=_{i=1}^{\mathrm{}}\frac{1}{i!}t^i𝐂[[t]]`$ be the power series of the function $`e^t1`$, and given $`f𝐂[x,y]`$ define
$$v(f)=\text{ord}_tf(t,p(t)).$$
This is a valuation on $`𝐂[x,y]`$, and therefore the ideals
$$𝔬_k=_{\text{def}}\{f𝐂[x,y]v(f)k\}$$
(which we may view as ideal sheaves on $`X=𝐂^2`$) form a graded family. Explicitly,
$$𝔬_k=(x^k,yp_k(x)),$$
where $`p_k(t)=_{i=1}^k\frac{1}{i!}t^i`$ is the $`k^{\text{th}}`$ Taylor polynomial of $`e^t1`$.
* Assume that $`X`$ is affine (and as always non-singular), so that ideal sheaves are identified with ideals in the coordinate ring $`𝐂[X]`$ of $`X`$. Given any non-zero ideal $`𝔞𝐂[X]`$, set
$$𝔞^{\{k\}}=\{f𝐂[X]Df𝔞\text{ differential operators }D\text{ on }X\text{ of order }<k\text{ }\}.$$
This determines a graded family $`𝔞^{\{\}}`$ which also reduces to the symbolic powers $`\{𝔮^{(k)}\}`$ when $`𝔞=𝔮`$ is radical.
* Let $`𝔞_{}=\{𝔞_k\}`$ be a graded family, and $`𝔟𝒪_X`$ a fixed ideal. Then the colon ideals
$$𝔯_k=(𝔞_k:𝔟^k)=_{\text{def}}\{f𝒪_Xf𝔟^k𝔞_k\}$$
form a graded family.
We now construct the asymptotic multiplier ideal associated to a graded family $`𝔞_{}`$.
###### Lemma 1.3.
Let $`𝔞_{}=\{𝔞_k\}`$ be a graded family of ideals, and fix $`\mathrm{}𝐍`$ plus a rational number $`c>0`$. Then for all positive integers $`p,n1`$ one has
$$𝒥(\frac{c}{p}𝔞_p\mathrm{})𝒥(\frac{c}{pn}𝔞_{pn\mathrm{}}).$$
###### Proof.
Let $`\mu :X^{}X`$ be a common log resolution of $`𝔞_p\mathrm{}`$ and $`𝔞_{pn\mathrm{}}`$, with
$$𝔞_p\mathrm{}𝒪_X^{}=𝒪_X^{}(F_p\mathrm{}),𝔞_{pn\mathrm{}}𝒪_X^{}=𝒪_X^{}(F_{pn\mathrm{}}).$$
Condition (4) implies that $`𝔞_p\mathrm{}^n𝔞_{pn\mathrm{}}`$, and hence $`nF_p\mathrm{}F_{pn\mathrm{}}`$ (i.e. the difference $`nF_p\mathrm{}F_{pn\mathrm{}}`$ is effective). Therefore
$$K_{X^{}/X}[\frac{c}{p\mathrm{}}F_p\mathrm{}]K_{X^{}/X}[\frac{c}{pn\mathrm{}}F_{pn\mathrm{}}],$$
and the statement follows. ∎
We assert next that the collection of multiplier ideals
(5)
$$\left\{𝒥(\frac{c}{p}𝔞_p\mathrm{})\right\}_{(p>0)}$$
has a unique maximal element. In fact, the existence of one maximal element follows from the ascending chain condition on ideals. On the other hand, if $`𝒥(\frac{c}{p}𝔞_p\mathrm{})`$ and $`𝒥(\frac{c}{q}𝔞_q\mathrm{})`$ are both maximal, then thanks to the Lemma they each coincide with $`𝒥(\frac{c}{pq}𝔞_{pq\mathrm{}})`$.
###### Definition 1.4.
Given a graded family of ideals $`𝔞_{}=\{𝔞_k\}`$, the asymptotic multiplier ideal at level $`\mathrm{}`$ associated to $`c>0`$ and $`𝔞_{}`$, written $`𝒥(c𝔞_{\mathrm{}})`$, is the maxmial element of the collection of ideals appearing in (5). In other words,
(6)
$$𝒥(c𝔞_{\mathrm{}})=𝒥(\frac{c}{p}𝔞_p\mathrm{})\text{ for }\text{sufficiently divisible }p0.\mathit{}$$
Assuming as we are that $`𝔞_k(0)`$ for $`k0`$, one can show that there is an integer $`p_0=p_0(𝔞_{},\mathrm{})`$ such that $`𝒥(c𝔞_{\mathrm{}})=𝒥(\frac{c}{p}𝔞_p\mathrm{})`$ for all $`pp_0`$. We use this fact only to observe that one does not actually need the divisibility condition in (6).
###### Remark 1.5.
Note that $`𝒥(c𝔞_{\mathrm{}})`$ depends not just on the particular ideal $`𝔞_{\mathrm{}}`$, but on all the ideals $`𝔞_p\mathrm{}`$ for $`p0`$. The double vertical lines should serve as a reminder of this point.
###### Example 1.6.
* If $`𝔞_k=𝔞^k`$ is the trivial graded family consisting of powers of a fixed ideal $`𝔞`$, then $`𝒥(c𝔞_{\mathrm{}})=𝒥(c𝔞^{\mathrm{}})=𝒥(c\mathrm{}𝔞)`$.
* When $`𝔟_k=𝔟\left(|kD|\right)`$ is the family of base ideals associated to a big divisor $`D`$, then $`𝒥(c𝔟_{\mathrm{}})=𝒥(c\mathrm{}D)`$ is the asymptotic multiplier ideal constructed for instance in and . These ideals have played an important role in recent work on linear series.
* Let $`𝔮𝒪_X`$ be a radical ideal. We denote the asymptotic multiplier ideal at level $`\mathrm{}`$ associated to the symbolic powers $`𝔮^{()}=\{𝔮^{(k)}\}`$ by $`𝒥(c𝔮^{(\mathrm{})})`$. Thus $`𝒥(c𝔮^{(\mathrm{})})=𝒥(\frac{c}{p}𝔮^{(p\mathrm{})})`$ for $`p0`$.
* Consider the ideals $`𝔬_k𝐂[x,y]`$ constructed in Example 1.2 (v) associated to the valuation $`v(f)=\text{ord}_tf(t,e^t1)`$. Then $`𝒥(𝔬_{\mathrm{}})=𝐂[x,y]`$ for every $`\mathrm{}`$. This can be checked directly using the observation that each $`𝔬_k`$ contains a polynomial whose divisor is a smooth curve. From a more sophisticated point of view, the triviality of the multiplier ideal in question is implied by Theorem B plus the fact that the colength of $`𝔬_k`$ in $`𝐂[x,y]`$ grows linearly rather than quadratically in $`k`$.
For our purposes the essential properties of these multiplier ideals are given by
###### Proposition 1.7.
Let $`𝔞_{}=\{𝔞_k\}`$ be a graded family of ideals on the smooth variety $`X`$, and fix $`\mathrm{}1`$. Then:
* $`𝔞_{\mathrm{}}𝒥(𝔞_{\mathrm{}})`$.
* For every $`m1`$ one has the inclusion
$$𝒥(𝔞_m\mathrm{})𝒥(𝔞_{\mathrm{}})^m.$$
###### Proof.
Since the relative canonical bundle $`K_{X^{}/X}`$ is effective, it follows from the definition that $`𝔞𝒥(𝔞)`$ for any ideal $`𝔞𝒪_X`$. Then using Lemma 1.3 we find that
$$𝔞_{\mathrm{}}𝒥(𝔞_{\mathrm{}})𝒥(\frac{1}{p}𝔞_p\mathrm{}).$$
Taking $`p0`$, this gives (i). For (ii), fix $`p0`$ and use the subadditivity relation (3) to deduce:
$`𝒥(𝔞_m\mathrm{})`$ $`=𝒥(\frac{1}{p}𝔞_{pm\mathrm{}})`$
$`=𝒥(\frac{m}{pm}𝔞_{pm\mathrm{}})`$
$`𝒥(\frac{1}{pm}𝔞_{pm\mathrm{}})^m`$
$`=𝒥(𝔞_{\mathrm{}})^m,`$
as asserted. ∎
###### Remark 1.8.
Note that it need not be true in general that $`𝒥(𝔞_m\mathrm{})𝒥(𝔞_{\mathrm{}})^m`$. This explains why it is crucial to pass to the asymptotic ideals.
## 2. Applications
Our concrete results follow from the following general statement — which appears as Theorem B in the Introduction — concerning the multiplicative behavior of graded families of ideals:
###### Theorem 2.1.
Let $`𝔞_{}=\{𝔞_k\}`$ be a graded family of ideals on a smooth complex variety $`X`$, and suppose that $`𝔟𝒪_X`$ is an ideal such that $`𝒥(𝔞_{\mathrm{}})𝔟`$ for some index $`\mathrm{}𝐍`$. Then
$$𝔞_m\mathrm{}𝔟^m$$
for every integer $`m1`$.
Proof. This is an immediate consequence of Proposition 1.7, which implies that
$$𝔞_m\mathrm{}𝒥(𝔞_m\mathrm{})𝒥(𝔞_{\mathrm{}})^m.\mathit{}$$
The first application is to symbolic powers:<sup>4</sup><sup>4</sup>4See §3 for a more general statement.
###### Theorem 2.2.
Let $`X`$ be a smooth complex variety, and $`ZX`$ a reduced subscheme all of whose irreducible components have codimension $`e`$ in $`X`$. Put $`𝔮=_Z`$, and fix an integer $`\mathrm{}e`$. Then
$$𝔮^{(m\mathrm{})}\left(𝔮^{(\mathrm{}+1e)}\right)^m$$
for every $`m1`$. In particular, taking $`\mathrm{}=e`$ one has
$$𝔮^{(me)}𝔮^m\text{ for all }m1.$$
###### Proof.
It suffices by Theorem 2.1 to show that
(\*)
$$𝒥(𝔮^{(\mathrm{})})𝔮^{(\mathrm{}+1e)}.$$
But membership in the ideal on the right is tested locally at a general point of each irreducible component of $`Z`$. So we can assume after shrinking $`X`$ that $`Z`$ is smooth and irreducible, of codimension $`e`$, and in this case (\*) is clear. For then $`𝔮^{(k)}=𝔮^k`$ for all $`k`$, and $`𝔮`$ is resolved by taking $`\mu :X^{}=\text{Bl}_Z(X)X`$ to be the blow-up of $`X`$ along $`𝔮`$. Writing $`EX^{}`$ for the corresponding exceptional divisor, one has
$$K_{X^{}/X}=(e1)E\text{and}𝔮^{\mathrm{}}𝒪_X^{}=𝒪_X^{}(\mathrm{}E).$$
Consequently
$$𝒥(𝔮^{(\mathrm{})})=\mu _{}𝒪_X^{}\left(K_{X^{}/X}\mathrm{}E\right)=\mu _{}𝒪_X^{}\left((\mathrm{}+1e)E\right)=𝔮^{l+1e},$$
as asserted. ∎
###### Example 2.3.
The first non-trivial case of Theorem 2.2 is the following. Let $`T𝐏^2`$ be a finite set of points, viewed as a reduced algebraic subset of the plane, and let $`I𝐂[X,Y,Z]`$ be the homogeneous ideal of $`T`$. If $`F𝐂[X,Y,Z]`$ is a homogeneous polynomial having multiplicity $`2m`$ at every point of $`T`$, then $`FI^m`$. (Apply Theorem 2.2 to the affine cone over $`T`$ in $`𝐂^3`$.) In spite of the very classical nature of this statement we do not know a direct elementary proof.
###### Remark 2.4.
The statement of Theorem 2.2 can fail on singular varieties. For example Huneke points out that counter-examples arise already when $`Z`$ is a line on a quadric cone $`X`$ in $`𝐂^3`$. However Hochster and Huneke give some statements valid also on singular ambient spaces.
###### Remark 2.5.
Using familiar arguments, one can deduce from Theorem 2.2 that the corresponding statement holds for excellent regular local rings containing a field of characteristic zero. However Hochster and Huneke have shown that in fact the analogue of (2.2) holds in any regular local ring containing a field. Therefore we do not dwell on the question finding the most general situation in which the arguments of the present paper apply.
We conclude with a result which renders effective and extends in certain directions a formulation due to Hübl and Swanson (, (1.4)) of a theorem of Izumi :
###### Theorem 2.6.
Let $`\nu :YX`$ be a proper birational map between smooth complex varieties. Let $`EY`$ be a prime divisor, set
$$\mathrm{}=1+\text{ord}_E(K_{Y/X})$$
and for $`k1`$ put $`𝔬_k=\nu _{}𝒪_Y(kE)`$. Fix an irreducible subvariety $`ZX`$ such that $`Z\nu (E)`$ and denote by $`𝔭=_Z`$ the ideal of $`Z`$. Then
$$𝔬_\mathrm{}m𝔭^m\text{ for all }m1.$$
###### Remark 2.7.
The result discussed in — which holds in considerably more general settings, but without the explicit determination of the coefficient $`\mathrm{}`$ of $`m`$ — deals with the situation in which $`E`$ maps to a point. It was in trying to understand this result that we were led to the statements about symbolic powers.
###### Proof of Theorem 2.6.
We can assume without loss of generality that $`E`$ is $`\nu `$-exceptional and that $`Z=\nu (E)`$, so that $`\nu _{}𝒪_Y(E)=𝔭`$. Applying (2.1) to the graded family $`𝔬_{}=\{𝔬_k\}`$ (Example 1.2(iv)), it suffices to prove that $`𝒥(𝔬_{\mathrm{}})𝔭`$. We suppose to this end that we’ve fixed a large integer $`p0`$ such that the multiplier ideal $`𝒥(𝔬_{\mathrm{}})=𝒥(\frac{1}{p}𝔬_p\mathrm{})`$ in question is computed on a log resolution $`\mu :X^{}X`$ of $`𝔬_p\mathrm{}`$ dominating $`\nu :YX`$. Then $`E`$ gives rise to a prime divisor on $`E^{}`$ on $`X^{}`$ — viz. the proper transform of $`E`$ — with
$$\text{ord}_E^{}(K_{X^{}/X})=\text{ord}_E(K_{Y/X})=\mathrm{}1,$$
and one has $`𝔬_k=\mu _{}𝒪_X^{}(kE^{})`$ for every $`k1`$.
Let $`F`$ be the effective Cartier divisor on $`X^{}`$ defined in the usual way by writing $`𝔬_p\mathrm{}𝒪_X^{}=𝒪_X^{}\left(F\right)`$. Since $`𝔬_p\mathrm{}=\mu _{}𝒪_X^{}\left(p\mathrm{}E^{}\right)`$, we see that $`E^{}`$ appears with coefficient $`p\mathrm{}`$ in $`F`$. Consequently
$$\text{ord}_E^{}\left(K_{X^{}/X}\left[\frac{1}{p}F\right]\right)(\mathrm{}1)\mathrm{}=1,$$
and therefore
$$𝒥(𝔬_{\mathrm{}})=\mu _{}𝒪_X^{}\left(K_{X^{}/X}[\frac{1}{p}F]\right)\mu _{}𝒪_X^{}\left(E^{}\right)=𝔭,$$
as required. ∎
## 3. Generalizations
In the preprint , which appeared shortly after the first version of the present paper, Hochster and Huneke use the theory of tight closure to extend Theorem 2.2 to arbitrary regular local rings containing a field. They also observe that it is sufficient to assume that $`𝔮𝒪_X`$ is unmixed. In this section we indicate how one applies Theorem 2.1 to treat unmixed ideals.
We start by recalling the definition of symbolic powers in this more general setting. Assume for simplicity of exposition that the smooth complex variety $`X`$ is affine. Given an ideal $`𝔮𝐂[X]`$, fix a primary decomposition
(\*)
$$𝔮=𝔮_1\mathrm{}𝔮_h$$
of $`𝔮`$, and let $`Y_i=\text{Zeroes}\left(\sqrt{𝔮_i}\right)`$ be the subvarieties of $`X`$ corresponding to the associated primes $`𝔭_i=\sqrt{𝔮_i}`$ of $`𝔮`$. Recall that $`𝔮`$ is unmixed if none of the associated primes $`𝔭_i`$ are embedded (or equivalently if there are no inclusions among the $`Y_i`$). Then the symbolic powers $`𝔮^{(k)}𝐂[X]`$ of $`𝔮`$ are defined as follows. For each associated subvariety $`Y_i`$ of $`𝔮`$, there is a natural map $`\varphi _i:𝐂[X]𝒪_{Y_i}X`$ from the coordinate ring of $`X`$ to the local ring of $`X`$ along $`Y_i`$. We then set
$$𝔮^{(k)}=\underset{i=1}{\overset{h}{}}\varphi _i^1\left(𝔮^k𝒪_{Y_i}X\right).$$
In other words, $`f𝔮^{(k)}`$ if and only if there is an element $`s𝐂[X]`$, not lying in any of the associated primes $`𝔭_i`$ of $`𝔮`$, such that $`fs𝔮^k`$.
Theorem 2.2 then admits the following
###### Variant.
Let $`𝔮𝐂[X]`$ be an unmixed ideal, and assume that every associated subvariety $`Y_i`$ of $`𝔮`$ has codimension $`e`$ in $`X`$. Then $`𝔮^{(me)}𝔮^m`$ for all natural numbers $`m1`$.
###### Sketch of Proof.
The symbolic powers $`𝔮^{()}=\{𝔮^{(k)}\}`$ again form a graded family of ideals, so Theorem 2.1 will apply as soon as we establish that $`𝒥(𝔮^{(e)})𝔮`$. Referring to the primary decomposition (\*), it is enough to show that
(\**)
$$𝒥(𝔮^{(e)})𝔮_i\text{for each}1ih.$$
For a given index $`i`$, inclusion in $`𝔮_i`$ is tested at a generic point of $`Y_i`$. So having fixed $`i`$ we are free to replace $`X`$ by any open subset meeting $`Y_i`$. Therefore, by definition of the symbolic powers, we may assume after localizing that $`𝔮^{(k)}=𝔮^k`$. But in this case $`𝒥(𝔮^{(e)})=𝒥(𝔮^e)`$, and $`𝒥(𝔮^e)𝔮𝔮_i`$ thanks to a variant of a theorem of Skoda (cf. ). ∎
###### Remark 3.1.
While there are certain similarities of spirit between the arguments appearing here and those of Hochster and Huneke — e.g. both involve asymptotic constructions, and reduce to the situation in which $`𝔮^{(k)}=𝔮^k`$ — the precise connections between the two points of view remain quite mysterious. In the hopes of understanding these connections more clearly, it is interesting to observe that the properties of multiplier ideals used here can be “axiomatized” as follows. Given a graded family $`𝔮_{}=\{𝔮_k\}`$ what is required for the application to symbolic powers is the existence of ideals $`(𝔮_m)𝒪_X`$ satisfying the following properties:
1. $`(𝔮_m)`$ is a sheaf on $`X`$, i.e. it commutes with localization, and when $`𝔮_k=𝔞^k`$ is the trivial family consisting of powers of a fixed ideal $`𝔞`$, then Skoda’s theorem
$$(𝔞^n)𝔞$$
holds;<sup>5</sup><sup>5</sup>5One also could ask for more precise statements involving the codimensions of associated primes of $`𝔞`$.
2. $`𝔮_m(𝔮_m)\text{ for all }m`$;
3. One has the subadditivity relation:
$$(𝔮_\mathrm{}m)(𝔮_m)^{\mathrm{}}.$$
In our setting, the required ideals are of course given by the asymptotic multiplier ideals $`𝒥(𝔮_m)`$. However the existence of such ideals $``$ is a purely algebraic question, and it would be very interesting to give a construction e.g. using ideas from tight closure. The hope here is that such a construction might serve as a Rosetta stone to help in deciphering the connections between the methods of the present note and the theory of tight closure. ∎
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# 1 High resolution (𝜆/Δ𝜆={37,0000}) spectrum of the 𝑧_𝑒=3.50 quasar Q1159+123, taken with the Keck High Resolution Spectrograph (exposure time 8 h). The data are taken from Songaila (1998). The Ly𝛼forest is clearly seen in absorption blueward of the atomic hydrogen Ly𝛼emission line from the quasar (the broad peak at 5470Å), and is produced by resonant Ly𝛼scattering in gas clouds along the line-of-sight between us and the quasar. The spectrum shows a Lyman-limit system just shortward of 4150Å, i.e. close to the quasar emission redshift. Most of the features between the Ly𝛼and C IVemission (the other broad peak just below 7000Å) are C IVintergalactic absorption lines.
THE INTERGALACTIC MEDIUM<sup>1</sup><sup>1</sup>1A review for the Encyclopedia of Astronomy and Astrophysics (Institute of Physics Publishing).
About half a million years after the Big Bang, the ever-fading cosmic blackbody radiation cooled below 3000 K and shifted first into the infrared and then into the radio, and the smooth baryonic plasma that filled the Universe became neutral. The Universe then entered a “dark age” which persisted until the first cosmic structures collapsed into gravitationally-bound systems, and evolved into stars, galaxies, and black holes that lit up the Universe again. Some time between redshift of 7 and 15, stars within protogalaxies created the first heavy elements; these systems, together perhaps with an early population of quasars, generated the ultraviolet radiation that reheated and reionized the cosmos. The history of the Universe during and soon after these crucial formative stages is recorded in the all-pervading intergalactic medium (IGM), which is believed to contain most of the ordinary baryonic material left over from the Big Bang. Throughout the epoch of structure formation, the IGM becomes clumpy and acquires peculiar motions under the influence of gravity, and acts as a source for the gas that gets accreted, cools, and forms stars within galaxies, and as a sink for the metal enriched material, energy, and radiation which they eject. Observations of absorption lines in quasar spectra at redshifts up to 5 provide invaluable insight into the chemical composition of the IGM and primordial density fluctuation spectrum of some of the earliest formed cosmological structures, as well as of the ultraviolet background radiation that ionizes them.
COSMOLOGICAL REIONIZATION
At epochs corresponding to $`z1000`$, the IGM is expected to recombine and remain neutral until sources of radiation develop that are capable of reionizing it. The detection of transmitted flux shortward of the Ly$`\alpha `$wavelength in the spectra of sources at $`z5`$ implies that the hydrogen component of this IGM was ionized at even higher redshifts. There is some evidence that the double reionization of helium may have occurred later, but this is still controversial. It appears then that substantial sources of ultraviolet photons were already present when the Universe was less than 7% of its current age, perhaps quasars and/or young star-forming galaxies: an episode of pre-galactic star formation may provide a possible explanation for the widespread existence of heavy elements (like carbon, oxygen, and silicon) in the IGM, while the integrated radiation emitted from quasars is likely responsible for the reionization of intergalactic helium. Establishing the epoch of reionization and reheating is crucial for determining its impact on several key cosmological issues, from the role reionization plays in allowing protogalactic objects to cool and make stars, to determining the small-scale structure in the temperature fluctuations of the cosmic background radiation. Conversely, probing the reionization epoch may provide a means for constraining competing models for the formation of cosmic structures, and of detecting the onset of the first generation of stars, galaxies, and black holes in the Universe.
INTERGALACTIC HYDROGEN DENSITY
The proper mean density of hydrogen nuclei at redshift $`z`$ may be expressed in standard cosmological terms as:
$`\overline{n}_\mathrm{H}`$ $`=`$ $`(\rho _{\mathrm{crit}}/m_\mathrm{H})(1Y)\mathrm{\Omega }_b(1+z)^3`$ (1)
$`=`$ $`(1.1\times 10^5\mathrm{cm}^3)(1Y)\mathrm{\Omega }_bh^2(1+z)^3,`$ (2)
where $`Y`$ is the primordial He abundance by mass, $`\rho _{\mathrm{crit}}=3H_0^2/(8\pi G)`$ is the critical density, $`\mathrm{\Omega }_b=\rho _b/\rho _{\mathrm{crit}}`$ is the current baryonic density parameter, and $`H_0=100h\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ is the present-day Hubble constant. Standard nucleosynthesis models together with recent observations of Deuterium yield $`Y=0.247\pm 0.02`$, and $`\mathrm{\Omega }_bh^2=0.0193\pm 0.0014`$. Thus,
$$\overline{n}_\mathrm{H}=(1.6\times 10^7\mathrm{cm}^3)\left(\frac{\mathrm{\Omega }_bh^2}{0.019}\right)(1+z)^3.$$
(3)
As some of the baryons had already collapsed into galaxies at $`z=25`$, the value of $`\mathrm{\Omega }_bh^2`$ $`=0.019`$ should strictly be considered as an upper limit to the intergalactic density parameter.
Because of the overwhelming abundance of hydrogen, the ionization of this element is of great importance for determining the physical state of the IGM. Popular cosmological models predict that most of the intergalactic hydrogen was reionized by the first generation of stars or quasars at $`z=715`$. The case that has received the most theoretical studies is one where hydrogen is ionized by the absorption of photons, $`H+\gamma p+e`$ (as opposite to collisional ionization $`H+ep+e+e`$) shortward of $`912`$Å; that is, with energies exceeding $`13.6`$eV, the energy of the Lyman edge. The process of reionization began as individual sources started to generate expanding $`\mathrm{II}`$regions in the surrounding IGM; throughout an $`\mathrm{II}`$region, H is ionized and He is either singly or doubly ionized. As more and more sources of ultraviolet radiation switched on, the ionized volume grew in size. The reionization ended when the cosmological $`\mathrm{II}`$regions overlapped and filled the intergalactic space.
PHOTOIONIZATION EQUILIBRIUM
At every point in a optically thin, pure hydrogen medium of neutral density $`n_{\mathrm{HI}}`$, the photoionization rate per unit volume is
$$n_{\mathrm{HI}}_{\nu _L}^{\mathrm{}}\frac{4\pi J_\nu \sigma _\mathrm{H}(\nu )}{h_P\nu }𝑑\nu ,$$
(4)
where $`J_\nu `$ is the mean intensity of the ionizing radiation (in energy units per unit area, time, solid angle, and frequency interval), and $`h_P`$ is the Planck constant. The photoionization cross-section for hydrogen in the ground state by photons with energy $`h_P\nu `$ (above the threshold $`h_P\nu _L=13.6`$eV) can be usefully approximated by
$$\sigma _\mathrm{H}(\nu )=\sigma _L(\nu /\nu _L)^3,\sigma _L=6.3\times 10^{18}\mathrm{cm}^2.$$
(5)
At equilibrium, this is balanced by the rate of radiative recombinations $`p+eH+\gamma `$ per unit volume,
$$n_en_p\alpha _A(T),$$
(6)
where $`n_e`$ and $`n_p`$ are the number densities of electrons and protons, and $`\alpha _A=\sigma _nv_e`$ is the radiative recombination coefficient, i.e. the product of the electron capture cross-section $`\sigma _n`$ and the electron velocity $`v_e`$, averaged over a thermal distribution and summed over all atomic levels $`n`$. At the commonly encountered gas temperature of $`10^4`$K, $`\alpha _A=4.2\times 10^{13}`$cm$`^3`$s<sup>-1</sup>.
Consider, as an illustrative example, a point in an intergalactic $`\mathrm{II}`$region at (say) $`z=6`$, with density $`\overline{n}_\mathrm{H}=(1.6\times 10^7\mathrm{cm}^3)(1+z)^3=5.5\times 10^5`$cm<sup>-3</sup>. The $`\mathrm{II}`$region surrounds a putative quasar with specific luminosity $`L_\nu =10^{30}(\nu _L/\nu )^2\mathrm{ergs}\mathrm{s}^1\mathrm{Hz}^1`$, and the point in question is at a distance of $`r=3`$Mpc from the quasar. To a first approximation, the mean intensity is simply the radiation emitted by the quasar reduced by geometrical dilution,
$$4\pi J_\nu =\frac{L_\nu }{4\pi r^2}.$$
(7)
We then have for the photoionization timescale:
$$t_{\mathrm{ion}}=\left[_{\nu _L}^{\mathrm{}}\frac{4\pi J_\nu \sigma _\mathrm{H}(\nu )}{h_P\nu }𝑑\nu \right]^1=5\times 10^{12}\mathrm{s},$$
(8)
and for the recombination timescale:
$$t_{\mathrm{rec}}=\frac{1}{n_e\alpha _A}=5\times 10^{16}\mathrm{s}\frac{\overline{n}_\mathrm{H}}{n_e}.$$
(9)
As in photoionization equilibrium $`n_{\mathrm{HI}}/t_{\mathrm{ion}}=n_p/t_{\mathrm{rec}}`$, these values imply $`n_{\mathrm{HI}}/n_p10^4`$, that is, hydrogen is very nearly completely ionized.
A source radiating ultraviolet photons at a finite rate cannot ionize an infinite region of space, and therefore there must be an outer edge to the ionized volume (this is true unless, of course, there is a population of UV emitters and all individual $`\mathrm{II}`$regions have already overlapped). One fundamental characteristic of the problem is the very small value of the mean free path for an ionizing photon if the hydrogen is neutral, $`(\sigma _Ln_\mathrm{H})^1=0.9`$kpc at threshold, much smaller than the radius of the ionized region. If the source spectrum is steep enough that little energy is carried out by more penetrating, soft X-ray photons, we have one nearly completely ionized $`\mathrm{II}`$region, sepated from the outer neutral IGM by a thin transition layer or ‘ionization-front’. The inhomogeneity of the IGM is of primary importance for understanding the ionization history of the Universe, as denser gas recombines faster and is therefore ionized at later times than the tenuous gas in underdense regions. An approximate way to study the effect of inhomogeneity is to write the rate of recombinations as
$$n_en_p\alpha _A(T)=Cn_e^2\alpha _A(T)$$
(10)
(assuming $`T`$ is constant in space), where the brackets are the space average of the product of the local proton and electron number densities, and the factor $`C>1`$ takes into account the degree of clumpiness of the IGM. If ionized gas with electron density $`n_e`$ density filled uniformly a fraction $`1/C`$ of the available volume, the rest being empty space, the mean square density would be $`n_e^2=n_e^2/C=n_e^2C`$.
The IGM is completely reionized when one ionizing photon has been emitted for each H atom by the radiation sources, and when the rate of emission of UV photons per unit (comoving) volume balances the radiative recombination rate, so that hydrogen atoms are photoionized faster than they can recombine. The complete reionization of the Universe manifests itself in the absence of an absorption trough in the spectra of galaxies and quasars at high redshifts. If the IGM along the line-of-sight to a distant source were neutral, the resonant scattering at the wavelength of the Ly$`\alpha `$($`2p1s;h_P\nu _\alpha =10.2`$eV) transition of atomic hydrogen would remove all photons blueward of Ly$`\alpha `$off the line-of-sight. For any reasonable density of the IGM, the scattering optical depth is so large that detectable absorption will be produced by relatively small column (or surface) densities of intergalactic neutral hydrogen.
GUNN-PETERSON EFFECT
Consider radiation emitted at some frequency $`\nu _e`$ that lies blueward of Ly$`\alpha `$by a source at redshift $`z_e`$, and observed at Earth at frequency $`\nu _o=\nu _e(1+z_e)^1`$. At a redshift $`z`$ such that $`(1+z)=(1+z_e)\nu _\alpha /\nu _e`$, the emitted photons pass through the local Ly$`\alpha `$resonance as they propagates towards us through a smoothly distributed sea of neutral hydrogen atoms, and are scattered off the line-of-sight with a cross-section (neglecting stimulated emission) of
$$\sigma [\nu _o(1+z)]=\frac{\pi e^2}{m_ec}f\varphi [\nu _o(1+z)],$$
(11)
where $`f=0.4162`$ is the upward oscillator strength for the transition, $`\varphi `$ is the line profile function \[with normalization $`\varphi (\nu )𝑑\nu =1`$\], $`c`$ is the speed of light, and $`e`$ and $`m_e`$ are the electron charge and mass, respectively. The total optical depth for resonant scattering at the observed frequency is given by the line integral of this cross-section times the neutral hydrogen proper density $`n_{\mathrm{HI}}(z)`$,
$$\tau _{\mathrm{GP}}=_0^{z_e}\sigma [\nu _o(1+z)]n_{\mathrm{HI}}(z)\frac{d\mathrm{}}{dz}𝑑z,$$
(12)
where $`d\mathrm{}/dz=cH_0^1(1+z)^1[\mathrm{\Omega }_M(1+z)^3+\mathrm{\Omega }_K(1+z)^2+\mathrm{\Omega }_\mathrm{\Lambda }]^{1/2}`$ is the proper line element in a Friedmann-Robertson-Walker metric, and $`\mathrm{\Omega }_M`$, $`\mathrm{\Omega }_\mathrm{\Lambda }`$, and $`\mathrm{\Omega }_K=1\mathrm{\Omega }_M\mathrm{\Omega }_\mathrm{\Lambda }`$ are the matter, vacuum, and curvature contribution to the present density parameter. As the scattering cross-section is sharply peaked around $`\nu _\alpha `$, we can write
$$\tau _{\mathrm{GP}}(z)=\left(\frac{\pi e^2f}{m_ec\nu _\alpha }\right)\frac{n_{\mathrm{HI}}}{(1+z)}\frac{d\mathrm{}}{dz}.$$
(13)
In an Einstein-de Sitter ($`\mathrm{\Omega }_M=1`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$) Universe, this becomes
$$\tau _{\mathrm{GP}}(z)=\frac{\pi e^2f}{m_eH_0\nu _\alpha }\frac{n_{\mathrm{HI}}}{(1+z)^{3/2}}=6.6\times 10^3h^1\left(\frac{\mathrm{\Omega }_bh^2}{0.019}\right)\frac{n_{\mathrm{HI}}}{\overline{n}_\mathrm{H}}(1+z)^{3/2}.$$
(14)
The same expression for the opacity is also valid in the case of optically thin (to Ly$`\alpha `$scattering) discrete clouds as long as $`n_{\mathrm{HI}}`$ is replaced with the average neutral density of individual clouds times their volume filling factor.
In an expanding Universe homogeneously filled with neutral hydrogen, the above equations apply to all parts of the source spectrum to the blue of Ly$`\alpha `$. An absorption trough should then be detected in the level of the rest-frame UV continuum of the quasar; this is the so-called “Gunn-Peterson effect”. Between the discrete absorption lines of the Ly$`\alpha `$forest clouds, quasar spectra do not show a pronounced Gunn-Peterson absorption trough. The current upper limit at $`z_e5`$ is $`\tau _{\mathrm{GP}}<0.1`$ in the region of minimum opacity, implying from equation (14) a neutral fraction of $`n_{\mathrm{HI}}/\overline{n}_\mathrm{H}<10^6h`$. Even if 99% of all the cosmic baryons fragment at these epochs into structures that can be identified with quasar absorption systems, with only 1% remaining in a smoothly distributed component, the implication is a diffuse IGM which is ionized to better than 1 part in $`10^4`$.
In modern interpretations of the IGM, it is difficult to use the Gunn-Peterson effect to quantify the amount of ionizing radiation that is necessary to keep the neutral hydrogen absorption below the detection limits. This is because, in hierarchical clustering scenarios for the formation of cosmic structures (the Cold Dark Matter model being the most studied example), the accumulation of matter in overdense regions under the influence of gravity reduces the optical depth for Ly$`\alpha `$scattering considerably below the average in most of the volume of the Universe, and regions of minimum opacity occur in the most underdense areas (expanding ‘cosmic minivoids’).
A CLUMPY IGM
Owing to the non-linear collapse of cosmic structures, the IGM is well known to be highly inhomogeneous. The discrete gaseous systems detected in absorption in the spectra of high-redshift quasars blueward of the Ly$`\alpha `$emission line are assigned different names based on the appearance of their absorption features (see Figure 1).
The term “Ly$`\alpha `$forest” is used to denote the plethora of narrow absorption lines whose measured equivalent widths imply $`\mathrm{I}`$column densities ranging from $`10^{16}`$cm<sup>-2</sup> down to $`10^{12}`$cm<sup>-2</sup>. These systems, observed to evolve rapidly with redshift between $`2<z<4`$, have traditionally been interpreted as intergalactic gas clouds associated with the era of baryonic infall and galaxy formation, photoionized (to less than a neutral atom in $`10^4`$) and photoheated (to temperatures close to 20,000 K) by a ultraviolet background close to the one inferred from the integrated contribution from quasars. Recent spectra at high resolution and high signal-to-noise obtained with the Keck telescope have shown that most Ly$`\alpha `$forest clouds at $`z3`$ down to the detection limit of the data have undergone some chemical enrichment, as evidenced by weak, but measurable $`\mathrm{IV}`$lines. The typical inferred metallicities range from 0.3% to 1% of solar values, subject to uncertainties of photoionization models. Clearly, these metals were produced in stars that formed in a denser environmemt; the metal-enriched gas was then expelled from the regions of star formation into the IGM.
An intervening absorber at redshift $`z`$ having a neutral hydrogen column density exceeding $`2\times 10^{17}`$cm<sup>-2</sup> is optically thick to photons having energy greater than 13.6 eV, and produces a discontinuity at the hydrogen Lyman limit, i.e. at an observed wavelength of $`912(1+z)`$Å. These scarcer “Lyman-limit systems” (LLS) are associated with the extended gaseous haloes of bright galaxies near the line-of-sight, and have metallicities which appear to be similar to that of Ly$`\alpha `$forest clouds.
In “damped Ly$`\alpha `$systems’ the $`\mathrm{I}`$column is so large ($`N_{\mathrm{HI}}\mathrm{}>10^{20}`$cm<sup>-2</sup>, comparable to the interstellar surface density of spiral galaxies today) that the radiation damping wings of the Ly$`\alpha `$line profile become detectable. While relatively rare, damped systems account for most of the neutral hydrogen seen at high redshifts. The typical metallicities are about 10% of solar, and do not evolve significantly over a redshift interval $`0.5<z<4`$ during which most of today’s stars were actually formed.
Except at the highest column densities, discrete absorbers are inferred to be strongly photoionized. From quasar absorption studies we also know that neutral hydrogen accounts for only a small fraction, $`10\%`$, of the nucleosynthetic baryons at early epochs.
DISTRIBUTION OF COLUMN DENSITIES AND EVOLUTION
The bivariate distribution $`f(N_{\mathrm{HI}},z)`$ of $`\mathrm{I}`$column densities and redshifts is defined by the probability $`dP`$ that a line-of-sight intersects a cloud with column density $`N_{\mathrm{HI}}`$ in the range $`dN_{\mathrm{HI}}`$, at redshift $`z`$ in the range $`dz`$,
$$dP=f(N_{\mathrm{HI}},z)dN_{\mathrm{HI}}dz.$$
(15)
As a function of column, a single power-law with slope $`1.5`$ appears to provide at high redshift a surprisingly good description over 9 decades in $`N_{\mathrm{HI}}`$, i.e. from $`10^{12}`$ to $`10^{21}\mathrm{cm}^2`$. It is a reasonable approximation to use for the distribution of absorbers along the line-of-sight:
$$f(N_{\mathrm{HI}},z)=AN_{\mathrm{HI}}^{1.5}(1+z)^\gamma .$$
(16)
Ly$`\alpha `$forest clouds and Lyman-limit systems appear to evolve at slightly different rates, with $`\gamma =1.5\pm 0.4`$ for the LLS and $`\gamma =2.8\pm 0.7`$ for the forest lines. Let us assume, for simplicity, a single redshift exponent, $`\gamma =2`$, for the entire range in column densities. In the power-law model (16) the number $`N`$ of absorbers with columns greater than $`N_{\mathrm{HI}}`$ per unit increment of redshift is
$$\frac{dN}{dz}=_{N_{\mathrm{HI}}}^{\mathrm{}}f(N_{\mathrm{HI}}^{},z)𝑑N_{\mathrm{HI}}^{}=2AN_{\mathrm{HI}}^{0.5}(1+z)^2.$$
(17)
A normalization value of $`A=4.0\times 10^7`$ produces then $`3`$ LLS per unit redshift at $`z=3`$, and, at the same epoch, $`150`$ forest lines above $`N_{\mathrm{HI}}=10^{13.8}\mathrm{cm}^2`$, in reasonable agreement with the observations.
If absorbers at a given surface density are conserved, with fixed comoving space number density $`n=n_0(1+z)^3`$ and geometric cross-section $`\mathrm{\Sigma }`$, then the intersection probability per unit redshift interval is
$$\frac{dP}{dz}=\mathrm{\Sigma }n\frac{d\mathrm{}}{dz}=\mathrm{\Sigma }n_0(1+z)^3\frac{d\mathrm{}}{dz}.$$
(18)
If the Universe is cosmologically flat, the expansion rate at early epochs is close to the Einstein-de Sitter limit, and the redshift distribution for conserved clouds is predicted to be
$$\frac{dP}{dz}(1+z)^3\frac{d\mathrm{}}{dz}(1+z)^{1/2}.$$
(19)
The rate of increase of $`f(N_{\mathrm{HI}},z)`$ with $`z`$ in both the Ly$`\alpha `$forest and LLS is considerably faster than this, indicating rapid evolution. The mean proper distance between absorbers along the line-of-sight with columns greater than $`N_{\mathrm{HI}}`$ is
$$L=\frac{d\mathrm{}}{dz}\frac{dz}{dN}\frac{cN_{\mathrm{HI}}^{1/2}}{H_0\mathrm{\Omega }_M^{1/2}2A(1+z)^{4.5}}.$$
(20)
For clouds with $`N_{\mathrm{HI}}>10^{14}\mathrm{cm}^2`$, this amounts to $`L0.7h^1\mathrm{\Omega }_M^{1/2}`$Mpc at $`z=3`$. At the same epoch, the mean proper distance between LLS is $`L30h^1\mathrm{\Omega }_M^{1/2}`$Mpc.
INTERGALACTIC CONTINUUM OPACITY
Even if the bulk of the baryons in the Universe are fairly well ionized at all redshifts $`z\mathrm{}<5`$, the residual neutral hydrogen still present in the Ly$`\alpha `$forest clouds and Lyman-limit systems significantly attenuates the ionizing flux from cosmological distant sources. To quantify the degree of attenuation we have to introduce the concept of an effective continuum optical depth $`\tau _{\mathrm{eff}}`$ along the line-of-sight to redshift $`z`$,
$$e^\tau e^{\tau _{\mathrm{eff}}},$$
(21)
where the average is taken over all lines-of-sight. Negleting absorption due to helium, if we characterize the Ly$`\alpha `$forest clouds and LLS as a random distribution of absorbers in column density and redshift space, then the effective continuum optical depth of a clumpy IGM at the observed frequency $`\nu _o`$ for an observer at redshift $`z_o`$ is
$$\tau _{\mathrm{eff}}(\nu _o,z_o,z)=_{z_o}^z𝑑z^{}_0^{\mathrm{}}𝑑N_{\mathrm{HI}}f(N_{\mathrm{HI}},z)(1e^\tau ).$$
(22)
where $`\tau =N_{\mathrm{HI}}\sigma _H(\nu )`$ is the hydrogen Lyman continuum optical depth through an individual cloud at frequency $`\nu =\nu _o(1+z)/(1+z_o)`$. This formula can be easily understood if we consider a situation in which all absorbers have the same optical depth $`\tau _0`$ independent of redshift, and the mean number of systems along the path is $`\mathrm{\Delta }N=𝑑z𝑑N/𝑑z`$. In this case the Poissonian probability of encountering a total optical depth $`k\tau _0`$ along the line-of-sight (with $`k`$ integer) is $`p(k\tau _0)=e^{\mathrm{\Delta }N}\mathrm{\Delta }N^k/(\tau _0k!)`$, and $`e^\tau =e^{k\tau _0}p(k\tau _0)=\mathrm{exp}[\mathrm{\Delta }N(1e^{\tau _0})]`$.
If we extrapolate the $`N_{\mathrm{HI}}^{1.5}`$ power-law in equation (16) to very small and large columns, the effective optical depth becomes an analytical function of redshift and wavelength,
$$\tau _{\mathrm{eff}}(\nu _o,z_o,z)=\frac{4}{3}\sqrt{\pi \sigma _L}A\left(\frac{\nu _o}{\nu _L}\right)^{1.5}(1+z_o)^{1.5}\left[(1+z)^{1.5}(1+z_o)^{1.5}\right].$$
(23)
Due to the rapid increase with lookback time of the number of absorbers, the mean free path of photons at $`912`$Å becomes so small beyond a redshift of 2 that the radiation field is largely ‘local’. Expanding equation (23) around $`z`$, one gets $`\tau _{\mathrm{eff}}(\nu _L)0.36(1+z)^2\mathrm{\Delta }z`$. This means that at $`z=3`$, for example, the mean free path for a photon near threshold is only $`\mathrm{\Delta }z=0.18`$, and sources of ionizing radiation at higher redshifts are severely attenuated.
BIBLIOGRAPHY
For a recent Deuterium abundance measurement and its implications for the baryon density parameter see
Burles, S., & Tytler, D. 1999, Ap. J., 499, 699.
A short and hardly comprehensive list of references on issues related to the reionization of the IGM includes
Arons, J., & Wingert, D. W. 1972, Ap. J., 177, 1.
Shapiro, P. R., & Giroux, M. L. 1987, Ap. J., 321, L107.
Meiksin, A., & Madau, P. 1993, Ap. J., 412, 34.
Tegmark, M., Silk, J., Rees, M. J., Blanchard, A., Abel, T., & Palla, F. 1997, Ap. J., 474, 1.
Gnedin, N. Y., & Ostriker, J. P. 1997, Ap. J., 486, 581.
Haiman, Z., & Loeb, A. 1997, Ap. J., 483, 21.
Miralda-Escude’, J. 1998, Ap. J., 501, 15.
Madau, P., Haardt, F., & Rees, M. J. 1999, Ap. J., 514, 648.
The use of Ly$`\alpha `$resonant absorption as a sensitive probe of intergalactic neutral hydrogen was predicted indipendently by
Gunn, J. E., & Peterson, B. A. 1965, Ap. J., 142, 1633.
Shklovsky, I. S. 1965, Soviet Astron., 8, 638.
Scheuer, P. A. G. 1965, Nature, 207, 963.
The quoted upper limit to the Gunn-Peterson optical depth at $`z5`$ is from
Songaila, A., Hu, E. M., Cowie, L. L., & McMahon, R. G. 1999, Ap. J., 525, L5.
The high resolution quasar spectrum shown in the Figure was taken from
Songaila, A. 1998, A. J., 115, 2184,
who also discusses the evolution of metal abundances in the Ly$`\alpha `$forest.
An extensive discussion of the physics of the intergalactic medium can be found in
Peebles, P. J. E. 1993, Principles of Physical Cosmology (Princeton University Press), chapter 23.
The physics of an ionized hydrogen gas is covered in
Osterbrock, D. E. 1989, Astrophysics of Gaseous Nebulae and Active Galactic Nuclei (University
Science Books), chapters 2 and 3.
Our present understanding of the Ly$`\alpha `$forest is summarized in
Rauch, M. 1998, Annual Rev. of Astron. and Astrophys. 36, 267.
The use of hydrodynamic cosmological simulations to make quantitative prediction of the physical state of the IGM was pioneered by
Cen, R., Miralda-Escude’, J., Ostriker, J. P., & Rauch, M. 1994, Ap. J., 437, L9.
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# MSW Solutions to the Solar Neutrino Problem in Presence of Noisy Matter Density Fluctuations
## I Introduction
Solar neutrinos were first detected already three decades ago in the Homestake experiment and from the very beginning it was pointed out the puzzling issue of the deficit in the observed rate as compared to the theoretical expectation based on the standard solar model with the implicit assumption that neutrinos created in the solar interior reach the Earth unchanged, i.e. they are massless and have only standard properties and interactions. This discrepancy led to a change in the original goal of using solar neutrinos to probe the properties of the solar interior towards the study of the properties of the neutrino itself and it triggered an intense activity both theoretical as well as experimental, with new measurements being proposed in order to address the origin of the deficit.
On the theoretical side, enormous progress has been done in the improvement of solar modelling and calculation of nuclear cross sections. For example, helioseismological observations have now established that diffusion is occurring and by now most solar models incorporate the effects of helium and heavy element diffusion . From the experimental point of view the situation is now much richer. Four additional experiments to the original Chlorine experiment at Homestake have also detected solar neutrinos: the radiochemical Gallium experiments on $`pp`$ neutrinos, GALLEX and SAGE , and the water Cerenkov detectors Kamiokande and Super–Kamiokande . The latter have been able not only to confirm the original detection of solar neutrinos at lower rates than predicted by standard solar models, but also to demonstrate directly that the neutrinos come from the Sun by showing that recoil electrons are scattered in the direction along the Sun–Earth axis. Moreover, they have also provided us with useful information on the time dependence of the event rates during the day and night, as well as a measurement of the recoil electron energy spectrum. After 825 days of operation, Super–Kamiokande has also presented preliminary results on the seasonal variation of the neutrino event rates, an issue which will become important in discriminating the MSW scenario from the possibility of neutrino oscillations in vacuum . At the present stage, the quality of the experiments themselves and the robustness of the theory give us confidence that in order to describe the data one must depart from the Standard Model (SM) of particle physics interactions by endowing neutrinos with new properties. In theories beyond the SM, neutrinos may naturally have new properties, the most generic of which is the existence of mass. It is undeniable that the most popular explanation of the solar neutrino anomaly is in terms of neutrino masses and mixing leading to neutrino oscillations either in vacuum or via the matter-enhanced MSW mechanism .
The standard MSW analysis is based on a mean–field treatment of the solar background through which the neutrinos propagate. In this approximation the global analysis of the full neutrino data sample described above leads to the existence of three allowed regions in the $`\mathrm{\Delta }m^2`$ $`\mathrm{sin}^22\theta `$ parameter space for neutrino oscillations
* non-adiabatic-matter-enhanced oscillations or small mixing angle (SMA) region with $`\mathrm{\Delta }m^2=(0.4`$$`1)\times 10^5`$ eV<sup>2</sup> and $`\mathrm{sin}^2(2\theta )=(1`$$`10)\times 10^3`$, and
* large mixing (LMA) region $`\mathrm{\Delta }m^2=(0.2`$$`5)\times 10^4`$ eV<sup>2</sup> and $`\mathrm{sin}^2(2\theta )=0.6`$$`1`$.
* low mass solution (LOW) $`\mathrm{\Delta }m^2=(0.3`$$`2)\times 10^7`$ eV<sup>2</sup> and $`\mathrm{sin}^2(2\theta )=0.8`$$`1`$.
There are several works in the literature where corrections to such mean–field picture have been studied. The influence of periodic matter density fluctuations of given amplitude and fixed frequency above the average density on resonant neutrino conversion was investigated in Refs.. In Ref. a parametric resonance is found when the fixed frequency of the perturbation is close to the neutrino oscillation frequency. This approach however gives not answer to the physical origin of such fixed frequency perturbation.
More recently, the main approach to fluctuations that has been pursued is to model the matter density as a Gaussian random variable (white noise). In this approach the number of free parameters remains the same as in the case of fixed frequency perturbations – two, the perturbation amplitude and the correlation length– but the presence of white noise in the sun is doubtless since there are many mechanism to generate random density perturbations. For technical reasons these analysis where performed under the assumption that fluctuations have spatial correlations only over distances small compared to the neutrino oscillation lengths. Within this approximation the conclusions obtained were that such fluctuations in the solar electron density can significantly modify the MSW solutions to the solar neutrino problem (SNP) provided that their relative amplitude near the MSW resonance point can be as large as few percent.
In Ref. criticisms to these results were raised based on two facts: i) the unexistence of a plausible source for such $`\delta `$-correlated fluctuations in the vicinity of the MSW resonance point and ii) whatever its origin, the effect of the density perturbation is maximum in the regime where the short correlation–length approximation fails. In particular in Ref. they concentrate on helioseismological waves as origin of the perturbation and in particular on g-waves whose amplitude increases with the solar depth and can, in principle, reach the interesting values to affect neutrino propagation. However they conclude that such g-waves do not affect the MSW neutrino conversions since the wavelength for the lower modes, for which the largest amplitudes are possible, is much longer than the characteristic MSW neutrino oscillation length.
It has been recently argued , however, that in the magnetohydrodynamical (MHD) generalization of the Helioseismology the objection in Ref. does not hold. Assuming modest central large-scale magnetic fields ($`B_0=`$1–10–100 Gauss) one can find magneto-gravity eigenmodes with much shorter wave lengths for density perturbations $`\lambda _{MHD}2002000`$km that is comparable with the neutrino oscillation length at the MSW resonance for large and small mixing angles correspondingly.
In this paper we revisit the problem of the effect of matter of density fluctuations in the sun on the MSW solutions to the SNP. In our approach we assume no specific mechanism for generation of the fluctuation and we keep the amplitude and correlation length as independent parameters. There are two main differences in our analysis as compared to those in Refs. . First we do not work under the approximation that fluctuations have spatial correlations only over distances small compared to the neutrino oscillation lengths. Instead we solve numerically the evolution equation for the neutrino system including the full effect of the random matter density fluctuations of given amplitude and correlation length. Second, in order to establish the possible effect on the MSW solutions to the SNP we perform a global analysis of all the existing observables including not only the measured total rates but also the Super–Kamiokande measurement on the time dependence of the event rates during the day and night, as well as the recoil electron energy spectrum. In particular we include the regeneration effects when neutrinos cross the Earth which were neglected in Ref. .
The outline of the paper is as follows. In Sec. II we discuss our approach to the solution of the neutrino evolution in the presence of density fluctuations. In Sec. II A we briefly summarize the standard analytical approach based on the short correlation length approximation and in Sec. II B we discuss our numerical treatment and present our results for the survival probabilities as a function of the relevant oscillation and noise parameters. In Sec. II C we interpret our results for the enhancement of the survival probability in the language of parametric resonance in MSW conversions. Section III is devoted to the statistical analysis of the solar neutrino observables in the framework of the MSW solutions of the SNP in presence of the noisy density fluctuations. We study the variation of the allowed regions of the SNP for different combinations of observables when noise fluctuations with different correlation lengths are included. Our results are summarized in Figs. 58 and Table II. We show that even for noise levels as large as 4% the relative quality of the three allowed regions of the MSW solutions to the SNP depends on the value of the correlation length studied in the wide range $`L_0=70`$$`10^4`$ Km and that the three allowed regions of the MSW solutions of the SNP remain valid at the present level of solar neutrino experiments. Finally in Sec. IV we discuss our results and summarize our conclusions.
## II MSW solutions to SNP for noisy matter density in the Sun.
Some mechanisms to rise density perturbations in the Sun were proposed in literature in connection with the MSW solution to the SNP. One of them with g-waves in Helioseismology as a plausible source of matter noise fails being applied to the MSW neutrino conversions since the wavelength of such modes happens to be too long, $`\lambda _g0.1R_{}l_{osc}`$ to influence neutrino oscillations. This $`\lambda _g`$ is the wave length for low radial degree $`n3`$ for which largest g-mode amplitudes $`\delta \rho (r)/\rho _0(r)4`$ % are possible .
However, in the MHD generalization of the Helioseismology such objection does not apply. Assuming modest central large-scale magnetic fields ($`B_0=110100`$ Gauss) one can find magneto-gravity eigenmodes with much shorter wave lengths for density perturbations $`\lambda _{MHD}l_{osc}2002000`$ km comparable with the neutrino oscillation length at the MSW resonance for large and small mixing angles correspondingly.
Note that standard Helioseismology corrections to the standard solar model (SSM) (neglecting magnetic fields in the Sun) give density fluctuations deep in solar interior at a low level $`\delta \rho (r)/\rho _0(r)1`$% . This analysis is done using the solution of the inverted problem in Helioseismology based on the corresponding integral equation for which the kernels are built on the full set of p-mode waves calculated and observed on the photosphere . However, such analysis does not touch both g-modes that are still invisible on the surface of the Sun (from SOHO satellite) and has no relation to MHD modes found in .
### A Generalized Parke formula for averaged evolution equation
In this subsection we show that the Schrödinger equation approach for noisy matter generalized here for an arbitrary matter density perturbation correlator (not only $`\delta `$-correlator like in ) is equivalent to the Redfield evolution equation leading to the generalized Parke formula for the survival probability in the presence of noise .
Let us assume the presence of regular density waves $`Re\left({\displaystyle \underset{n}{}}C_n\delta \rho _n(t)\right)`$ exited somehow within the solar interior. They could be, for instance, the MHD matter density waves, $`\delta \rho _n(z)`$, which appear in the 1-dimensional solar model with the exponential matter background profile $`\rho _0(z)\mathrm{exp}(z/H)`$, $`H0.1R_{}`$, in the presence of gravity g = (0,0, -g) and an external constant magnetic field B = $`(B_0,0,0)`$. Such waves obey the dispersion relation $`\omega _n=\omega (n,B_0,k_x,k_y)`$ with the periods $`T=2\pi /\omega _n`$ few days and they are quite different from the g-modes in Helioseismology. In particular, unlike the helioseismological g-modes, they have very short wavelength along the z-axis $`\lambda _zR_{}/n\lambda _g`$ for large node numbers $`n1`$ acceptable in the model .
Assuming these density perturbations added to the SSM background density profile $`\rho _0(t)`$, we write the master Schrödinger equation for MSW conversions of two neutrino flavors, $`\nu _e\nu _y`$,
$$i\left(\begin{array}{c}\dot{\nu }_e\hfill \\ \dot{\nu }_y\hfill \end{array}\right)=\left(\begin{array}{cc}H_e& s_2\delta \\ s_2\delta & 0\end{array}\right)\left(\begin{array}{c}\nu _e\\ \nu _y\end{array}\right),$$
(1)
where in the diagonal entry $`H_e=V_{ey}(t)[1+Re\left({\displaystyle \underset{n}{}}C_n\varphi _n(t)\right)]2c_2\delta `$ any density eigenmode $`\xi _n(t)=C_n\varphi _n(t)=\delta \rho _n(t)/\rho _0(t)`$ has the small amplitude, $`\xi _n(t)1`$. $`c_2=\mathrm{cos}2\theta `$, $`s_2=\mathrm{sin}2\theta `$ and $`\delta =\mathrm{\Delta }m^2/4E`$ are the neutrino mixing parameters; $`V_{ey}(t)=G_F\sqrt{2}(\rho (t)/m_p)(1Y_n)`$ and $`V_{ey}=V_{es}(t)=G_F\sqrt{2}(\rho (t)/m_p)(13Y_n/2)`$ are the neutrino vector potentials in the Sun for active-active (y= x) and for active-sterile neutrino conversions correspondingly. They are given by the neutron abundance $`Y_n=m_pN_n(t)/\rho (t)`$ where for neutral matter the relation $`Y_e=Y_p=1Y_n`$ has been used and the SSM density profile $`\rho (t)`$ is given by BP98 model. In what follows we will discuss only conversion into active neutrinos.
Using the survival probability $`P_{ee}=\nu _e^{}\nu _e`$ and the auxiliary functions $`I=Im(\nu _e^{}\nu _y)`$, $`R=Re(\nu _e^{}\nu _y)`$ one can derive from the master equation above the equivalent system of dynamical equations. After averaging of those dynamical equations over small density perturbations $`\xi _n(t)1`$, such system takes the form:
$$i\frac{d}{dt}\left(\begin{array}{c}\hfill \\ \hfill \\ 𝒫\text{1/2}\hfill \end{array}\right)=\left(\begin{array}{ccc}2\kappa (t)& (V_{ey}(t)2\delta c_2)& 0\\ V_{ey}(t)2\delta c_2& 2\kappa (t)& \delta s_2\\ 0& +\delta s_2& 0\end{array}\right)\left(\begin{array}{c}\\ \\ 𝒫\text{1/2}\end{array}\right).$$
(2)
with $`𝒫=<𝒫_{}>`$, $`=<>`$, $`=<>`$. This system of equations is similar to Eq. (3.14) in but here the matter perturbation parameter $`\kappa (t)`$ is of the form:
$$\kappa (t)=\frac{1}{2}\underset{m,n}{}_{t_1}^tV_{ey}(t)V_{ey}(t_2)C_nC_m\varphi _n(t)\varphi _m(t_2)𝑑t_2,$$
(3)
where in averaging $`\mathrm{}`$ we use that: (i) $`C_n`$ are uncorrelated random variables which are Gaussian distributed with vanishing mean: $`C_n=0`$, (ii) different modes (with different node number) are uncorrelated, $`C_nC_m=A_n\delta _{nm}`$.
It is reasonable to assume that even for a regular eigenmode “n” excited somehow within the solar interior, different neutrinos emitted from different starting points $`t_0`$ in the core propagate to the Earth along parallel rays crossing different profile realizations of the same mode $`\xi _n(t)`$. The phase of the wave entering in $`C_n`$ is random, $`C_n=0`$, as well as a starting point $`t_0`$ is random for given t). In other words, the averaging over random phases is equivalent to the averaging over production point distribution (see the discussion in Sec. II B).
Thus, one can consider the sum of multimode regular perturbations as random density perturbations $`C_n\varphi _n(t)\xi (t)`$, $`\xi (t)=0`$, for which the $`\delta `$-correlated noise
$$\xi (t_1)\xi (t_2)=L_0\xi (t_1)^2\delta (t_1t_2)$$
is the particular case leading from Eq. (3) to $`\kappa (t)=V_{ey}^2(t)\xi ^2L_0/4`$ .
We can now identify the entries in the Hamiltonian Eq. (2) with the corresponding terms in Eq. (77) derived in Ref. through the averaging of the Redfield equation for the density matrix. One finds full coincidence with the notation in Ref. : $`\rho _1=`$, $`\rho _2=`$, $`\rho _3=𝒫\text{1/2}`$ and $`a(t)=\kappa (t)`$, $`2(M_3+b)=V_{ey}(t)\delta c_2`$, $`M_2=0`$, $`M_1=\delta s_2`$.
For slowly varying variables $`V_{ey}(t)`$ and $`\kappa (t)`$, one can diagonalize the Hamiltonian in Eq. (2) obtaining the generalized Parke formula (Eq. (86) in Ref. )
$$P_{ee}(t)=\frac{1}{2}+\left(\frac{1}{2}P_J\right)\mathrm{exp}\left(2_{t_0}^t\gamma _0(x)𝑑x\right)\mathrm{cos}2\theta _m(t_0)\mathrm{cos}2\theta _m(t)$$
(4)
where the effect of the density fluctuations is contained in the $`\gamma _0`$ factor
$$\gamma _0(x)=\frac{4\kappa (x)\delta ^2s_2^2}{4\delta ^2s_2^2+(V_{ey}(x)\delta c_2)^2}.$$
(5)
$`P_J`$ is the probability for level crossing as one passes the resonant point and $`\theta _m`$ is the mixing angle in matter.
### B Computer simulation of the master equation
There are some remarks to the analytic approach shown above that forced us to perform a numerical calculation. First, the generalized Parke formula in Eq. (5) does not describe neutrino conversions for large correlation lengths $`L_0l_{osc}`$ in an appropriate way giving a huge discrepancy with the results from numerical calculations . Alternatively, in Ref. , the survival probability was evaluated using the standard Parke formula for each neutrino ray and then averaging this result over 200 random density profiles of the Cell type of length $`L_0`$. Below we refer to this procedure as the “Cell” model. Their result for large $`L_0`$, however, presents a strong dependence on the resonance position within a Cell.
In our approach we directly solve the Schröedinger equation for the two–neutrino system with the following procedure. Not appealing to any origin of the matter density perturbations we consider different levels of the random density fluctuations $`\xi (r)\delta \rho (r)/\rho _0(r)=0`$ added to the background matter density $`\rho _0(r)`$ in SSM which we take to be the BP98 density
$$\rho (r)=\rho _0(r)[1+\xi (r)],$$
(6)
where the parameter $`\xi \sqrt{\xi (r)^2}`$ measures the amplitude of the perturbation. For a given value of $`\xi `$ and correlation length $`L_0`$ we generate a density profile of the type in Eq.(6). The function $`\xi (r)`$ is constructed as steps of constant values $`\xi \mathrm{\_}i`$ for each step $`i`$ of length $`L_0`$ from the production point to the edge of the Sun. The $`\xi _i`$ numbers are randomly generated following a Gaussian distribution of mean $`\xi (r)`$=0 and dispersion $`\xi `$. For illustration in Fig. 1 we show the numerical profile generated using this procedure for $`L_0=700`$ km and $`\xi =0.1`$.
Substituting this realization of the matter density into the Schrödinger equation for two neutrino flavors we have solved the Cauchy problem for different starting points $`r_0`$ within the solar core $`r_0R_{core}=0.3R_{}`$ with the same initial condition $`\nu _e(r_0)=1`$. The production points $`r_0`$ are chosen to be placed in knots of a $`30\times 60`$-net that covers the cross-section of the core hemisphere, $`r<0.3R_{}`$. In other words, $`0.01R_{}7000`$ km is the cell size of a k-rectangle chosen within core.
In this way, for $`30\times 60=1800`$ $`r_0`$-points, we have obtained a set of the complex wave functions, $`\nu _a(r,r_0)=\nu _a(r,r_0)\mathrm{exp}(i\mathrm{\Phi }(r,r_0))`$, a=e,x, from which one can easily get the survival probability at the surface of the Sun, $`P_{ee}(R_{},r_0)P_{ee}(R_{},r_0,s_2^2,\delta ,L_0,\xi )`$ and the survival probability on the day–side of the Earth after propagation in vacuum through the solar wind (e.g. for $`\nu _e\nu _\mu `$ oscillations),
$`P_{ee}^{day}(r_0)`$ $``$ $`P_{ee}^{day}(r_0,s_2^2,\delta ,L_0,\xi )=P_{ee}(R_{},r_0)+{\displaystyle \frac{s_2^2}{2}}[12P_{ee}(R_{},r_0)]`$ (8)
$`{\displaystyle \frac{1}{2}}s_2c_2[\nu _e^{}(R_{},r_0)\nu _\mu (R_{},r_0)+\nu _\mu ^{}(R_{},r_0)\nu _e(R_{},r_0)].`$
Then, assuming spherical symmetry, we have averaged the survival probability in Eq. (8) over the production points $`r_0`$. The averaging means the multiplication by the weight factor defined as the local $`\nu `$-source distribution $`S_i(r_{0k})`$, given by SSM BP98 for each $`r_{ok}`$ and each neutrino flux type i=pp, Be, pep,…, resulting in:
$$P_{ee}^{day,i}(s_2^2,\delta ,L_0,\xi )=\frac{1}{1800}\underset{k=1}{\overset{1800}{}}P_{ee}^{day}(r_{0k})S_i(r_{0k}).$$
(9)
Notice that Eq. (8) is equivalent to the standard expression for the day–side survival probability $`P_{ee}^{day}(r_0)=P_{e1}^{Sun}(r_0)P_{1e}^{Earth}+P_{e2}^{Sun}(r_0)P_{2e}^{Earth}=c^2\psi _ecs\psi _\mu ^2+s^2\psi _e+cs\psi _\mu ^2`$ where the complex wave functions $`\psi _a=\nu _a(R_{},r_0)`$ are given at the surface of the Sun and during the day $`P_{1e}^{Earth}=1P_{2e}^{Earth}=\mathrm{cos}^2\theta `$. During the night, solar neutrinos cross the Earth before reaching the detector and regeneration of $`\nu _e`$’s is possible . In order to take into account this effect we compute the probability $`P_{2e}^{Earth}`$ by integrating numerically the differential equation that describes the evolution of neutrino flavors in the Earth. This probability depends on the amount of Earth matter travelled by the neutrino, or in other words, in its arrival direction which is usually parametrized in terms of the zenith angle $`\mathrm{\Phi }`$. Thus, in general, the survival probability for a neutrino of given source $`i`$ arriving at a given zenith angle $`\mathrm{\Phi }`$ is given by
$`P_{ee}^{\mathrm{\Phi },i}(s_2^2,\delta ,L_0,\xi )`$ $`=`$ $`P_{ee}^{day,i}(s_2^2,\delta ,L_0,\xi )`$ (11)
$`+{\displaystyle \frac{(2P_{ee}^{day,i}(s_2^2,\delta ,L_0,\xi )1)\left(s^2P_{2e}^{Earth}(s_2^2,\delta ,L_0,\xi ;\mathrm{\Phi })\right)}{cos2\theta }},`$
Such probability remains a function of the fundamental neutrino parameters $`\delta `$ and $`s_2^2`$, as well as of two noise parameters $`\xi `$, $`L_0`$.
In principle we should also average over different realizations of the density perturbations with the same level of noise $`\xi `$ and correlation length $`L_0`$. We discuss next that our averaging procedure over the grid of production points $`r_0`$ is equivalent to the numerical average over an ensemble of electron densities.
Notice that for parallel rays which are directed along the z-axis to the Earth at a fixed distance $`z=z_0`$ from the center, the density profile Eq. (6) has different density amplitudes since $`z=r_1`$ only for one ray in equiliptics, for other k rays ($`k1`$) the hypotenuse is longer, $`r_k>z`$. Thus, rays are not equivalent to each other. This means that considering parallel rays and for any k-ray substituting the same final distance $`z=R_{}`$ into the wave functions $`\nu _a(R_{},r_{0k})`$ we automatically took into account different density profile realizations including a matter noise. Then integrating (summing) over $`r_0`$ in Eq. (9) we have averaged simultaneously over different realizations of noise.
Alternatively we may think of our averaging procedure in the following way. In our distribution of starting points in the grid we have several points $`j`$ with the same value of $`r_0`$ but located at different distances $`y_j`$ from the Sun–Earth $`z`$ axis. All these points encounter different realizations of the matter density in their way as they move in their path parallel to the $`z`$ axis since for each of them $`r_j(t)=\sqrt{z_j(t)^2+y_j^2}`$ is different and so it is the profile $`\rho `$ they are subject to at each time $`t`$. So our averaging over 1800 initial points can be understood as an integration over the starting point $`r_0`$ times an average over different matter density realizations $`j`$ for each $`r_0`$.
In Fig.2 we show the averaged survival probabilities $`<P_{ee}^{day}(s_2^2,\delta ,L_0,\xi )>`$ averaged for the $`{}_{}{}^{8}B`$ production point distribution, as a function of $`\delta =\mathrm{\Delta }m^2/4E`$ for two values of the mixing angle $`s_2^2=0.006`$ (SMA) and $`s_2^2=.78`$ (LMA) for level noise $`\xi =4`$ % and different values of the correlation length $`L_0`$. Also shown in the figure is the corresponding probability for the noiseless case. As shown in the figure, even for modest noise level $`\xi =4`$ % we find relatively large effects. This is specially the case for the SMA value. For LMA the effect is large only for short correlation lengths. We next discuss the interpretation of this behaviour.
### C Parametric resonance for MSW conversions in noisy matter
In Fig. 3 we show isolines $`\mathrm{\Delta }P_{ee}=[<P_{ee}^{day}(s_2^2,\delta ,L_0,\xi )><P_{ee}^{day}(s_2^2,\delta )>]`$ (averaged for the $`{}_{}{}^{8}B`$ production point distribution) in the plane $`\delta `$, $`L_0`$ for noise level $`\xi =4`$ %. $`P_{ee}^{day}(s_2^2,\delta )`$ is the survival probability in the absence of noise. Figure 3.a corresponds to fixed $`s_2^2=0.0063`$ while Fig. 3.b corresponds to $`s_2^2=0.79`$.
One can see in Fig. 3.a a wide spectrum of the domain sizes $`L_0`$, for which the noisy solution gives large difference for the $`\nu _e`$\- suppression, $`\mathrm{\Delta }P_{ee}=0.30.5`$, even for a modest noise level (=4 % )!. Largest enhancement occurs for values of $`\delta =(1`$–3$`)\times 10^6`$ eV<sup>2</sup>/MeV. For characteristic values of $`\mathrm{\Delta }m^2(0.4`$$`1)\times 10^5eV^2`$ in the SMA region the energy values for such enhancement in Fig. 3.a corresponds to $`E=\mathrm{\Delta }m^2(eV^2)/(4\delta )`$ MeV $`0.33`$–2 MeV in the interesting range for the existing solar neutrino experiments.
For the larger mixing angle the enhancement is smaller (see Fig. 3.b). The maximum effect of about $`(0.13)`$ appears somewhere near $`L_070100`$ km.
For both cases we can interpret this enhancement as being due to parametric effects in neutrino oscillations. The parametric resonance implies a synchronization between the system eigen–oscillations having MSW frequencies $`\omega =2\pi /l_m`$ and the parameter variations given by a changing size of density fluctuations $`L_0`$, $`\omega =2\pi /L_0`$. Here the neutrino oscillation length in medium $`l_m`$ is given by
$$l_m(r)=\frac{l_\nu }{\sqrt{(\mathrm{cos}2\theta l_\nu /l_0(r))^2+\mathrm{sin}^22\theta }},$$
(12)
where $`l_\nu =4\pi E/\mathrm{\Delta }m^2`$ is the oscillation length in vacuum, $`l_0(r)=2\pi m_p/\sqrt{2}G_F\rho (r)`$ is the refraction length.
The synchronization condition (parametric resonance condition) states that the phase acquired by the oscillating neutrino state on one fluctuation length $`L_0`$ should be a multiple of 2$`\pi `$
$$_0^{L_0}𝑑r\frac{2\pi }{l_m(r)}=2\pi k,k=1,2,3,\mathrm{}$$
(13)
One can simplify this condition for SMA MSW when the neutrino energy differs considerably from the MSW resonance energy, $`\mathrm{cos}2\theta l_\nu /l_0\mathrm{sin}2\theta `$. In this approximation substituting the mean density $`\rho =\overline{\rho }`$ we obtain from Eq. (13) the simple formula for parametric resonance in the case of SMA MSW oscillations
$$\frac{\mathrm{cos}2\theta }{l_\nu }\frac{1}{\overline{l_0}}=\frac{k}{L_0}.$$
(14)
According to Eq. (14), two types of resonance regimes may appear depending on the value of $`L_0`$ :
* A resonance at low energy when $`\delta \delta _R=G_F\rho _R/(\sqrt{2}m_p\mathrm{cos}2\theta )`$ ($`\overline{\rho }\rho _R`$). In this case, for $`L_0\overline{l_0}`$, one obtains the parametric resonance condition $`l_ml_\nu L_0`$.
* A resonance at energy well above the MSW resonance, $`\delta \delta _R`$ ($`\overline{\rho }>\rho _R`$). In this case the parametric resonance condition is $`l_ml_0L_0`$.
A more interesting case occurs when the synchronization condition in Eq. (13) is achieved in the proximity of the MSW resonance ($`\delta =\delta _R`$). In this case we can substitute $`l_m^{res}=l_\nu /\mathrm{sin}2\theta `$, and still assuming a short correlation length, $`L_0<\mathrm{\Delta }r`$, compared to the thickness of the resonant layer $`\mathrm{\Delta }r`$, we get much longer $`L_0`$-sizes
$$L_0=kl_m^{res}=\frac{k\pi }{\delta _R\mathrm{sin}2\theta }.$$
(15)
The parametric resonance is stronger for this case .
Since for SMA MSW conversions the thickness of the resonant layer is $`\mathrm{\Delta }r=\rho (d\rho /dr)^1/\mathrm{cot}2\theta 0.1R_{}/\mathrm{cot}2\theta 7000`$ km the condition in Eq. (15) remains valid for some first numbers k = 1,2,…. In other words, the condition 7000 km $`L_0=k\pi /\delta _R\mathrm{sin}2\theta `$ implies a lower bound on the possible values of $`\delta _R`$ for which the parametric resonance is possible. For instance for the mixing angle in Fig. 3.a $`\mathrm{sin}(2\theta )=0.08`$, the value $`\delta _R10^6`$ eV<sup>2</sup>/MeV is only marginally allowed for $`k=1`$.
Conversely from the MSW condition $`l_m^{res}=250`$ km(E/MeV)$`/(\mathrm{\Delta }m^2/10^5`$ eV<sup>2</sup>)$`/\mathrm{sin}2\theta )`$ substituting the SMA MSW parameter values $`\mathrm{\Delta }m^210^5`$ eV<sup>2</sup>, $`\mathrm{sin}2\theta 0.1`$ we find $`l_m^{res,SMA}2500(E`$/MeV) km. Moreover we see in Fig. 3.a that domain sizes $`L_0=7002000`$ km are appropriate for maximum enhancement which implies that the MSW resonant parametric condition Eq. (15) is mainly achieved for $`k=1`$ in the case of low neutrino energies like the pp-neutrinos seen in GALLEX and SAGE.
Let us note that the parametric resonant condition in Eq. (15) is rather opposite to what was assumed in Eq. (3.18) , $`L_0=0.1l_m`$. In other words, $`\delta `$-correlated noise, $`L_0l_m`$, is an analytic approximation which is appropriate only for small values of the correlation length. To illustrate this we plot in Fig. 4 the survival probability $`P_{ee}(s_2^2,\delta ,L_0,\xi )`$ as the function of $`L_0`$ for different values of the level noise and fixed $`\delta =10^6`$ eV<sup>2</sup>/MeV. Figure 4.a corresponds to $`s_2^2=0.01`$ ($`l_m^{res}7000`$ km) and Fig. 4.b to $`s_2^2=0.7`$ ($`l_m^{res}700`$ km).
One can see in Fig. 4 different heights of the “bumps” of the survival probability for different noise levels whereas the position of the peaks are similar cutting sharply at the right edge $`L_0=l_m^{res}`$ \- position of the strongest parametric resonance Eq. (15) for $`k=1`$. If $`\delta \delta _R`$, or energy E is far from the MSW resonance value $`E_R`$ the parametric resonance still takes place but from corresponding Eq.(14) we find that lower $`L_0l_m<l_m^{res}`$ are appropriate. Only if $`L_0l_m`$ the $`\delta `$-correlated regime starts.
It is interesting to compare our results in Fig. 4.a with the numerical results obtained for the “Cell” model plotted in Fig. 1 of Ref. where authors averaged the ordinary Parke formula over random ensemble of 200 density profiles of Cell type (dashed, dot-dashed and solid thick lines on that figure). These 200 profiles correspond to 200 neutrino creation sites within the core only in contrast with 1800 points in our direct numerical simulation of the Schrödinger equation Eq. (1).
Notice that in that figure $`\delta =2\delta _{Fig.\text{4}}`$. We find that both figures present similar behaviour in the short correlation length regime and a comparable enhancement for $`L_010^2`$$`10^4`$ which ends at $`L_0l_m^{res}`$. But for $`L_0l_m^{res}`$ the results in Fig. 1 of Ref. present a strong dependence on the position of the resonance within the cell. The result of our numerical calculation is closer to the “Cell” model result with randomly distributed cell positions (dot-dashed line in Fig. 1 of Ref. ) and the mean $`P_{ee}(L_0)`$ tends to the noiseless MSW survival probability as expected for large $`L_0`$.
Concluding this section we want to remark that the numerical approach presented here with the averaging of the solutions $`P_{ee}(r_0)`$ over noise realizations is quite different from any previous approaches with the averaging of the Schrödinger equation itself before obtaining of a solution. Only the straightforward numerical solution of the Schrödinger equation is an appropriate way to tackle the problem in the general case of arbitrary correlation lengths.
## III Fits: Results
### A Data and Techniques
In order to study the possible values of neutrino masses and mixing for the oscillation solution of the solar neutrino problem in the presence of noisy matter density fluctuations in the sun, we have used data on the total event rates measured in the Chlorine experiment at Homestake , in the two Gallium experiments GALLEX and SAGE and in the water Cerenkov detectors Kamiokande and Super–Kamiokande shown in Table 5. Apart from the total event rates, we have in this last case the zenith angle distribution of the events and the electron recoil energy spectrum, all measured with their recent 825-day data sample .
For the calculation of the theoretical expectations we use the BP98 standard solar model of Ref. . The general expression of the expected event rate in the presence of oscillations in experiment $`i`$ is given by $`R_i^{th}`$ :
$`R_i^{th}`$ $`=`$ $`{\displaystyle \underset{k=1,8}{}}\varphi _k{\displaystyle }dE_\nu \lambda _k(E_\nu )\times [\sigma _{e,i}(E_\nu )P_{ee}^k(s_2^2,\delta ,L_0,\xi )`$ (17)
$`+\sigma _{x,i}(E_\nu )(1P_{ee}^k(s_2^2,\delta ,L_0,\xi ))].`$
where $`E_\nu `$ is the neutrino energy, $`\varphi _k`$ is the total neutrino flux and $`\lambda _k`$ is the neutrino energy spectrum (normalised to 1) from the solar nuclear reaction $`k`$ with the normalization given in Ref. . Here $`\sigma _{e,i}`$ ($`\sigma _{x,i}`$) is the $`\nu _e`$ ($`\nu _x`$, $`x=\mu ,\tau `$) interaction cross section in the Standard Model with the target corresponding to experiment $`i`$. For the Chlorine and Gallium experiments we use improved cross sections $`\sigma _{e,i}(E)`$ from Ref. . For the Kamiokande and Super–Kamiokande experiment we calculate the expected signal with the corrected cross section as explained below. $`P_{ee}^k`$ is is the yearly averaged $`\nu _e`$ survival probability at the detector, as given by Eq. (11) after averaging over arrival directions $`\mathrm{\Phi }`$. Note that $`P_{ee}^k`$ is a function of the oscillation parameters as well as the noise parameters.
We have also included in the fit the experimental results from the Super–Kamiokande Collaboration on the zenith angle distribution of events taken on 5 night periods and the day averaged value . We compute the expected event rate in the period $`a`$ in the presence of MSW oscillations as,
$`R_{sk,a}^{th}`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Delta }\tau _a}}{\displaystyle _{\tau (\mathrm{cos}\mathrm{\Phi }_{min,a})}^{\tau (\mathrm{cos}\mathrm{\Phi }_{max,a})}}d\tau {\displaystyle \underset{k=1,8}{}}\varphi _k{\displaystyle }dE_\nu \lambda _k(E_\nu )\times [\sigma _{e,sk}(E_\nu )P_{ee}^k(s_2^2,\delta ,L_0,\xi ;\tau )`$ (19)
$`+\sigma _{x,sk}(E_\nu )(1P_{ee}^k(s_2^2,\delta ,L_0,\xi ;\tau ))],`$
where $`\tau `$ measures the yearly averaged length of the period $`a`$ normalized to 1, so $`\mathrm{\Delta }\tau _a=\tau (\mathrm{cos}\mathrm{\Phi }_{max,a})\tau (\mathrm{cos}\mathrm{\Phi }_{min,a})=`$ .500, .086, .091, .113, .111, .099 for the day and five night periods. Notice that the dependence of $`P_{ee}^i(s_2^2,\delta ,L_0,\xi ;\tau )`$ on $`\tau `$ comes only from the dependence of the Earth regeneration probability $`P_{2e}^{Earth}`$ on the different Earth matter profile crossed by the neutrino during the five night periods.
The Super-Kamiokande Collaboration has also measured the recoil electron energy spectrum. In their published analysis after 504 days of operation they present their results for energies above 6.5 MeV using the Low Energy (LE) analysis in which the recoil energy spectrum is divided into 16 bins, 15 bins of 0.5 MeV energy width and the last bin containing all events with energy in the range 14 MeV to 20 MeV. Below 6.5 MeV the background of the LE analysis increases very fast as the energy decreases. Super–Kamiokande has designed a new Super Low Energy (SLE) analysis in order to reject this background more efficiently so as to be able to lower their threshold down to 5.5 MeV. In their 825-day data they have used the SLE method and they present results for two additional bins with energies between 5.5 MeV and 6.5 MeV. In our study we use the experimental results from the Super–Kamiokande Collaboration on the recoil electron spectrum divided in 18 energy bins, including the results from the LE analysis for the 16 bins above 6.5 MeV and the results from the SLE analysis for the two low energy bins below 6.5 MeV. The general expression of the expected rate in a bin in the presence of oscillations, $`R^{th}`$, is similar to that in Eq. (17), with the substitution of the cross sections with the corresponding differential cross sections folded with the finite energy resolution function of the detector and integrated over the electron recoil energy interval of the bin, $`T_{\text{min}}TT_{\text{max}}`$:
$$\sigma _{\alpha ,sk}(E_\nu )=_{T_{\text{min}}}^{T_{\text{max}}}𝑑T_0^{\frac{E_\nu }{1+m_e/2E_\nu }}𝑑T^{}Res(T,T^{})\frac{d\sigma _{\alpha ,sk}(E_\nu ,T^{})}{dT^{}}.$$
(20)
The resolution function $`Res(T,T^{})`$ is of the form :
$$Res(T,T^{})=\frac{1}{\sqrt{2\pi }(0.47\sqrt{T^{}\text{(MeV)}})}\mathrm{exp}\left[\frac{(TT^{})^2}{0.44T^{}(\text{MeV})}\right],$$
(21)
and we take the differential cross section $`d\sigma _\alpha (E_\nu ,T^{})/dT^{}`$ from .
In the statistical treatment of all these data we perform a $`\chi ^2`$ analysis for the different sets of data, following closely the analysis of Ref. with the updated uncertainties given in Refs. , as discussed in Ref. . We thus define a $`\chi ^2`$ function for the three set of observables $`\chi _{\text{rates}}^2`$, $`\chi _{\text{zenith}}^2`$, and $`\chi _{\text{spectrum}}^2`$ where in both $`\chi _{\text{zenith}}^2`$, and $`\chi _{\text{spectrum}}^2`$ we allow for a free normalization in order to avoid double-counting with the data on the total event rate which is already included in $`\chi _{\text{rates}}^2`$. In the combinations of observables we define the $`\chi ^2`$ of the combination as the sum of the different $`\chi ^2`$’s. In principle such analysis should be taken with a grain of salt as these pieces of information are not fully independent; in fact, they are just different projections of the double differential spectrum of events as a function of time and energy. Thus, in our combination we are neglecting possible correlations between the uncertainties in the energy and time dependence of the event rates.
### B MSW Regions
We present next the results of the allowed regions in the two–parameter space $`\mathrm{\Delta }m^2`$, $`\mathrm{sin}^2(2\theta )`$ for the analysis of the different combination of observables. In building these regions, for a given set of observables and a certain value of the noise parameters $`L_0`$ and $`\xi `$ (in what follows we present the results for $`\xi =4\%`$ which is a reasonable large value for the noise amplitude) we compute for any point in the parameter space of two–neutrino oscillations the expected values of the observables and with those and the corresponding uncertainties we construct the function $`\chi ^2(\mathrm{\Delta }m^2,\mathrm{sin}^2(2\theta );L_0,\xi )_{obs}`$. We find its minimum in the full two-dimensional space of MSW oscillations. The allowed regions for a given CL are then defined as the set of points satisfying the condition:
$$\chi ^2(\mathrm{\Delta }m^2,\mathrm{sin}^2(2\theta );L_0,\xi )_{obs}\chi _{min,obs}^2(L_0,\xi )\mathrm{\Delta }\chi ^2\text{(CL, 2 dof)}$$
(22)
where, for instance $`\mathrm{\Delta }\chi ^2(`$CL, 2 dof)=4.6, 6.1, and 9.2 for CL=90, 95, and 99 % respectively. In Figs. 58 we plot the corresponding allowed regions for $`\xi =4`$ % and for five values of the correlation length $`L_0=`$70, 200, 700, 2000, and 10<sup>4</sup> km. For comparison in the last panel we also show the regions in the absence of random noise.
Figure 5 shows the results of the fit to the observed total rates only. We see in the figure that for any value of the correlation length we always find the three allowed regions, SMA, LMA and LOW as in the standard noiseless MSW analysis although their extend in $`\mathrm{\Delta }m^2`$ and $`\mathrm{sin}^2(2\theta )`$ varies with the value of the correlation length $`L_0`$. The presence of noise modifies the shape and size of the allowed regions through two different although related effects. Due to the modification of the shape of the electron survival probability, the value of the expected rates for a given point on the MSW plane are different once the noise is included what modifies the value of $`\chi ^2`$ for that given point. This leaves also to a shift on the value of $`\chi _{min}^2`$ used to define the regions.
In Table II we show the values of the local $`\chi _{min}^2`$ in the three regions for the different values of the correlation length. First thing we notice is that for the analysis of the rates only in the LMA region the value of $`\chi _{min}^2`$ is always lower in the presence of noise. This can be understood by looking at Fig. 2.b where we see that noise increases the survival probability for $`2\times 10^7\delta 3\times 10^6`$, ie for lower values of the neutrino energies relevant for the Gallium experiments. This increases the rate at the Gallium experiments which is underestimated in the standard LMA solution. For the same reason the LMA become larger and shifted to lower $`\mathrm{\Delta }m^2`$ values and it extends to smaller angles. For larger $`L_0`$, noiseless LMA solution is obtained.
For the SMA solution we find that unless for very short or very long correlation lengths, for which $`\chi _{min}^2`$ is almost the same as in the absence of noise, for 100 km $`L_0`$ few 1000 km, the presence of noise leads to an increase on the value of $`\chi _{min}^2`$, or in other words to a worse description of the data in the SMA region. In Fig. 2.a we see that the presence of noise leads to an increase on the survival probability for $`2\times 10^7\delta 3\times 10^6`$, so there is not enough suppression of the Berilium neutrinos in the relevant $`\mathrm{\Delta }m^2`$ values of the SMA solution. For this reason the rate at Chlorine experiment is increased and the fit worsens as compared to the noiseless case. The worsening is maximal for $`L_0`$ of few thousand km when the effect of noise is maximally resonantly enhanced as discussed in Sec. II C. In this case the LMA gives a better description of the measured rates. Also the SMA region is shifted towards larger values of the mixing angle. Finally we notice that the LOW solution is basically unmodified by the presence of noise as it mainly corresponds to values of $`\delta 10^7`$ which are little affected by the presence of perturbations as seen in Fig. 2.b.
### C Zenith angle Dependence and Super–Kamiokande Spectrum
Figure 6 shows the regions allowed by the fit of both total rates and the Super–Kamiokande zenith angular distribution. Also plotted is the excluded region at 99 % CL from the zenith angular measurement. As seen in the figure the shape of the excluded region is very little dependent of $`L_0`$, as expected. The zenith angular dependence measures the regeneration effect on the neutrino survival probability when the neutrinos travel through the Earth matter and it is expected to be independent on the details of the sun matter modelling. The main effect of the inclusion of the day–night variation data is to cut down the lower part of the LMA region for any value of $`L_0`$. Since for short correlation lengths the LMA region had been shifted towards lower $`\mathrm{\Delta }m^2`$ values, the inclusion of the zenith angle distribution data leads to a reduction of the size of the LMA region for short $`L_0`$. In this way, for instance, for $`L_0=70`$ km the LMA region at 99 % CL extends only in the range $`1.5(\mathrm{\Delta }m^2/10^5)`$ eV$`{}_{}{}^{2}7`$– to be compared with $`1.5(\mathrm{\Delta }m^2/10^5)`$ eV$`{}_{}{}^{2}100`$ in the absence of noise–. The SMA is also reduced in size as compared to the noiseless case. We also find that once the zenith angle information is included the LMA becomes a better solution (lower $`\chi _{min}^2`$) for $`L_0`$=700–2000 km as seen in Table II. For these intermediate $`L_07002000`$ km, when the LMA region obtained from the analysis of the rates extended to smaller angles, some tiny regions in between the LMA and the SMA are still allowed at the 99 % CL.
In Fig. 6 we show the regions allowed by the fit of both total rates and the Super–Kamiokande recoil electron energy spectrum. Also plotted is the excluded region at 99 % CL from the spectrum measurement. The main effect of the inclusion of the spectrum data is to improve the quality of the LMA solution as compared to the SMA. Comparing with Fig. 5 we observe that the regions become larger after the inclusion of the spectrum data. In particular the LMA region extends to larger values of $`\mathrm{\Delta }m^2`$. This behaviour is also observed in the absence of noise. This is mainly due to the flattening of the $`\chi ^2`$ function after the inclusion of the spectrum data and it is independent on the presence of noise. Notice that the larger sensitivity of the mean $`P_{ee}`$ to the presence of noise occurs for low energy neutrinos as discussed in Sec. II. Having a threshold energy, $`T_{th}=5.5MeV`$, Super–Kamiokande is insensitive to such low energy effects.
### D Global
Figure 8 displays the results for the allowed regions from the global analysis of the solar neutrino data including the data on the total event rates, the zenith angular dependence and the recoil electron energy spectrum. We find that, after all the observables are included the three regions remain valid. The shape and size of final LMA solution is very little affected by the presence of noise. The SMA solution is maximally deformed for correlation lengths few 100 km $`L_0`$ few 1000 km. It is for these values also that the SMA gives a worse description of the observables. This behaviour is mainly driven by the effect on the total event rates which is the observable more sensitive to the presence of noise. As discussed above both the zenith angle distribution data and the spectrum data have very little discriminating power on the noise parameters. One also finds that the tiny 99% CL “islands” in between the SMA and LMA are allowed for $`L_07002000`$ km.
Finally we notice that the LOW solution is basically unmodified by the presence of the noise. For the neutrino energies accessible at existing experiments, the LOW region corresponds to values of $`\delta 10^7`$ which are little affected by the presence of noise.
## IV Discussion
In this paper we have studied the effect of density matter fluctuations in the sun on the MSW solutions to the SNP. Assuming no specific mechanism for generation of the fluctuations we have kept the amplitude and correlation length as independent parameters.
Our analysis is performed under no assumption on the relative size of the correlation length of fluctuations as compared to the neutrino oscillation length. To perform such a study we have solved numerically the evolution equation for the neutrino system including the full effect of the random matter density fluctuations of given amplitude and correlation length. Our procedure is to generate a realization of the density profile for given values of the perturbation amplitude and correlation length and then to solve numerically the evolution equation for the neutrino states for that given realization of the density profile and different neutrino production points and finally to average the obtained survival probability over different density realizations (with the same amplitude and correlation length) and neutrino production points. This numerical approach of averaging the solutions $`P_{ee}(r_0)`$ over noise realizations is different from any previous approaches with the averaging of the Schrödinger equation itself before obtaining the solution.
Our results for the survival probabilities are presented in Figs. 24. We find that the effects are larger for small mixing angles. The larger the mixing angle the shorter the correlation length needed to observe an effect. For the SMA the larger effects occur for correlation lengths in the range few 100 km $`L_0`$ few 1000 km. They can be understood as due to a parametric resonance occuring when the phase acquired by the oscillating neutrino state on one fluctuation length $`L_0`$ is a multiple of 2$`\pi `$. This resonance is maximal when this condition is verified close to the MSW resonance. We find that this resonant parametric condition is mainly achieved for low neutrino energies such as the pp-neutrinos seen in GALLEX and SAGE.
Next, in order to establish the possible effect of the presence of noise on the MSW solutions to the SNP we have performed a global analysis of all the existing observables including not only the measured total rates but also the Super–Kamiokande measurement on the time dependence of the event rates during the day and night, as well as the recoil electron energy spectrum. The result of such analysis is presented in Figs. 58 where we plot the allowed regions for MSW neutrino oscillations in the framework of two–neutrino mixing with the Sun density profile generated from the BP98, after including random noise with amplitude $`\xi =4\%`$ and different correlation lengths $`L_0`$ (70, 200, 700, 2000 y 10000 km). The main conclusions are that the total rates are the most sensitive observables to the presence of noise. On the other hand when the many degrees of freedom corresponding to the Super–Kamiokande spectrum are included the dependence of the allowed mixing parameters on the matter noise is smoothed. This is caused by the larger sensitivity of the mean $`P_{ee}`$ to the noise for low energy neutrinos. Due to its higher energy threshold, the Super–Kamiokande experiment is mostly insensitive to these effects. For the same reason one expects that the Borexino experiment would be more suitable to place bounds both on the level of neutrino noise $`\sqrt{<\xi >^2}`$ and on the correlation length $`L_0`$.
###### Acknowledgements.
This work was supported by DGICYT under grants PB98-0693 and PB97-1261, and by the TMR network grant ERBFMRXCT960090 of the European Union. M.C. Gonzalez-Garcia wish to thank the Phenomenology Institute for their kind hospitality during her visit. The work of A.A. Bykov, V.Yu. Popov and V.B. Semikoz was partially supported through RFBR grant 00-02-16271 and for V.B.S. by Iberdrola excellence grant.
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# 1 Introduction
## 1 Introduction
Recently, higher dimensional theories with extra dimensions have revived and have vastly been discussed from various points of view . In many such scenarios, nontrivial backgrounds, such as magnetic flux, vortices, domain walls and branes, turn out to be a key ingredient. It would be of great importance to study physical consequences caused by the nontrivial backgrounds thoroughly.
In this letter, we shall concentrate on a simple situation that “magnetic” flux passes through a circle $`S^1`$. Physically, this system may equivalently be described by the system without flux but with fields obeying twisted boundary conditions for the $`S^1`$-direction. In the following, we shall take the latter point of view for a technical reason. Even though the situation we consider is very simple, physical consequences caused by boundary conditions turn out to be unexpectedly rich, as we will see later. The parameter space of a model on $`M^{D1}S^1`$ is, in general, wider than that of the model on $`M^D`$, and is spanned by twist parameters specifying boundary conditions , in addition to parameters appearing in the action. One of characteristic features of such models is the appearance of a critical radius of $`S^1`$, at which some of symmetries are broken/restored. The spontaneous breakdown of the translational invariance for the $`S^1`$-direction is another characteristic feature .
The paper is organized as follows: In the next section, we discuss a general feature of scalar field theories on $`M^{D1}S^1`$ and allowed boundary conditions. In Section 3, the $`O(N)`$ $`\varphi ^4`$ model on $`M^{D1}S^1`$ with the antiperiodic boundary condition is studied. In Section 4, general twisted boundary conditions are investigated, and the spontaneous symmetry breaking caused by nonvanishing vacuum expectation values is classified. In Section 5, the model is reanalyzed from a $`(D1)`$-dimensional field theory point of view. Section 6 is devoted to conclusions and discussions.
## 2 A General Discussion
In this section, we shall discuss a general feature of scalar field theories compactified on $`M^{D1}S^1`$. Let us consider an action which consists of $`N`$ real scalar fields $`\varphi _i`$ $`(i=1,\mathrm{},N)`$<sup>3</sup><sup>3</sup>3Repeated indices are generally summed, unless otherwise indicated.
$$S=d^{D1}x_0^{2\pi R}𝑑y\left\{\frac{1}{2}_A\varphi _i(x^\nu ,y)^A\varphi _i(x^\nu ,y)V(\varphi )\right\},$$
(1)
where the index $`A`$ runs from $`0`$ to $`D1`$, and $`x^\nu `$ $`(\nu =0,\mathrm{},D2)`$ and $`y`$ are the coordinates on $`M^{D1}`$ and $`S^1`$, respectively. The radius of $`S^1`$ is denoted by $`R`$. Suppose that the action has a symmetry $`G`$. Since $`S^1`$ is multiply-connected, we can impose a twisted boundary condition on $`\varphi _i`$ such as
$$\varphi _i(x^\nu ,y+2\pi R)=U_{ij}\varphi _j(x^\nu ,y).$$
(2)
The matrix $`U`$ must belong to $`G`$, otherwise the action would not be single-valued. If $`U`$ is not proportional to the identity matrix, the symmetry group $`G`$ will be broken to its subgroup $`H`$, which consists of all the elements of $`G`$ commuting with $`U`$. Note that this symmetry breaking caused by the boundary condition is not spontaneous but explicit.
In order to find the vacuum configuration of $`\varphi _i(x^\nu ,y)`$, one might try to minimize the potential $`V(\varphi )`$. This would, however, lead to wrong vacua in the present model . To find the true vacuum configuration, it is important to take account of the kinetic term in addition to the potential term. This is because the translational invariance could be broken and the vacuum configuration might be coordinate-dependent. Thus, the vacuum configuration will be obtained by solving a minimization problem of the following functional:
$$[\varphi ,R]_0^{2\pi R}𝑑y\left\{\frac{1}{2}\left(\frac{d\varphi _i(y)}{dy}\right)^2+V(\varphi )\right\},$$
(3)
where we have assumed that the translational invariance of the uncompactified $`(D1)`$-dimensional Minkowski space-time is unbroken.<sup>4</sup><sup>4</sup>4This is true, at least, at the classical level.
In general, solving the minimization problem may not be an easy task because we must minimize the functional $`[\varphi ,R]`$ with the boundary condition (2). Although we have no general procedure to solve the minimization problem, we can present candidates of the vacuum configuration of $`\varphi _i(y)`$ for some class of twisted boundary conditions. Suppose that $`G`$ is a continuous symmetry and that the matrix $`U`$ in Eq.(2) can be expressed as $`U=e^X`$, where $`X`$ belongs to the algebra of $`G`$. ($`U`$ should continuously be connected to the identity in $`G`$.) Then, a candidate of the vacuum configuration will be given by
$$\overline{\varphi }_i(y)=(e^{\frac{y}{2\pi R}X})_{ij}v_j,$$
(4)
where $`v_i`$ $`(i=1,\mathrm{},N)`$ are constants. Note that $`\overline{\varphi }_i(y)`$ satisfy the desired boundary condition (2). Even if $`U`$ cannot continuously be connected to the identity in $`G`$, we could find a configuration such as Eq.(4) by restricting some of $`v_i`$ to zero. In fact, we will see later that the vacuum configuration can be written into the form (4) in the $`O(N)`$ $`\varphi ^4`$ model (except for $`N=1)`$.
## 3 $`O(N)`$ $`\varphi ^4`$ Model with the Antiperiodic Boundary Condition
We shall now investigate the $`O(N)`$ $`\varphi ^4`$ model whose potential is given by
$$V(\varphi )=\frac{\mu ^2}{2}\varphi _i\varphi _i+\frac{\lambda }{8}\left(\varphi _i\varphi _i\right)^2.$$
(5)
Since the phase structure is trivial for a positive squared mass, we will assume $`\mu ^2>0`$ in the following analysis. The boundary condition for $`\varphi _i(y)`$ is taken to be antiperiodic, i.e.
$$\varphi _i(y+2\pi R)=\varphi _i(y)\text{for }i=1,\mathrm{},N\text{ .}$$
(6)
General twisted boundary conditions will be discussed in the next section. Since $`U=\mathrm{𝟏}`$, the twisted boundary condition (6) does not break the $`O(N)`$ symmetry, and hence the unbroken symmetry $`H`$, which is consistent with the boundary condition, is $`O(N)`$ itself.
Let us first consider the case of even $`N`$. In this case, it may be convenient to introduce the $`N/2`$ complex fields by
$$\mathrm{\Phi }_a(y)\frac{e^{i\frac{y}{2R}}}{\sqrt{2}}\left(\varphi _{2a1}(y)+i\varphi _{2a}(y)\right)\text{for }a=1,\mathrm{},\frac{N}{2}\text{ .}$$
(7)
It should be noticed that $`\mathrm{\Phi }_a(y)`$ obey the periodic boundary condition, i.e.
$$\mathrm{\Phi }_a(y+2\pi R)=+\mathrm{\Phi }_a(y)\text{for }a=1,\mathrm{},\frac{N}{2}.$$
(8)
Inserting Eq.(8) into $`[\varphi ,R]`$, we may write
$$[\varphi ,R]=^{(1)}[\mathrm{\Phi },R]+^{(2)}[\mathrm{\Phi },R],$$
(9)
where
$`^{(1)}[\mathrm{\Phi },R]`$ $``$ $`{\displaystyle _0^{2\pi R}}𝑑y\left\{\left|{\displaystyle \frac{d\mathrm{\Phi }_a}{dy}}\right|^2{\displaystyle \frac{i}{2R}}\left(\mathrm{\Phi }_a^{}{\displaystyle \frac{d\mathrm{\Phi }_a}{dy}}{\displaystyle \frac{d\mathrm{\Phi }_a^{}}{dy}}\mathrm{\Phi }_a\right)\right\},`$ (10)
$`^{(2)}[\mathrm{\Phi },R]`$ $``$ $`{\displaystyle _0^{2\pi R}}𝑑y\left\{\left({\displaystyle \frac{1}{4R^2}}\mu ^2\right)\left|\mathrm{\Phi }_a\right|^2+{\displaystyle \frac{\lambda }{2}}\left(\left|\mathrm{\Phi }_a\right|^2\right)^2\right\}.`$ (11)
Our strategy to find the vacuum configuration, which minimizes the functional (9), is as follows: We shall first look for configurations which minimize each of $`^{(1)}[\mathrm{\Phi },R]`$ and $`^{(2)}[\mathrm{\Phi },R]`$, and then construct configurations which minimize both of them simultaneously.
Let us first look for configurations which minimize $`^{(1)}[\mathrm{\Phi },R]`$. To this end, we may expand $`\mathrm{\Phi }_a(y)`$ in the Fourier-series, according to the boundary condition (8), as
$$\mathrm{\Phi }_a(y)=\underset{n𝐙}{}\phi _a^{(n)}e^{i\frac{n}{R}y}\text{for }a=1,\mathrm{},\frac{N}{2}.$$
(12)
Inserting Eq.(12) into $`^{(1)}[\mathrm{\Phi },R]`$, we find
$$^{(1)}[\mathrm{\Phi },R]=\frac{2\pi }{R}\underset{n𝐙}{}\left[\left(n+\frac{1}{2}\right)^2\left(\frac{1}{2}\right)^2\right]\left|\phi _a^{(n)}\right|^2.$$
(13)
Since $`(n+1/2)^2(1/2)^20`$ for all $`n𝐙`$, $`^{(1)}[\mathrm{\Phi },R]`$ is positive semi-definite. The configuration which gives $`^{(1)}[\mathrm{\Phi },R]=0`$ is found to be of the form $`\mathrm{\Phi }_a(y)=\phi _a^{(0)}+\phi _a^{(1)}e^{i\frac{y}{R}}(a=1,\mathrm{},N/2)`$, where $`\phi _a^{(0)}`$ and $`\phi _a^{(1)}`$ are arbitrary complex constants. Let us next look for configurations which minimize $`^{(2)}[\mathrm{\Phi },R]`$. We find that the configuration which minimizes $`^{(2)}[\mathrm{\Phi },R]`$ is $`\mathrm{\Phi }_a(y)=0`$ for $`R1/(2\mu )`$ and $`\left|\mathrm{\Phi }_a(y)\right|^2=(\mu ^21/(2R)^2)/\lambda `$ for $`R>1/(2\mu )`$. Combining the above two results and performing an appropriate orthogonal $`O(N)`$ transformation, we conclude that in terms of $`\varphi _i`$ the vacuum configuration, which minimizes both of $`^{(1)}[\mathrm{\Phi },R]`$ and $`^{(2)}[\mathrm{\Phi },R]`$ simultaneously, can take to be of the form
$$\varphi _i(x^\nu ,y)=\{\begin{array}{cc}(0,0,\mathrm{},0)\hfill & \text{for }R\frac{1}{2\mu }\hfill \\ (v\mathrm{cos}(\frac{y}{2R}),v\mathrm{sin}(\frac{y}{2R}),0,\mathrm{},0)\hfill & \text{for }R>\frac{1}{2\mu }\text{ ,}\hfill \end{array}$$
(14)
where $`v=\sqrt{2(\mu ^21/(2R)^2)/\lambda }`$. It follows that for $`R1/(2\mu )`$ the $`O(N)`$ symmetry is unbroken, while for $`R>1/(2\mu )`$ the spontaneous symmetry breaking occurs and the $`O(N)`$ symmetry is broken to $`O(N2)`$. It is interesting to contrast this result with that of the $`O(N)`$ $`\varphi ^4`$ model with the periodic boundary condition, for which the $`O(N)`$ symmetry is spontaneously broken to $`O(N1)`$ irrespective of $`R`$.
We now proceed to the case of odd $`N`$. In this case, we cannot apply the same method, as was done above, to find the vacuum configuration because we cannot take a complex basis such as Eq.(7) for odd $`N`$ and because the twist matrix $`U=\mathrm{𝟏}`$ cannot continuously be connected to the identity matrix. Nevertheless, we can show that the problem to find the vacuum configuration for odd $`N`$ reduces to that for even $`N`$ (expect for $`N=1`$). The trick is to add an additional real field $`\varphi _{N+1}(y)`$ satisfying the antiperiodic boundary condition to the action in order to form the $`O(N+1)`$ $`\varphi ^4`$ model. It follows from the previous analysis that the vacuum configuration will be found to be of the form (14) since $`N+1`$ now becomes an even integer. The fact that the configuration space spanned by $`\{\varphi _i(y)`$, $`i=1,\mathrm{},N+1\}`$ contains that by $`\{\varphi _i(y)`$, $`i=1,\mathrm{},N\}`$ implies that the vacuum for odd $`N`$ is also given by Eq.(14), and hence the spontaneous symmetry breaking from $`O(N)`$ to $`O(N2)`$ can occur for $`R>1/(2\mu )`$. The exception is the model with $`N=1`$. In this case, there is no continuous symmetry and the $`O(1)`$ model has only a discrete symmetry of $`G=H=Z_2`$. The $`O(1)`$ $`\varphi ^4`$ model has been investigated in Ref. and the vacuum configuration has been found to be
$$\varphi (x^\nu ,y)=\{\begin{array}{cc}0\hfill & \text{for }R\frac{1}{2\mu }\hfill \\ \frac{2k\mu }{\sqrt{\lambda (1+k^2)}}\mathrm{sn}(\frac{\mu }{\sqrt{1+k^2}}(yy_0),k)\hfill & \text{for }R>\frac{1}{2\mu }\text{ .}\hfill \end{array}$$
(15)
Here, $`\mathrm{sn}(u,k)`$ is the Jacobi elliptic function whose period is $`4K(k)`$, where $`K(k)`$ denotes the complete elliptic function of the first kind. The $`y_0`$ is an integration constant and the parameter $`k`$ $`(0k<1)`$ is determined by the relation $`\pi R\mu =\sqrt{1+k^2}K(k)`$. Thus, the $`Z_2`$ symmetry is unbroken for $`R1/(2\mu )`$, while it is broken spontaneously for $`R>1/(2\mu )`$.
## 4 General Twisted Boundary Conditions
In this section, we shall construct the vacuum configurations of the $`O(N)`$ $`\varphi ^4`$ model on $`M^{D1}S^1`$ for general twisted boundary conditions and clarify the phase structure.
To discuss general boundary conditions, it is convenient to transform the matrix $`U`$ in Eq.(2) by means of a real orthogonal transformation into the normal form. It is known that any matrix $`U`$ belonging to $`O(N)`$ can be transformed, by an orthogonal transformation, into a block diagonal form whose diagonal elements are one of $`1`$, $`1`$ and a two dimensional rotation matrix . In this basis, we may arrange the boundary conditions for $`\varphi _i(y)`$ as follows:
$`\varphi _a^{(\alpha _0)}(y+2\pi R)`$ $`=`$ $`+\varphi _a^{(\alpha _0)}(y)\text{for }a=1,\mathrm{},L_0\text{ ,}`$
$`\left(\begin{array}{c}\varphi _{2b_k1}^{(\alpha _k)}(y+2\pi R)\\ \varphi _{2b_k}^{(\alpha _k)}(y+2\pi R)\end{array}\right)`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{cos}(2\pi \alpha _k)& \mathrm{sin}(2\pi \alpha _k)\\ \mathrm{sin}(2\pi \alpha _k)& \mathrm{cos}(2\pi \alpha _k)\end{array}\right)\left(\begin{array}{c}\varphi _{2b_k1}^{(\alpha _k)}(y)\\ \varphi _{2b_k}^{(\alpha _k)}(y)\end{array}\right)`$
for $`b_k=1,\mathrm{},\frac{L_k}{2}`$ and $`k=1,\mathrm{},M1`$ ,
$`\varphi _c^{(\alpha _M)}(y+2\pi R)`$ $`=`$ $`\varphi _c^{(\alpha _M)}(y)\text{for }c=1,\mathrm{},L_M,`$ (16)
where $`L_0+L_1+\mathrm{}+L_{M1}+L_M=N`$ and $`0=\alpha _0<\alpha _1<\mathrm{}<\alpha _{M1}<\alpha _M=1/2`$. The above boundary conditions explicitly break the $`O(N)`$ symmetry down to
$$H=O(L_0)\times U(\frac{L_1}{2})\times \mathrm{}\times U(\frac{L_{M1}}{2})\times O(L_M)$$
(17)
which is the subgroup of $`O(N)`$ commuting with the twist matrix $`U`$.
Let us first consider the case of $`L_00`$. In this case, $`\varphi _a^{(\alpha _0)}(y)`$ $`(a=1,\mathrm{},L_0)`$ satisfy the periodic boundary condition. Then, it is easy to show that the vacuum configuration can, without loss of generality, be taken into the form
$$\varphi _1^{(\alpha _0)}(x^\nu ,y)=\sqrt{\frac{2}{\lambda }}\mu ,$$
(18)
and other fields vanish. Thus, the symmetry $`H`$ in Eq.(17) is spontaneously broken to <sup>5</sup><sup>5</sup>5For $`L_0=1`$, $`O(L_0=1)`$ means $`Z_2`$ and the $`Z_2`$ symmetry is broken completely.
$$I=O(L_01)\times U(\frac{L_1}{2})\times \mathrm{}\times U(\frac{L_{M1}}{2})\times O(L_M),$$
(19)
irrespective of the value of the radius $`R`$.
Let us next consider the case of $`L_0=0`$ and $`N=\text{even}`$. It is then convenient to introduce the $`N/2`$ complex fields as
$$\mathrm{\Phi }_{b_l}^{(\alpha _l)}(y)\frac{e^{i\frac{\alpha _l}{R}y}}{\sqrt{2}}\left(\varphi _{2b_l1}^{(\alpha _l)}(y)+i\varphi _{2b_l}^{(\alpha _l)}(y)\right)\text{for }b_l=1,\mathrm{},\frac{L_l}{2}\text{ and }l=1,\mathrm{},M.$$
(20)
Inserting Eqs.(20) into $`[\varphi ,R]`$, we may rewrite it into the form
$$[\varphi ,R]=^{(1)}[\mathrm{\Phi },R]+^{(2)}[\mathrm{\Phi },R]+^{(3)}[\mathrm{\Phi },R],$$
(21)
where
$`^{(1)}[\mathrm{\Phi },R]`$ $``$ $`{\displaystyle _0^{2\pi R}}𝑑y\left\{\left|{\displaystyle \frac{d\mathrm{\Phi }_{b_l}^{(\alpha _l)}}{dy}}\right|^2i{\displaystyle \frac{\alpha _l}{R}}\left(\mathrm{\Phi }_{b_l}^{(\alpha _l)}{\displaystyle \frac{d\mathrm{\Phi }_{b_l}^{(\alpha _l)}}{dy}}{\displaystyle \frac{d\mathrm{\Phi }_{b_l}^{(\alpha _l)}}{dy}}\mathrm{\Phi }_{b_l}^{(\alpha _l)}\right)\right\},`$
$`^{(2)}[\mathrm{\Phi },R]`$ $``$ $`{\displaystyle _0^{2\pi R}}𝑑y\left\{\left[\left({\displaystyle \frac{\alpha _1}{R}}\right)^2\mu ^2\right]\left|\mathrm{\Phi }_{b_l}^{(\alpha _l)}\right|^2+{\displaystyle \frac{\lambda }{2}}\left(\left|\mathrm{\Phi }_{b_l}^{(\alpha _l)}\right|^2\right)^2\right\},`$
$`^{(3)}[\mathrm{\Phi },R]`$ $``$ $`{\displaystyle _0^{2\pi R}}𝑑y\left[\left({\displaystyle \frac{\alpha _l}{R}}\right)^2\left({\displaystyle \frac{\alpha _1}{R}}\right)^2\right]\left|\mathrm{\Phi }_{b_l}^{(\alpha _l)}\right|^2.`$ (22)
Since $`(\alpha _1)^2<(\alpha _l)^2`$ for $`l=2,\mathrm{},M`$, it is not difficult to show that in terms of the fields (20) the vacuum configuration which minimizes every $`^{(j)}[\mathrm{\Phi },R]`$ $`(j=1,2,3)`$ simultaneously can, without loss of generality, be taken into the form
$$\mathrm{\Phi }_{b_l}^{(\alpha _l)}(x^\nu ,y)=\{\begin{array}{cc}0\hfill & \text{for }R\frac{\alpha _1}{\mu }\hfill \\ \frac{v}{\sqrt{2}}\delta _{\alpha _l,\alpha _1}\delta _{b_l,1}\hfill & \text{for }R>\frac{\alpha _1}{\mu }\hfill \end{array}$$
(23)
with $`v=\sqrt{2(\mu ^2(\alpha _1/R)^2)/\lambda }`$. It follows that for $`R\alpha _1/\mu `$ the symmetry $`H`$ with $`L_0=0`$ is unbroken, while for $`R>\alpha _1/\mu `$ it is spontaneously broken to <sup>6</sup><sup>6</sup>6 For $`L_1=\mathrm{}=L_{M1}=0`$, the symmetry $`H`$ is $`O(N)`$ and is broken to $`O(N2)`$ for $`R>1/(2\mu )`$, as shown in the previous section.
$$I=U(\frac{L_1}{2}1)\times U(\frac{L_2}{2})\times \mathrm{}\times U(\frac{L_{M1}}{2})\times O(L_M).$$
(24)
Let us finally investigate the case of $`L_0=0`$ and $`N=\text{odd}`$. To find the vacuum configuration, we may perform the trick used in the previous section: We add an additional real field $`\varphi _{N+1}(y)`$ which satisfies the antiperiodic boundary condition to the action. Then, the resulting model may become the $`O(N+1)`$ model, which has been analyzed just above since $`N+1`$ is now even. The result of the $`O(N+1)`$ model will tell us that the vacuum configuration for the $`O(N)`$ model with odd $`N`$ can be taken into the same form as Eq.(23) (except for $`N=1`$).<sup>7</sup><sup>7</sup>7 Since for $`N=1`$ the possible boundary condition is either periodic or antiperiodic, the $`O(1)`$ model has no new phase more than discussed in the previous section. It follows that for $`R\alpha _1/\mu `$ the symmetry $`H`$ with $`L_0=0`$ is unbroken, while for $`R>\alpha _1/\mu `$ the spontaneous symmetry breaking occurs and the symmetry $`H`$ is broken to $`I`$ given in Eq.(24).
## 5 Reanalysis with Kaluza-Klein Modes
In the previous sections, we have succeeded to reveal the phase structure of the twisted $`O(N)\varphi ^4`$ model. In this section, we shall reanalyze the model from a $`(D1)`$-dimensional field theory point of view, and discuss Nambu-Goldstone modes associated with the broken symmetries and also the symmetry breaking of the translational invariance for the $`S^1`$-direction.
To avoid inessential complexities, we shall restrict our considerations to the case of $`L_0,L_M=`$ even. The $`N`$ real fields (16) can then form the $`N/2`$ complex fields which are expanded in the Fourier-series as
$$\frac{1}{\sqrt{2}}\left(\varphi _{2b_l1}^{(\alpha _l)}(x^\nu ,y)+i\varphi _{2b_l}^{(\alpha _l)}(x^\nu ,y)\right)=\underset{n𝐙}{}\phi _{b_l,n}^{(\alpha _l)}(x^\nu )e^{i(\frac{n+\alpha _l}{R})y}$$
(25)
for $`l=0,1,\mathrm{},M`$ and $`b_l=1,2,\mathrm{},L_l/2`$. Inserting Eq.(25) into Eq.(3), we have, up to the quadratic terms with respect to $`\phi _{b_l,n}^{(\alpha _l)}`$,
$$_0[\phi ,R]=2\pi R\underset{l=0}{\overset{M}{}}\underset{b_l=1}{\overset{L_l/2}{}}\underset{n𝐙}{}m_{l,n}^2|\phi _{b_l,n}^{(\alpha _l)}|^2,$$
(26)
where $`m_{l,n}^2`$ are the squared masses of the Kaluza-Klein modes $`\phi _{b_l,n}^{(\alpha _l)}`$ and are given by
$$m_{l,n}^2=\mu ^2+\left(\frac{n+\alpha _l}{R}\right)^2.$$
(27)
The second term in Eq.(27) is the Kaluza-Klein mass, which comes from the “kinetic” term $`\frac{1}{2}(_y\varphi _i(y))^2`$ and which gives a positive contribution to the squared mass term.
For $`L_00`$, the squared mass $`m_{0,0}^2`$ for the modes $`\phi _{b_0,0}^{(\alpha _0)}`$ is always negative irrespective of $`R`$. This observation suggests that $`\phi _{b_0,0}^{(\alpha _0)}`$ acquire non-vanishing vacuum expectation values, so that the $`O(L_0)`$ symmetry is spontaneously broken. This is consistent with the results obtained in the previous section. Taking Eq.(18) into account, we should replace the fields $`\phi _{b_l,n}^{(\alpha _l)}`$ by $`\stackrel{~}{\phi }_{b_l,n}^{(\alpha _l)}+\frac{\mu }{\sqrt{\lambda }}\delta _{l,0}\delta _{b_l,1}\delta _{n,0}`$ and then find that all the squared masses for $`\stackrel{~}{\phi }_{b_l,n}^{(\alpha _l)}`$ become positive semi-definite, as they should be. The $`L_01`$ massless modes, $`\mathrm{Im}\stackrel{~}{\phi }_{1,0}^{(\alpha _0)}`$ and $`\stackrel{~}{\phi }_{b_0,0}^{(\alpha _0)}(b_0=2,3,\mathrm{},L_0/2)`$, are found to appear and turn out to correspond to the Nambu-Goldstone modes associated with the broken generators of $`O(L_0)/O(L_01)`$.
For $`L_0=0`$, all the squared masses in Eq.(27) are positive for $`R<\alpha _1/\mu `$. The $`m_{1,0}^2`$ vanishes at $`R=\alpha _1/\mu `$ and becomes negative for $`R>\alpha _1/\mu `$. This is a signal of the phase transition and is consistent with the results obtained in the previous section. Taking Eq.(23) into account, we should replace the fields $`\phi _{b_l,n}^{(\alpha _l)}`$ by $`\stackrel{~}{\phi }_{b_l,n}^{(\alpha _l)}+\frac{v}{\sqrt{2}}\delta _{l,1}\delta _{b_l,1}\delta _{n,0}`$ for $`R>\alpha _1/\mu `$ and then find that all the squared masses become positive semi-definite, as they should be. The $`L_11`$ massless modes, $`\mathrm{Im}\stackrel{~}{\phi }_{1,0}^{(\alpha _1)}`$ and $`\stackrel{~}{\phi }_{b_1,0}^{(\alpha _1)}(b_1=2,3,\mathrm{},L_1/2)`$, are found to appear and turn out to correspond to the Nambu-Goldstone modes associated with the broken generators of $`U(\frac{L_1}{2})/U(\frac{L_1}{2}1)`$. If $`L_M=N`$, the additional $`N2`$ massless modes, $`\stackrel{~}{\phi }_{b_M,1}^{(\alpha _M)}(b_M=2,3,\mathrm{},N/2)`$, appear and all the massless modes turn out to form the Nambu-Goldstone modes associated with the broken generators of $`O(N)/O(N2)`$.
We shall finally discuss the symmetry breaking of the translational invariance for the $`S^1`$-direction. For $`L_00`$, the vacuum expectation values of the fields are coordinate-independent, so that the translational invariance is unbroken. For $`L_0=0`$, the vacuum expectation values depend on the coordinate $`y`$ for $`R>\alpha _1/\mu `$, so that the translational invariance for the $`S^1`$-direction is spontaneously broken. It may be instructive to point out that the translational invariance for the $`S^1`$-direction can be reinterpreted as a global $`U(1)`$ symmetry, which is in fact possessed by the theory after compactification. To see this, we note that the translations $`yy+ϵR`$ in Eq.(25) can equivalently be realized by the following $`U(1)`$ transformations:
$$U(1):\phi _{b_l,n}^{(\alpha _l)}e^{i(n+\alpha _l)ϵ}\phi _{b_l,n}^{(\alpha _l)}$$
(28)
from which we may assign a $`U(1)`$ charge $`n+\alpha _l`$ to the field $`\phi _{b_l,n}^{(\alpha _l)}`$. Thus, the spontaneous breakdown of the translational invariance for the $`S^1`$-direction may be understood as that of the $`U(1)`$ symmetry. For $`L_00`$, some of $`\phi _{b_0,0}^{(\alpha _0)}`$ acquire non-vanishing vacuum expectation values but have no $`U(1)`$ charges, so that the $`U(1)`$ symmetry is unbroken. For $`L_0=0`$, some of $`\phi _{b_1,0}^{(\alpha _1)}`$ acquire non-vanishing vacuum expectation values for $`R>\alpha _1/\mu `$. Since $`\phi _{b_1,0}^{(\alpha _1)}`$ have the nonzero $`U(1)`$ charge $`\alpha _1`$, the $`U(1)`$ symmetry would be broken for $`R>\alpha _1/\mu `$. However, the following modified $`U(1)^{}`$ symmetry, which is a combination of the $`U(1)`$ symmetry and the $`O(N)`$ symmetry, survives as a symmetry even for $`R>\alpha _1/\mu `$:
$$U(1)^{}:\phi _{b_l,n}^{(\alpha _l)}e^{inϵ}\phi _{b_l,n}^{(\alpha _l)}.$$
(29)
This is because $`\phi _{b_1,0}^{(\alpha _1)}`$ now have zero $`U(1)^{}`$ charge. Hence, no new Nambu-Goldstone modes are produced other than those found before.
## 6 Conclusions and Discussions
We have studied the $`O(N)`$ $`\varphi ^4`$ model compactified on $`M^{D1}S^1`$ with the general twisted boundary conditions. Since $`S^1`$ is multiply-connected, the model can be parametrized by not only the mass and the coupling appearing in the action but also the twist matrix appearing in the boundary condition (2). Thus, the parameter space of the $`O(N)`$ model on $`M^{D1}S^1`$ is much wider than that on $`M^D`$. We have succeeded to reveal the rich phase structure and to classify the patterns of the symmetry breaking/restoration thoroughly.
In this letter, our analysis has been restricted to the classical level, and has not taken quantum corrections into account. When the radius $`R`$ of $`S^1`$ is large, $`R`$-dependent quantum corrections might be small. But when $`R`$ is smaller than the inverse of the mass, the leading correction to the squared mass turns out to be proportional to $`1/R^2`$ for $`D=4`$ and hence could drastically change the phase structure at the classical level. Furthermore, the introduction of gauge fields leads to a new interesting feature: Twisted boundary conditions in the directions of the gauge symmetry can dynamically be determined through the Hosotani mechanism . It would be of great interest to analyze $`R`$-dependent quantum corrections in gauge field theories and the phase structure of symmetries systematically. The work on these subjects will be reported elsewhere.
Acknowledgments
We would like to thank to H. Hatanaka, M. Tachibana and K. Takenaga for valuable discussions.
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# Field Theory of 𝑁 Entangled Polymers
## I Introduction
The discovery of the intimate relationship between the statistical mechanics of long polymer molecules and certain field theories describing critical systems has been crucial in achieving the present understanding of the physics of unentangled polymer chains .
Under certain experimental conditions, polymers form knotted configurations which remain stable in time. In this case, the topological relationships among the molecules become relevant. The principal theoretical tools for investigationg these were set up by S. Edwards and Brereton and Shaw , and formulated by one of the authors (HK) with the help of the topological field theory of Chern and Simons . While the tools are available, a satisfactory theoretical treatment of an ensemble of topologically linked polymers is still missing, in spite of much work .
Definite progress has recently been achieved by the present authors by mapping the statistical mechanics of two fluctuating polymers to a topological Ginzburg-Landau model, and applying Feynman diagram techniques to calculate the average of the square winding number. In the present work we extend the field theoretic formulation of to a system of $`N`$ polymers.
Excluded volume forces, although physically important, have so far been neglected, and will be ignored also in this work.
We start in Section II from the path integral description of polymers in which these are viewed as Brownian trajectories. The topological constraints are imposed using the simplest link invariant; the Gauss link integral for pairs of trajectories. The topological interactions are complicated, and several nontrivial steps are necessary to arrive at a useful field theoretic description of the system. For this purpose we introduce in Section III a set of auxiliary abelian Chern-Simons fields. These allow us to make the theory local and to convert the polymer path integral to a Markoffian form. Actually, there are several possible abelian Chern-Simons field theories which could accomplish this task, differing from each other by the number of fields. However, we prove that these are equivalent after exploiting the field equations and the freedom of performing linear transformations of the fields.
The arbitrariness in choosing the auxiliary topological field theory is used in Section VI to overcome a technical problem which was absent in the two polymer case: the coupling constants of the interactions between Chern-Simons fields and polymer trajectories are related to the parameters for the topological constraints by non-linear algebraic equations. The latter are too complicated to be solved analytically apart from particular cases, in which the parametrization simplifies considerably. A particular choice of the Chern-Simons fields removes this problem.
With these methods we convert the polymer path integral to a Markoffian form and it becomes possible to complete the mapping of the polymer problem into a field theoretical model of topological entanglement. This is achieved in Section VI by exploiting the method of replica .
## II Path Integral Approach to Topological Polymers
Let $`P_1,\mathrm{},P_n`$ be a set of topologically linked polymers of lengths $`L_1,\mathrm{}L_N`$ respectively. In order to keep our treatment as general as possible we consider both, open and closed chains of polymers. Strictly speaking, the entanglement of open chains is not really topological. Over a long time scale, it is always possible to disentangle open polymers, and the topology of the system is not conserved. However, if we restrict ourselves to intermediate time scales short compared to that of complete molecular rearrangements, open chains will behave almost like closed chains. In this approximate sense we shall be able to apply our topological methods also to open chains. Moreover, the study of the difference between the behaviors of open and closed chains will make it possible to gain some insight in the nature of topological interactions . Thus, we shall start our considerations with open polymers $`P_i`$ with $`i=1,\mathrm{},N`$, whose end points lie at $`𝐱^i,𝐲^i`$. For coinciding end points $`𝐱^i=𝐲^i,`$ they become closed polymers running through $`𝐱^i`$.
The relevant quantity to describe the statistical mechanics of the polymers is the configurational probability $`G_{\{m\}}(\stackrel{}{𝐱},\stackrel{}{𝐲};\stackrel{}{L})`$, which measures the probability to find the polymers $`P_i`$ with end points at $`𝐱^i,𝐲^i`$ in a given topological configuration $`\{m\}`$. To compactify the notations, we have collected the set of all end points in $`N`$-dimensional multi-vectors $`\stackrel{}{𝐱}=(\stackrel{}{𝐱}^1,\mathrm{},𝐱^N),\stackrel{}{y}=(𝐲^1,\mathrm{}𝐲^N)`$, with an analogous vector notation for the lenghts of the polymers: $`\stackrel{}{L}=(L_1,\mathrm{},L_N)`$. To distinguish different topological configurations, we shall use the Gauss link invariant:
$`\chi (P_i,P_j)=`$ (2)
$`{\displaystyle \frac{1}{4\pi }}{\displaystyle _0^{L_i}}𝑑s_i{\displaystyle _0^{L_j}}𝑑s_j\dot{𝐱}(s_i)\left[\dot{𝐱}^j(s_j)\times {\displaystyle \frac{𝐱^i(s_i)𝐱^j(s_j)}{|𝐱^i(s_i)𝐱^j(s_j)|^3}}\right],`$
defined for each pair $`P_i,P_j`$ of polymers with $`ij`$.
The integral (2) is well-defined also for open trajectories, but it becomes a real topological invariant only for a pair of closed polymers with $`𝐱^i𝐲^i`$$`𝐱^i𝐲^i`$. In the latter case it counts how many times $`P_i`$ winds up around $`P_j`$.
Following the approach of Edwards , the configurational probability can be expressed as a path integral over all possible configurations $`𝐱^i(s_i)`$ with $`0s_iL_i`$ and periodic boundary conditions $`𝐱^i(0)=𝐲^i`$ and $`𝐱^i(L_i)=𝐱^i`$:
$`G_{\{m\}}(\stackrel{}{𝐱},\stackrel{}{𝐲}_j\stackrel{}{L})`$ $`=`$ $`{\displaystyle _{𝐲^1}^{𝐱^1}}𝒟𝐱^{}(s_1)\mathrm{}{\displaystyle _{𝐲^N}^{𝐱^N}}𝒟𝐱^N(s_N)e^{(𝒜_0+𝒜_{ev})}`$ (3)
$`\times `$ $`{\displaystyle \underset{i=1}{\overset{N1}{}}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{j=2}{j>i}}{\overset{N}{}}}\delta \left(\chi (P_i,P_j)m_{ij}\right),`$ (4)
where $`𝒜_0`$ is the euclidean action of a random-chain
$`𝒜_0={\displaystyle \frac{3}{2a}}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle _0^{L_i}}\dot{𝐱}_{}^{i}{}_{}{}^{2}𝑑s_i,`$ (5)
and
$`𝒜_v={\displaystyle \frac{1}{2a^2}}{\displaystyle \underset{i,j=1}{\overset{N}{}}}{\displaystyle _0^{L_i}}𝑑s_i{\displaystyle _0^{L_i}}𝑑s_jv_{ij}^0\delta ^{(3)}\left(𝐱^i(s_i)𝐱^j(s_j)\right)`$ (6)
(7)
the steric repulsion between the chain elements. The parameter $`a`$ denotes the length of the chain elements and $`v_{ij}^0`$ represents an $`N\times N`$ matrix of coupling constants with the dimension of a volume.
The $`\delta `$-functions in the integrand of Eq. (4) enforce the topological constraints that the pairs of chains $`P_i,P_j`$ wind around each other a number of times $`m_{ij}`$. Since we are describing open polymers up to this point, these numbers are continuous. Only for closed polymers will they become integer numbers. Then the Dirac $`\delta `$-functions in Eq. (4) have to be replaced by Kronecker $`\delta `$’s.
In the following, it will be convenient to introduce an auxiliary probability $`G_{\{\lambda \}}(\stackrel{}{𝐱},\stackrel{}{𝐲};\stackrel{}{L})`$, from which the original $`G_{\{m\}}(\stackrel{}{𝐱},\stackrel{}{𝐲};\stackrel{}{L})`$ is obtained by a Fourier transformation with respect to the topological numbers:
$`G_{\{m\}}(\stackrel{}{𝐱},\stackrel{}{𝐲};\stackrel{}{L})={\displaystyle _{\mathrm{}}^+\mathrm{}}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{j=2}{j>i}}{\overset{N}{}}}{\displaystyle \frac{d\lambda _{ij}}{2\pi }}e^{i\lambda _{ij}m_{ij}}G_{\{\lambda \}}(\stackrel{}{𝐱},\stackrel{}{𝐲};\stackrel{}{L}).`$ (8)
(9)
The auxiliary probability has the advantage that its path integral representation
$`G_{\{\lambda \}}(\stackrel{}{𝐱},\stackrel{}{𝐲};\stackrel{}{L})`$ $`=`$ $`{\displaystyle _{𝐲^1}^{𝐱^1}}𝒟𝐱^1(s_1)\mathrm{}{\displaystyle _{𝐲^N}^{𝐱^N}}𝒟𝐱^N(s_N)`$ (10)
$`\times `$ $`e^{(𝒜_0+𝒜_v+𝒜_{\mathrm{top}})},`$ (11)
accounts for the topological constraints among the polymers by a source like term:
$`𝒜_{\mathrm{top}}=i{\displaystyle \underset{i=1}{\overset{N1}{}}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{j=2}{j>i}}{\overset{N}{}}}\chi (P_i,P_j)\lambda _{ij}`$ (12)
Note that if in a formulation for closed polymers, where $`m_{ij}`$ are integer numbers and the Dirac $`\delta `$-functions in Eq. (4) are Kronecker symbols, the Fourier variable $`\lambda _{ij}`$ would be cyclic with a range $`\lambda _{ij}(0,2\pi )`$. Then we would use angular variables $`\phi _{ij}`$ rather than $`\lambda _{ij}`$, and with Eq. (9) as
$`G_{\{m\}}(\stackrel{}{𝐱};\stackrel{}{L})`$ $`=`$ $`{\displaystyle _0^{2\pi }}{\displaystyle \underset{i=1}{\overset{N1}{}}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{j=2}{j>i}}{\overset{N}{}}}{\displaystyle \frac{d\phi _{ij}}{2\pi }}e^{i\phi _{ij}m_{ij}}G_{\{\lambda \}}(\stackrel{}{𝐱};\stackrel{}{L}).`$ (13)
where
$`G_{\{m\}}(\stackrel{}{𝐱};\stackrel{}{L})=\underset{\stackrel{}{𝐱}\stackrel{}{𝐲}}{lim}G_{\{m\}}(\stackrel{}{𝐱},\stackrel{}{𝐲};\stackrel{}{L}).`$ (14)
Returning to (11), we see that it has the form of a path integral over the trajectories of a system of $`N`$ particles performing a random walk. Thus we may exploit the duality between particles and fields, which is valid in statistical mechanics as well as quantum mechanics, and express the path integral in terms of fields. Although the general techniques are well-known , the specific task here is complicated by the presence of the topological term (12), which is non-Markoffian and rather complicated. We shall solve this problem by introducing auxiliary fields, which allow us to rewrite the right-hand side of (11) in a more tractable form.
## III Auxiliary fields and the decoupling of trajectories
Before coming to this we first reformulate the excluded volume interaction $`𝒜_v`$ in Eq. (4) is a standard way, since its non-Markoffian character. The trajectories are coupled to each other by the two-body potential in Eq. (7). The interactions can be disentangled by introducing $`N`$ real scalar fields $`\varphi _1\mathrm{}\varphi _N`$ with the euclidean action
$`𝒜_\varphi ={\displaystyle \frac{a^2}{2}}{\displaystyle \underset{i,j=1}{\overset{N}{}}}{\displaystyle d^3x\varphi _i(𝐱)[(v^0)^1]_{ij}\varphi _j(𝐱)}.`$ (15)
It is obviously possible to rewrite the exponential of $`𝒜_v`$ in (4) with the help of the following identity :
$`e^{𝒜_v}=`$ (17)
$`{\displaystyle 𝒟\varphi _i\mathrm{}𝒟\varphi _Ne^{𝒜_\varphi }\underset{i=1}{\overset{N}{}}\times \mathrm{exp}\left[i_0^{L_i}𝑑s_i\varphi _i\left(𝐱^i(s_i)\right)\right]}.`$
On the right-hand side, each trajectory $`𝐱^i(s_i)`$ interacts only with a single field $`\varphi _i(𝐱)`$, so that the polymers move in individual random fields. Thus, the contribution of the excluded volume forces is converted to a Markoffian form.
## IV Topological Interactions
Let us now turn to the topological interactions, where we search for auxiliary fields to simplify $`𝒜_{\mathrm{top}}`$ in the path integral (11). Such auxiliary fields are provided by Chern-Simons theories . Let $`A_\mu ^\alpha `$ with $`\alpha =1,\mathrm{},N^{}`$ be a set of $`N^{}`$ Chern-Simons fields with euclidean spatial indices $`\mu =1,2,3`$. They will allow us to write the identity
$`e^{𝒜_{\mathrm{top}}}={\displaystyle \frac{𝒟𝐀e^{i𝒜_{\mathrm{CS}}+i\mathrm{\Sigma }_{i=1}^N\mathrm{\Sigma }_{\beta =1}^N^{}{\scriptscriptstyle d^3xh_{i\beta }𝐉^i𝐀^\beta }}}{𝒟𝐀e^{i𝒜_{\mathrm{CS}}}}}`$ (18)
(19)
where $`𝒜_{\mathrm{cs}}`$ is a the Chern-Simons action
$`𝒜_{\mathrm{CS}}={\displaystyle \underset{\alpha ,\beta =1}{\overset{N^{}}{}}}{\displaystyle d^3x𝐀^\alpha \left(\times 𝐀^\beta \right)g_{\alpha \beta }}`$ (20)
and $`𝐉^i`$ are currents
$`𝐉^i(𝐱)={\displaystyle _0^{L_i}}𝑑s_i\dot{𝐱}^i(s_i)\delta ^{(3)}\left(𝐱𝐱^i(s_i)\right).`$ (21)
The measure of functional integration $`𝒟𝐀`$ is short for the product $`_{\alpha =1}^N^{}𝒟𝐀^\alpha `$. We assume $`g_{\alpha \beta }`$ to be a suitable $`N^{}\times N^{}`$ -symmetric matrix to be specified later, which possesses an inverse $`(g^1)^{\alpha \beta }`$, while $`h_{i\alpha }`$ is a $`N\times N^{}`$ matrix.
The right-hand side of Eq. (19) has the described Markoffian form. To calculate it explicitly, we quantize the Chern-Simons fields $`A_\mu ^\alpha `$ in the Lorentz gauge, where they are completely transverse, and have the correlation functions
$`G_{\mu \nu }^{\alpha \beta }(𝐱,𝐲)`$ $``$ $`A_\mu ^\alpha (𝐱)A_\nu ^\beta (𝐲)`$ (22)
$`=`$ $`{\displaystyle \frac{(g^1)^{\alpha \beta }}{4\pi }}\epsilon _{p\nu \rho }{\displaystyle \frac{(xy)^\rho }{|𝐱𝐲|^3}}.`$ (23)
After some calculations we find
$`{\displaystyle \frac{𝒟𝐀e^{i𝒜_{\mathrm{CS}}+i\mathrm{\Sigma }_{i=1}^N\mathrm{\Sigma }_{\alpha =1}^N^{}{\scriptscriptstyle d^3xh_{i\alpha }𝐉^i𝐀^\alpha }}}{𝒟𝐀e^{i𝒜_{\mathrm{CS}}}}}`$ (24)
$`=\mathrm{exp}\left[{\displaystyle \frac{i}{4}}{\displaystyle \underset{i,j=1}{\overset{N}{}}}h_{i\alpha }(g^1)^{\alpha \beta }h_{j\beta }\chi (P_i,P_j)\right].`$ (25)
The right-hand side is equal to that of Eq. (19) with the topological action (12) if we satisfy the equation
$`h_{i\alpha }(g^1)^{\alpha \beta }h_{j\beta }=2\lambda _{ij}.`$ (26)
Since the matrix $`\lambda _{ij}`$ has vanishing diagonal elements, this condition automatically eliminates the appearance of non-topological terms of the form $`\chi (P_i,P_i)`$ with $`i=1,\mathrm{},N`$.
At this point, it is convenient to introduce the linear combinations of the Chern-Simons fields
$`𝐂^i={\displaystyle \underset{\alpha =1}{\overset{N^{}}{}}}h_{i\alpha }𝐀^\alpha ,`$ (27)
such that Eq. (19) becomes
$`e^{𝒜_{\mathrm{top}}}={\displaystyle \frac{𝒟𝐀e^{i𝒜_{\mathrm{CS}}+i_{i=1}^N{\scriptscriptstyle d^3x𝐉^i𝐂^i}}}{𝒟𝐀e^{i𝒜_{\mathrm{CS}}}}}`$ (28)
The absence of the non-topological contributions $`\chi (P_i,P_i)`$ in (25) is ensured by requiring purely off-diagonal correlation with
$`C_\mu ^i(𝐱)C_\nu ^i(𝐲)=0,i=1,\mathrm{},N.`$ (29)
The linking of a polymer $`P_i`$ with a different polymer $`P_j`$ is described by the off-diagonal correlation functions. the related propagator should be different from zero:
$`C_\mu ^i(𝐱)C_\nu ^j(𝐲)0ij`$ (30)
This condition allows us to conclude that in the most general situation in which all elements $`\lambda _{ij}`$ are non-vanishing, the number $`N^{}`$ of Chern-Simons fields should be at least equal to $`N`$.
Indeed let us assume the contrary: $`N^{}<N`$. In this case we can always rearrange the indices such that the first $`N^{}`$ fields $`𝐂^\sigma (𝐱)`$ with $`\sigma =1,\mathrm{},N^{}`$ are independent. The remaining fields
$`\stackrel{~}{C}_\mu ^\tau (𝐱)C_\mu ^{N^{}+\tau }(𝐱),\tau =1,\mathrm{},NN^{},`$ (31)
should then be linear combinations of the first $`N^{}`$ fields.
$`\stackrel{~}{C}_\mu ^\tau (𝐱)={\displaystyle \underset{\sigma =1}{\overset{N^{}}{}}}s_\sigma ^\tau C_\mu ^\sigma (𝐱)`$ (32)
From the property (29), we find the correlation functions
$`\stackrel{~}{C}_\mu ^\tau (𝐱)\stackrel{~}{C}_\nu ^\tau (𝐲)={\displaystyle \underset{\sigma ,\sigma ^{}=1}{\overset{N^{}}{}}}s_\sigma ^\tau s_\sigma ^{}^\tau C_\mu ^\sigma (𝐱)C_\nu ^\sigma ^{}(𝐲)=0`$ (33)
By hypothesis, however, we have
$`C_\mu ^\sigma (𝐱)C_\nu ^\sigma ^{}(𝐲)\lambda _{\sigma \sigma ^{}}\epsilon _{\mu \nu \rho }{\displaystyle \frac{(xy)^\rho }{|𝐱𝐲|^3}}0.`$ (34)
As a consequence, the most general solution of Eq. (33) is:
$`s_{\overline{\sigma }}^\tau 0,s_\sigma ^\tau =0\text{if}\sigma \overline{\sigma },`$ (35)
where $`\overline{\sigma }`$ is a fixed integer with $`1\overline{\sigma }N^{}`$. Therefore, the fields $`\stackrel{~}{𝐂}^\tau (𝐱)`$ and $`𝐂^{\overline{\sigma }}(𝐱)`$ coincide apart from an irrelevant factor $`s_{\overline{\sigma }}^\tau `$. In this way, if $`N^{}<N`$, we obtain for $`N^{}<N`$ the following contradiction:
$`0`$ $`=`$ $`C_\mu ^{\overline{\sigma }}(𝐱)C_\nu ^{\overline{\sigma }}(y)=s_{\overline{\sigma }}^\tau C_\mu ^{\overline{\sigma }}(𝐱)C_\nu ^{\overline{\sigma }}(𝐲)`$ (36)
$`=`$ $`\stackrel{~}{C}_\mu ^\tau (𝐱)C_\nu ^{\overline{\sigma }}(𝐲)=C_\mu ^{N^{}+\tau }(𝐱)C_\nu ^{\stackrel{~}{\sigma }}(𝐲)`$ (37)
$``$ $`\lambda _{N+\tau \overline{\sigma }}\epsilon _{\mu \nu \rho }{\displaystyle \frac{(xy)^\rho }{|𝐱𝐲|^3}}0`$ (38)
In the opposite case of $`N^{}>N`$, it is always possible to reduce the number of Chern-Simons fields to $`N`$. To show this we consider the action
$`𝒜_{\mathrm{CS}}^J=𝒜_{\mathrm{CS}}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle d^3xJ^i𝐂^i}`$ (39)
appearing in Eq. (28). In the above action we perform a change of variables in which a number $`N`$ of fields $`𝐀^\alpha `$ is expressed as a linear combination of the fields $`𝐂^i`$’s and of the remaining $`𝐀^\alpha `$’s. Without loss of generality, we may suppose (27) to be invertible, with the solutions
$`𝐀^i={\displaystyle \underset{j=1}{\overset{N}{}}}(h^1)^{ij}\left(𝐂^j{\displaystyle \underset{\alpha =N+1}{\overset{N^{}}{}}}h_{j\alpha }𝐀^\alpha \right).`$ (40)
Substituting this into (39), we find
$`𝒜_{\mathrm{CS}}^J={\displaystyle }d^3x\{C^{i\mu }J_\mu ^i+\epsilon ^{\mu \nu \rho }[{\displaystyle \underset{i,j=1}{\overset{N}{}}}M_{ij}C_\mu ^i_\nu C_\rho ^j`$ (41)
$`2{\displaystyle \underset{\alpha =N+1}{\overset{N^{}}{}}}{\displaystyle \underset{i=1}{\overset{N}{}}}N_{\alpha i}A_\mu ^\alpha _\nu C_\rho ^i+{\displaystyle \underset{\alpha ,\beta =N+1}{\overset{N^{}}{}}}O_{\alpha \beta }A_\mu ^\alpha _\nu A_\rho ^\beta ]\}`$ (42)
where the constant coefficients $`M_{ij},N_{\alpha i}`$ and $`O_{\alpha \beta }`$ are functions of the matrix elements $`g^{\alpha \beta }`$ and $`h_{i\alpha }`$. The mixed terms in $`𝒜_{\mathrm{CS}}^J`$, which are proportional to $`N_{\alpha i}`$, are eliminated by introducing the new field variables
$`A^{}{}_{\mu }{}^{\alpha }=A_\mu ^\alpha {\displaystyle \underset{\beta =N+1}{\overset{N^{}}{}}}{\displaystyle \underset{i=1}{\overset{N}{}}}(O^i)^{\alpha \beta }N_{\beta i}C^i.`$ (44)
It is easy to see that the redundant fields $`A^{}^\alpha `$ with $`\alpha =N+1,\mathrm{},N^{}`$ can now be integrated out is Eq. (28), so that we arrive at a Chern-Simons field theory with $`N`$ fields $`𝐂^1,\mathrm{},𝐂^N`$.
This result could be expected from the fact that only the fields $`𝐂^i`$ have sources in the action (39), while the remaining fields are free, and may be eliminated via the equations of motion. As a consequence, the freedom in choosing the number of Chern-Simons fields is only apparent, because any Chern-Simons field theory with $`N^{}>N`$ substituting Eq. (28) is equivalent to a Chern-Simons field theory with $`N`$ fields only. In principle, there is still some arbitrariness in the choice of the matrix elements $`g^{\alpha \beta }`$ and $`h_{i\alpha }`$ once $`N^{}`$ has been fixed, since their values are only constrained by (26). However, this arbitrariness reflects merely the possibility of performing linear transformations of the $`A^\beta `$’s, and it is thus irrelevant.
To conclude this section, we compute the denominator in the right hand side of (28). For this purpose we note that since $`g_{\alpha \beta }`$ is a $`N^{}\times N^{}`$ symmetric matrix, it can always be expressed as follows:
$`g_{\alpha \beta }={\displaystyle \underset{\beta =1}{\overset{N^{}}{}}}\eta _{\alpha \gamma }\eta _{\beta \beta }`$ (45)
where $`\eta _{\alpha \beta }`$ is again a $`N^{}\times N^{}`$ symmetric matrix. Performing in Eq. (20) the substitution:
$`𝐚_\alpha ={\displaystyle \underset{\beta =1}{\overset{N^{}}{}}}\eta _{\alpha \beta }𝐀^\beta `$ (46)
the Chern-Simons action becomes
$`𝒜_{\mathrm{CS}}={\displaystyle \underset{\alpha =1}{\overset{N^{}}{}}}{\displaystyle d^3x\epsilon ^{\mu \nu \rho }a_\mu ^\alpha _\nu a_\rho ^\alpha }`$ (47)
i.e., the dependence on $`g_{\alpha \beta }`$ disappears. Thus:
$`{\displaystyle 𝒟𝐀e^{i𝒜_{\mathrm{CS}}(𝐀^\alpha )}}=\left[det(g_{\alpha \beta })\right]^{1/2}c`$ (48)
where $`c=𝒟𝐚e^{i𝒜_{\mathrm{CS}}(𝐚^\alpha )}`$ is an irrelevant constant factor.
## V Field Theory
In the previous section we have seen that the path integral over the polymer trajectories can be converted to a Markoffian form via auxiliary fields. In rewriting the topological interactions there is some freedom in choosing the auxiliary fields by varying their number. Also the parameters $`g_{\alpha \beta }`$ and $`h_{\alpha i}`$ are not completely fixed by the system of equations (26).
On the other hand, it has been shown that all abelian Chern-Simons field theories for which the relevant identity (28) is satisfied are equivalent after exploiting the equations of motion and performing linear transformations of the vector fields.
The simplest way to ensure (28) is to choose
$`N^{}=Ng_{ij}={\displaystyle \frac{\kappa \lambda _{ij}}{4\pi }}h_{ij}={\displaystyle \frac{\kappa ^{1/2}}{4\pi }}\lambda _{ij}`$ (49)
for $`i,j=1,\mathrm{},N`$. This choice, however, makes the Fourier transformation (9) from the auxiliary probability $`G_{\{\lambda \}}(\stackrel{}{𝐱},\stackrel{}{𝐲};\stackrel{}{L})`$ to the original configurational probability $`G_{\{m\}}(\stackrel{}{𝐱},\stackrel{}{𝐲};\stackrel{}{L})`$ too difficult for an analytic treatment. As a matter of fact, starting from Eq. (49), we find with the help of (48):
$`e^{A_{\mathrm{top}}}=\left[(det\lambda _{ij})^{1/2}c\right]^1{\displaystyle 𝒟𝐀}`$ (50)
$`\times \mathrm{exp}\{{\displaystyle }d^3x{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{j=1}{ji}}{\overset{N}{}}}\lambda _{ij}[{\displaystyle \frac{\kappa }{4\pi }}\epsilon ^{\mu \nu \rho }A_\mu ^i_\nu A_\rho ^j`$ (51)
$`+{\displaystyle \frac{\kappa ^{1/2}}{4\pi }}{\displaystyle _0^{L_i}}d𝐱^i(s_i)𝐀^j(𝐱^i(s_i))]\}.`$ (52)
The denominator in the right hand side contains the term $`\left[det(\lambda _{ij})\right]^{1/2}`$, which complicates the integration over the parameters $`\lambda _{ij}`$ in (9).
Moreover, in the ansatz (49) the requirement that the matrix $`g_{ij}`$ should be invertible cannot be guaranteed for all possible matrices $`\lambda _{ij}`$.
The situation does not improve if we choose
$`N^{}=Ng_{ij}={\displaystyle \frac{\kappa }{4\pi }}\delta _{ij},h_{ij}={\displaystyle \frac{\kappa ^{1/2}}{4\pi }}\eta _{ij}.`$ (53)
Here the elements $`\eta _{ij}`$ have a complicated dependence on the variables $`\lambda _{ij}`$, via the algebraic equations (26).
To solve these difficulties, we exploit the freedom of enlarging the number of topological vector fields. The simplest Chern-Simons field theory for our purpose contains $`N^{}=2(N1)`$ fields $`𝐀^1,\mathrm{}𝐀^{N1}`$ and $`𝐁^1,\mathrm{},𝐁^{N1}`$. The action $`𝒜_{\mathrm{CS}}`$ is given by:
$`𝒜_{\mathrm{CS}}=\kappa {\displaystyle \underset{i=1}{\overset{N1}{}}}\epsilon ^{\mu \nu \rho }{\displaystyle d^3xA_\mu ^i_\nu B_\rho ^i}`$ (54)
Equation (28) becomes now
$`e^{𝒜_{\mathrm{top}}}=c^1{\displaystyle 𝒟𝐀𝒟𝐁e^{i𝒜_{\mathrm{CS}}}\mathrm{exp}\left\{\underset{i=1}{\overset{N}{}}d^3x𝐂^i𝐉^i\right\}},`$ (55)
(56)
where the currents $`J_\mu ^i(𝐱)`$ have been already defined in Eq. (21) and the fields $`𝐂^i`$ of Eq. (27) have the following explicit expressions:
$`𝐂^1=𝐁^1,𝐂^N=\kappa {\displaystyle \underset{i=1}{\overset{N1}{}}}\lambda _{Ni}𝐀^i`$ (57)
$`𝐂^i=\kappa {\displaystyle \underset{j=i}{\overset{i1}{}}}\lambda _{ij}𝐀^j+𝐁^i,i=2,\mathrm{},N1.`$ (58)
The factor $`c^1`$ in Eq. (56) is an irrelevant constant independent of $`\lambda _{ij}`$.
Using Eq. (17) and (56) in the expression of the auxiliary probability (11) we obtain the path integral representation
$`G_{\{\lambda \}}(\stackrel{}{𝐱},\stackrel{}{𝐲};\stackrel{}{L})={\displaystyle \underset{i=1}{\overset{N}{}}}G(𝐱^i,𝐲^i;L_i|\varphi _i,𝐂^i),`$ (59)
where
$`G(𝐱^i,𝐲^i;L_i|\varphi _i,𝐂^i)={\displaystyle _{𝐲^i}^{𝐱^i}}𝒟𝐱^i(s_i)`$ (60)
$`\times e^{\left[_0^{L_i}𝑑s_i_{\varphi _i}(𝐱^i(s_i))+_0^{L_i}𝑑s_i\dot{𝐱}^i(s_i)𝐂^i(𝐱^i(s_i))\right]},`$ (61)
and
$`_{\varphi _i}(𝐱^i(s_i))={\displaystyle \frac{3}{2a}}𝐱^i{}_{}{}^{}{}_{}{}^{2}(s_i)+i\varphi _i(x^i(s_i))`$ (62)
In Eq. (59), the expression in the expectation symbol $``$ must be averaged with respect to the fields $`\varphi _i`$ with $`i=1,\mathrm{},N`$ and the Chern-Simons fields $`𝐀^j,𝐁^j`$ with $`j=1,\mathrm{},N1`$.
The path integral in Eq. (61) describes a Markoffian random walk of a particle immersed in an electromagnetic field $`(𝐂^i,i\varphi _i)`$. In analogy with the evolution kernel of a particle in quantum mechanics, $`G(𝐱^i,𝐲^i,L_i|\varphi _i,𝐂^i)`$ satisfies a Schrödinger-like equation
$`[{\displaystyle \frac{}{L_i}}{\displaystyle \frac{a}{6}}𝐃_{𝐢}^{}{}_{}{}^{2}+i\varphi _i]G(𝐱^i,𝐲^i;L_i|\varphi _i,𝐂^i)=\delta (L_i)\delta (𝐱^i𝐲^i)`$ (63)
(64)
where
$`𝐃_i=\mathbf{}+i𝐂^i.`$ (65)
It is now convenient to consider its Laplace-transformed in the length parameter $`L_i`$:
$`G_{\{\lambda \}}(\stackrel{}{𝐱},\stackrel{}{𝐲};\stackrel{}{\mu })`$ (67)
$`={\displaystyle _0^+\mathrm{}}𝑑L_1\mathrm{}𝑑L_Ne^{_{i=1}^N\mu _iL_i}G_{\{\lambda \}}(\stackrel{}{𝐱},\stackrel{}{𝐲};\stackrel{}{L}).`$
The parameters $`\mu _i`$ in Boltzmann-like factors control the growth of the polymers. Applying the Laplace transformations to both sides of Eq. (59), we find
$`G_{\{\lambda \}}(\stackrel{}{𝐱},\stackrel{}{𝐲};\stackrel{}{\mu })={\displaystyle \underset{i=1}{\overset{N}{}}}G(𝐱^i,𝐲^i;\mu ^i|\varphi _i,𝐂^i),`$ (68)
where $`G(𝐱^i,𝐲^i;\mu _i|\varphi _i,𝐂^i)`$ is the Laplace transformed correlation function (61), obeying the stationary equation
$`\left[\mu _i{\displaystyle \frac{a}{6}}𝐃_i^2+i\varphi _i\right]G(𝐱^i,𝐲^i;\mu _i|\varphi _i,𝐂^i)=\delta (𝐱^i𝐲^i).`$ (69)
Equation (69) can be solved in terms of fluctuating polymer fields $`\psi _i^{},\psi _i`$ as
$`G(𝐱^i,𝐲^i;\mu _i|\varphi _i,𝐂^i)={\displaystyle \frac{1}{Z_i}}{\displaystyle 𝒟\psi _i𝒟\psi _i^{}}`$ (70)
$`\times \psi _i(𝐱^i)\psi _i^{}(𝐲^i)e^{𝒜_{\mathrm{pol}}[\psi _i^{},\psi _i]}`$ (71)
where $`𝒜_{\mathrm{pol}}`$ is the polymer action
$`𝒜_{\mathrm{pol}}[\psi _i^{},\psi _i]={\displaystyle d^3x\left[\frac{a}{6}|𝐃^i\psi _i|^2+(\mu _i+i\varphi _i)|\psi _i|^2\right]}`$ (72)
and $`Z_i`$ is the associated partition function
$`Z_i={\displaystyle 𝒟\psi _i𝒟\psi _i^{}e^{𝒜_{\mathrm{pol}}[\psi _i^{},\psi _i]}}`$ (73)
The integrations over the auxiliary fields $`\varphi _i`$ is complicated by the presence of the factor $`Z_i^1`$ in Eq. (71), which makes them non-gaussian. This problem can be solved exploiting the method or replicas.
## VI Replica Formulation
For each pair of fields $`\psi _i,\psi _i^{}`$ we introduce a set of $`n_i`$ replica field $`\psi _i^{a_i},\psi _i^{a_i}`$ with $`a_i=1,\mathrm{},n_i.`$ It is convenient to group the replica fields in $`n_i`$tuplets
$$\mathrm{\Psi }_i=(\psi _i^1,\mathrm{},\psi _i^{n_i}),\mathrm{\Psi }_i^{}=(\psi _i^1,\mathrm{},\psi _i^{n_i}).$$
With these fields the correlation function (71) may be rewritten as follows:
$`G(𝐱^i,𝐲^i;\mu _i|\varphi _i,𝐂^i)`$ (74)
$`=\underset{n_i0}{lim}{\displaystyle 𝒟\mathrm{\Psi }_i𝒟\mathrm{\Psi }_i^{}\psi _i^{\overline{a}_i}(𝐱^i)\psi _i^{\overline{a}_i}(𝐲^i)e^{𝒜_{\mathrm{rep}}[\mathrm{\Psi }_i,\mathrm{\Psi }_i^{}]}}`$ (75)
where we have set
$`{\displaystyle 𝒟\mathrm{\Psi }_i𝒟\mathrm{\Psi }_i^{}}={\displaystyle \underset{a_i=1}{\overset{n_i}{}}𝒟\psi _i^{a_i}𝒟\psi _i^{a_i}}`$ (76)
and defined the replica field action
$`𝒜_{\mathrm{rep}}[\mathrm{\Psi }_i,\mathrm{\Psi }_i^{}]`$ $`=`$ $`{\displaystyle \underset{a_i=1}{\overset{n_i}{}}}\left[{\displaystyle \frac{a}{6}}|𝐃_i\psi _i^{a_i}|^2+(\mu _i+i\varphi _i)|\psi _i^{a_i}|^2\right]`$ (77)
$``$ $`{\displaystyle \frac{a}{6}}|𝐃_i\psi _i|^2+(\mu _i+i\varphi _i)|\mathrm{\Psi }_i|^2`$ (78)
The index $`\overline{a}_i`$ in (75) is a fixed replica index chosen arbitrarily in the range $`1a_in_i`$.
The limit of zero replica number in (75) is performed by an analytic extrapolation. The path integral on the right hand side is calculated for integer values of $`n_i`$ and the result is then extrapolated analytically to the point $`n_i=0`$.
Combining everything, the auxiliary probability (68) has the functional integral representation
$`G_{\{\lambda \}}(\stackrel{}{𝐱},\stackrel{}{𝐲};\stackrel{}{\mu })=\underset{n_i,\mathrm{},n_N0}{lim}{\displaystyle 𝒟𝐀𝒟𝐁\left[\underset{i=1}{\overset{N}{}}𝒟\mathrm{\Psi }_i𝒟\mathrm{\Psi }_i^{}𝒟\varphi _i\right]}`$ (79)
$`\times \left[{\displaystyle \underset{i=1}{\overset{N}{}}}\psi _i^{\overline{a}_i}(𝐱^i)\psi _i^{\overline{a}_i}(𝐲^i)\right]e^{i𝒜_{\mathrm{CS}}𝒜_\varphi \mathrm{\Sigma }_{i=1}^N𝒜_{\mathrm{rep}}[\mathrm{\Psi }_i,\mathrm{\Psi }_i]}`$ (80)
Note that in this expression the integration over the auxiliary fields $`\varphi _i`$ has become Gaussian. After a suitable rescaling of the fields $`\psi _i^{a_i},\psi _i^{a_i}`$, we give the final result in the following form:
$`G_{\{\lambda \}}(\stackrel{}{𝐱},\stackrel{}{𝐲};\stackrel{}{\mu })=\underset{n_1,\mathrm{}n_N0}{lim}{\displaystyle }𝒟𝐀𝒟𝐁[{\displaystyle \underset{i=1}{\overset{N}{}}}𝒟\mathrm{\Psi }_i𝒟\mathrm{\Psi }_i^{}`$ (81)
$`\times \psi _i^{\overline{a}_i}(𝐱)^i\psi _i^{\overline{a}_i}(𝐲^i)]e^{𝒜_{\mathrm{tot}}}`$ (82)
where
$`𝒜_{\mathrm{tot}}=i𝒜_{\mathrm{CS}}+{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle }d^3x[\mathrm{\Psi }_i^{}(𝐃_i^2+m_i^2)\mathrm{\Psi }_i`$ (83)
$`+{\displaystyle \underset{i,j=1}{\overset{N}{}}}{\displaystyle \frac{2M^2v_{ij}^0}{a^2}}|\mathrm{\Psi }_i|^2|\mathrm{\Psi }_j|^2`$ (84)
and
$`m_i^2=2M\mu _i`$ (85)
With respect to Eq. (80), the fields $`\psi _i^{a_i},\psi _i^{a_i}`$ have been rescaled by a factor $`\sqrt{2/M}`$, so that they acquire the canonical dimension $`[\psi _i^{a_i}]=[\psi _i^{a_i}]=1/2`$ of usual scalar fields.
Following , we have introduced a mass parameter $`M`$ and a parameter playing a similar role as the Planck constant in quantum mechanical path integrals $`\mathrm{}=\mu _a/3`$. The value of $`M`$ has been fixed with respect to the step length $`a`$ by requiring the condition $`\mathrm{}=1`$.
In this way, the action (84) becomes that of a standard field theory with unit Planck constant. Moreover, $`𝒜_{\mathrm{tot}}`$ is a quadratic form with respect to the parameters $`\lambda _{ij}`$, so that the inverse Fourier transformations leading to the original configurational probability can be performed after a diagrammatic evaluation of the correlation functions (82).
The actual calculations are left to a future publication.
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# A Novel Spin-Statistics Theorem in (2+1)𝑑 Chern-Simons Gravity
## I Introduction
In general relativity, although the metric of spacetime is a dynamical entity determined by Einstein’s field equations, the underlying topology is not a priori determined. On a closer inspection, however, one actually finds that once one imposes that spacetime should possess some physically reasonable geometrical conditions, the presence of non-trivial topology is constrained. Simple examples are the well-known constraints on the spacetime topology in Robertson-Walker models. Also, in classical general relativity, when some standard types of energy conditions are valid, non-trivial spatial topology may lead to singularities in spacetime: Gannon’s theorem (see also ) implies that, in a spacetime satisfying the weak energy condition, if one attempts to develop Cauchy initial data on a spatial 3-manifold <sup>*</sup><sup>*</sup>*More precisely, a partial Cauchy surface regular near infinity-see for the appropriate definitions. with a non-simply connected topology, the corresponding Cauchy development will be geodesically incomplete to the past or to the future. The so-called active topological censorship theorem formulated more recently states that in a globally hyperbolic, asymptotically flat spacetime obeying an averaged null energy condition (ANEC), every causal curve beginning and ending at the boundary at infinity can be homotopically deformed to that boundary. Therefore, an external observer near that boundary would not be able to probe the non-simply connectedness of spacetime. This result has been extended to more general contexts than the asymptotically flat case, such as asymptotically anti-de Sitter spacetimes (see and references therein).
In spite of such results, there is still much room left for investigation of the physical consequences of having a non-trivial spatial topology, especially in quantum theory. On the one hand, even in the classical case one can have non-trivial compact spatial topologies, which evade the conditions of the above cited theorems and also have physical interest, and on the other hand, in quantum theory, the energy conditions to prove these theorems are often violated: for example Wald and Yurtsever show that ANEC is violated by the renormalized stress tensor of free fields in generic curved spacetimes. Indeed, its is the existence of this so-called quantum “exotic” matter that permits the violation of the classical area theorem by evaporating black holes , and the existence of “traversable” wormholes, in spite of the above mentioned theorems (see, e.g., for an extensive account). Moreover, it is widely believed that quantum gravity effects will alter the topology of spacetime at Planck scales (“spacetime foam”). Indeed, some semiclassical calculations indicate that a configuration with the presence of wormholes is energetically favored over the euclidean one .
Topological geons, which are the subject of this paper, are topological structures with some remarkable properties. They were first studied by Friedmann and Sorkin , as “localized excitations of spatial topology”, or “lumps” of non-trivial topology in an otherwise Euclidean spatial background. The idea was to view such entities as particles much in the same way as solitons in a field theory. The presence of geons can give rise to half-integer spin states and fermionic or even fractional statistics, in pure (i.e., without matter) quantum gravity . It is common in the literature refer to such solitonic states as geon states. We follow this usage here.
Geons being soliton-like objects, we can talk about their spin and statistics. In , it was shown that such states could violate the usual spin-statistics theorem, in $`(3+1)d`$ and $`(2+1)d`$, if the spatial topology is assumed not to change in time, or more precisely if the topology of the spacetime $`M`$ is of the form $`M=\mathrm{\Sigma }\times \mathrm{IR}`$. On a spacetime of the form $`M=\mathrm{\Sigma }\times \mathrm{IR}`$, the topology of a spatial slice is well-captured by the geons on $`\mathrm{\Sigma }`$. For example, in the $`(2+1)d`$ context that we are interested in this paper, the topology of an orientable, connected surface $`\mathrm{\Sigma }`$ representing space, with at most one asymptotic region, is completely specified by the number of handles. Each handle corresponds to a geon in this simple context. Accordingly, topology changes are always associated with creation and annihilation of geons. It has been suggested that the standard spin-statistics relation can be recovered if geons can be created and annihilated, in other words, topology change may be required in order to establish the full spin-statistics theorem for geons. In this paper we seek instead a relation between spin and statistics assuming a fixed spatial topology.
To appreciate the importance of having or not having a spin-statistics connection for geons, one must recall that in ordinary quantum field theories in Minkowski spacetime, the particles which arise when we second quantize, for example, have this connection naturally. Now, in a hypothetical quantum theory of gravity, one could think of geons as a “particle”, representing the excitations of the topology itself. It seems therefore natural to ask whether they share this connection with normal particles. We find that in the formalism we develop here a different, weaker version of the spin-statistics connection arises, instead of the normal one.
Before we describe our approach to this situation, we examine more carefully what is meant by spin and statistics. Let us assume that we have a configuration space $`Q`$ describing a pair of identical geons. One such configuration can be visualized as two handles on the plane. Now, the quantization of two geons on the plane is not unique. One has to choose some hermitian vector bundle $`B_k`$ over $`Q`$ whose square-integrable sections (with a suitable measure) serve to define the domains of appropriate observables , and are the “wave functions” in the quantum theory. The index $`k`$ labels inequivalent quantizations. The space of these sections is the quantum Hilbert space $`_k`$ of the two-geon system. Physical operations can be implemented as operators on $`_k`$. If we perform a $`2\pi `$-rotation of one of the geons, described by an operator $`C_{2\pi }`$, then its eigenstate will change by a phase $`e^{i2\pi S}`$, where $`S`$ is the spin. Just like particles in $`(2+1)d`$, geons can carry fractional spin, i.e, $`S`$ can be any real number . Similarly, if we exchange the position of the two geons, the wave function will change by the action of an operator $``$ that we call the statistics operator. The standard spin-statistics relation would tell us that the action of $``$ on a two-geon system should be equivalent to acting with the operator $`C_{2\pi }`$ on one of the geons. Note that there is no a priori reason for this relation to hold since $`C_{2\pi }`$ and $``$ correspond to two independent diffeomorphisms of $`\mathrm{\Sigma }`$. Now one can ask if such a relation is true for each quantization procedure parametrized by $`k`$. The results of shed some light on the problem. The authors show that some quantizations violate the spin-statistics theorem, but leave open the question of which are the ones that do not. Furthermore, as emphasized in , the list of quantum theories derived in is completely based on kinematic considerations. In other words, only the diffeomorphism constraint is imposed, whereas the Hamiltonian constraint, which gives the dynamical features of gravity, is not considered at the quantum level. Imposing the latter would further restrict the states, and in this sense some of the values of $`k`$ may not be dynamically allowed.
In this letter we show that, at least for $`(2+1)d`$ gravity in the first order formalism, there is a generalization of the standard spin-statistics connection relating $``$ and $`C_{2\pi }`$, even for a fixed spatial topology, i.e., for spacetime manifolds of the form $`\mathrm{\Sigma }\times \mathrm{IR}`$. We shall consider $`\mathrm{\Sigma }`$ to be a one-point compactified two-manifold, i.e., we compactify the spatial manifold with one asymptotic region by adding a “point at infinity”. In the quantization scheme given in , one considers the mapping class group $`M_\mathrm{\Sigma }`$ (the group of “large” spatial diffeomorphisms, not connected to the identity of $`Diff(\mathrm{\Sigma })`$) and finds a vector bundle $`B_k`$ for each unitary irreducible representation of $`M_\mathrm{\Sigma }`$. Then, one sees no relation between $``$ and $`C_{2\pi }`$ for a generic $`k`$. The physical significance of this procedure is as follows. Physical states in quantum gravity obey the diffeomorphism constraint, meaning that they are invariant under “small” diffeomorphisms, i.e., the diffeomorphisms connected to the identiy of $`Diff(\mathrm{\Sigma })`$, which are the ones generated by this constraint. The diffeomorphism constraint means that “small”diffeos should be regarded as gauge, but leaves one free to consider the states either as invariant under the “large” diffeos (those not connected to the identity of $`Diff(\mathrm{\Sigma })`$), in which case the “large”diffeomorphisms are also viewed as gauge, or just “covariant”, i.e., transforming by an unitary representation of the mapping class group. In this approach, “large” diffeos are regarded as a symmetry of the theory. We adopt the latter view in this work, the former being a special case of this view.
We will look at $`M_\mathrm{\Sigma }`$ as part of a larger algebra $`𝒜`$ of operators describing the quantum theory of geons. It contains the group algebra of $`M_\mathrm{\Sigma }`$. Let us give an intuitive account of $`𝒜`$. We start by considering the classical (reduced) configuration space $`\stackrel{~}{Q}`$ of $`(2+1)d`$ gravity in the first order formalism which is based on the $`SO(2,1)`$ gauge group. It is well-known that this is the space of flat $`SO(2,1)`$ bundles over the space manifold $`\mathrm{\Sigma }`$. As we will discuss in more detail in the body of the paper, this space admits a natural measure. The wave functions are then taken to be square-integrable functions with respect to this measure. We now describe the algebra $`𝒜`$ used for quantization. In building this algebra, we consider only the minimum needed to investigate the spin-statistics connection. First, we comment on its general structure. Its first component consists of the operators of “position” type on the space $`\stackrel{~}{Q}`$ and corresponds to the commutative algebra $`(\stackrel{~}{Q})`$ of continuous functions of compact support $`f:\stackrel{~}{Q}`$. Next we consider the operators corresponding to the symmetries of the theory. The gauge group $`SO(2,1)`$ acting on $`\stackrel{~}{Q}`$ induces an action on functions. Again, instead of $`SO(2,1)`$, we take its group algebra $`𝒢`$. Finally, we also include the algebra $`𝒰`$ of (suitable) remaining operators acting on $`(\stackrel{~}{Q})`$. In other words, $`𝒜`$ has the structure
$$𝒜=(𝒰𝒢)(\stackrel{~}{Q}),$$
(1)
We then choose the algebra $`𝒰`$ to be the group algebra of $`M_\mathrm{\Sigma }`$. It contains all the operations necessary to investigate the spin-statistics connection.
Another important feature is that the first order formalism naturally takes into account the dynamical constraints. The possible quantizations are given by irreducible $``$-representations $`\mathrm{\Pi }_r`$ of $`𝒜`$, where the index $`r`$ parameterizes inequivalent quantizations. We show that there is a large class of quantizations $`\mathrm{\Pi }_r`$ such that statistics is totally determined by spin according to the formula
$$\mathrm{\Pi }_r()=e^{i(2\pi S\theta [r])}\mathrm{II},$$
(2)
on state vectors of spin $`S`$. Here the extra phase $`\theta [r]`$ is completely fixed by the choice of the representation $`\mathrm{\Pi }_r`$.
The rest of the letter is organized as follows. In Section II we briefly review the first order formalism of general relativity and deduce the classical configuration space and the group actions thereon. We then proceed to the construction of the algebra. The geon algebra can be viewed as an example of a transformation group algebra, first studied by Glimm , and the representation theory of this algebra is known. In Section III we analyze more closely the structure of the algebra and classify the irreducible $``$-representations. We then show how a class of states in these representaions possess a spin-statistics connection, namely those states which are eigenstates of a certain charge operator. These states are then argued to be the true physical states, due to a superselection rule. We end the paper with some final remarks.
## II The connection formalism
In the first order formalism, one takes as fundamental variables a triad $`e^{(3)a}=e_\mu ^{(3)a}dx^\mu `$, possibly degenerate and an $`SO(2,1)`$ connection one-form $`A^{(3)a}=\frac{1}{2}ϵ^{abc}\omega _{\mu bc}^{(3)}dx^\mu `$, where $`\omega _{bc}^{(3)a}`$ is the spin connection In our notation, the superscipt (3) on the upper right denote fields on the three-dimensional spacetime $`M`$, of the form $`\mathrm{\Sigma }\times \mathrm{IR}`$ and fields without superscipt correspond to their pullbacks to $`\mathrm{\Sigma }`$.. The Einstein-Hilbert action takes the form
$$S=_Me^{(3)a}F_a^{(3)}+\text{boundary terms},$$
(1)
where $`F_a^{(3)}=d_MA_a^{(3)}+\frac{1}{2}ϵ_{abc}A^{(3)b}A^{(3)c}`$ is the usual curvature for the connection $`A^{(3)}`$. In our convention, Lorentz spacetime indices are represented by Greek letters, and spatial indices by Latin letters $`i,j=1,2`$. Internal $`SO(2,1)`$ indices are represented by Latin letters $`a,b=0,1,2`$. Boundary terms arise in the cases in which the spatial manifold $`\mathrm{\Sigma }`$ is non-compact, or compact with bondary, and are of course zero for closed $`\mathrm{\Sigma }`$.
Upon variation of the action (1) with respect to $`A^{(3)}`$ and $`e^{(3)}`$, we find the equations of motion
$`F^{(3)a}`$ $`=`$ $`0;`$ (2)
$`D_Me^{(3)a}`$ $`=`$ $`0,`$ (3)
where $`D_M`$ denotes covariant differentiation with respect to the connection $`A^{(3)}`$. Let us consider the equations of motion (2) in coordinates. Since $`M`$ is taken to be of the form $`\mathrm{\Sigma }\times \mathrm{IR}`$, we can use a “space + time” splitting. We then obtain the following set of equations for the spatial components:
$`F_{ij}^a`$ $`=`$ $`0,`$ (4)
$`D_{[i}e_{j]}^a`$ $`=`$ $`0,`$ (5)
which are nothing but the pullback of the equations (2) to $`\mathrm{\Sigma }`$ by the natural inclusion $`\mathrm{\Sigma }\mathrm{\Sigma }\times \mathrm{IR}:x(x,0)`$. The covariant differentiation is now with respect to the pullback $`A`$ of the connection $`A^{(3)}`$. Note that eqs. (4) do not involve time derivatives of the basic fields: they are just constraints on the fields $`e^a`$ and $`A_a`$ on $`\mathrm{\Sigma }`$ at any given time, and initial data are a set of basic fields on $`\mathrm{\Sigma }`$ satisfying these constraints. The remaining equations are the time evolution equations for $`e^a`$ and $`A_a`$. Since we shall not make explicit use of the latter, we omit them here.
$`A_{aj}`$ and $`ϵ^{ij}e_i^a`$, $`i=1,2`$ are canonically conjugate variables defined on $`\mathrm{\Sigma }`$. The pairs $`(e^a,A^a)`$ obeying the constraints span the (reduced) phase space $`𝒫`$ of the theory, which is just the cotangent bundle of the space of $`SO(2,1)`$ connections on $`\mathrm{\Sigma }`$. The canonical symplectic structure is given by the Poisson brackets coming from (1). The only non-vanising ones are:
$$\{A_i^a(x),e_j^b(y)\}=\frac{1}{2}\delta _{ab}ϵ_{ij}\delta ^{(2)}(xy),$$
(6)
where $`x,y\mathrm{\Sigma }`$.
The quantum theory in the “position representation” would be described by wave functionals $`\psi [A]`$. The constraints can be easily imposed before quantization, and one then quantizes only the physical degrees of freedom. When $`\mathrm{\Sigma }`$ is a closed (i.e., compact and boundaryless) 2-surface, the constraints imply that the physical configuration space $`Q`$ is given by the moduli space of flat connections,i.e., the set of equivalence classes of flat connections on $`\mathrm{\Sigma }`$ under gauge transformations. When $`\mathrm{\Sigma }`$ is non-compact, however, one has to specify how fields behave asymptotically. This choice gives rise to boundary terms in (1) , and the physical configuration space is the space of those flat connections which have the appropriate asymptotic behavior, modulo those gauge transformations which preserve this behavior.
The full analysis becomes considerably more complicated in the non-compact case because of the asymptotic considerations involved. To simplify matters we just perform a one-point-compactification of $`\mathrm{\Sigma }`$, by adding a point $`p_{\mathrm{}}`$, the “point at infinity”, since the boundary terms in (1) will play no role here. “Rotations” of geons will be considered to be about this point, and we also fix a frame there. Thus, $`\mathrm{\Sigma }`$ is topologically taken to be a closed suface with a marked point and a frame attached there.
Again, just like in the usual closed case, the configuration space is the space of all flat connections. However, gauge transformations which are not trivial at infinity are not a symmetry of the theory. Therefore, in our case, configurations which differ by a gauge transformation which is not trivial at $`p_{\mathrm{}}`$ should not be viewed as equivalent. The reduced configuration space in this case is theerefore the moduli space space of flat connections modulo gauge transformations which are trivial at $`p_{\mathrm{}}`$.
Note also that we only need regular flat initial data on $`\mathrm{\Sigma }`$ to define the configuration space $`Q`$, and to quantize. We make no assumption as to geodesic completeness, and in particular, the formalism can accommodate geodesically incomplete classical solutions. This is important, because in classical general relativity, Gannon’s theorems imply that singularities must arise due to the multiple connectivity of $`\mathrm{\Sigma }`$, at least when $`\mathrm{\Sigma }`$ is non-compact, under certain mild physical assumptions. Even if the formation of singularities occurs in our case, this seems not to interfere with the quantization procedure, at least formally. On the other hand, precisely because of this independence, it is not clear at this point what are the implications, if any, of such singularities in the quantum theory.
A connection $`A`$ on $`\mathrm{\Sigma }`$ is determined by its holonomies. For each closed curve $`\gamma `$ based at $`p_{\mathrm{}}`$ compute the holonomy $`W([\gamma ])=Pe^{_\gamma A}`$. This quantity is invariant under gauge transformations that are identity at $`p_{\mathrm{}}`$. Since $`A`$ is flat, $`W([\gamma ])`$ is invariant under small deformations of $`\gamma `$ preserving $`p_0`$. In other words, it depends only on the homotopy class $`[\gamma ]`$ of loop $`\gamma `$. In fact, $`W`$ gives a homomorphism $`\pi _1(\mathrm{\Sigma })SO(2,1)`$.
Let $`\stackrel{~}{Q}`$ be the set of all such maps. We recall that $`W([\gamma ])`$ changes to $`gW([\gamma ])g^1,gSO(2,1)`$, under gauge transformations that are not identity (and equal $`g`$) at $`p_{\mathrm{}}`$. For closed surfaces with no marked point, one must make an identification $`WgWg^1`$ to get the moduli space of flat connections. In other words, $`Q=\stackrel{~}{Q}/SO(2,1)`$.
In our case, $`\mathrm{\Sigma }`$ is a two-dimensional surface with a marked point $`p_{\mathrm{}}`$, which is chosen to be our base point. Gauge transformations which are not trivial at $`p_{\mathrm{}}`$, taking a value $`g`$ (say) at $`p_{\mathrm{}}`$, change $`W`$ to $`gWg^1`$ as before, but, as explained, these are no longer equivalent. We call this action of $`SO(2,1)`$ by conjugation the gauge action. It corresponds to a Lorentz transformation of our chosen, fixed frame at $`p_{\mathrm{}}`$. The group $`Diff^{\mathrm{}}(\mathrm{\Sigma })`$ of orientation-preserving spatial diffeomorphisms (diffeos) which are trivial at $`p_{\mathrm{}}`$ (and leave a frame there fixed) acts on the holonomies $`W`$ by changing the curve $`\gamma `$. Its subgroup $`Diff_0^{\mathrm{}}(\mathrm{\Sigma })Diff^{\mathrm{}}(\mathrm{\Sigma })`$, connected to the identity (the group of small diffeos) cannot change the homotopy class of $`\gamma `$. Therefore the formulation is already invariant by small diffeos, and the physical configuration space is $`\stackrel{~}{Q}`$. Large diffeos, on the other hand, act nontrivially on the holonomies. So, we can work with the quotient group $`M_\mathrm{\Sigma }=Diff^{\mathrm{}}(\mathrm{\Sigma })/Diff_0^{\mathrm{}}(\mathrm{\Sigma })`$, known as the mapping class group. In particular, the elements $`C_{2\pi }`$ and $``$ are large diffeos. For the sake of simplicity, we will denote the elements of $`Diff^{\mathrm{}}(\mathrm{\Sigma })`$ and its classes in $`M_\mathrm{\Sigma }`$ by the same letters. An important fact is that elements of $`M_\mathrm{\Sigma }`$ commute with the gauge action.
## III The Geon Algebra
The algebra $`𝒜`$ used for quantization has the structure
$$𝒜=(𝒰𝒢)(\stackrel{~}{Q}),$$
(1)
where $`𝒢`$ is the group algebra of $`SO(2,1)`$ and $`(\stackrel{~}{Q})`$ is the space of complex-valued, continuous functions with compact support on $`\stackrel{~}{Q}`$. We choose the algebra $`𝒰`$ to be the group algebra of $`M_\mathrm{\Sigma }`$. $`𝒜`$ contains all the operations necessary to investigate the spin-statistics connection.
Let us give an explicit presentation of $`𝒜^{(1)}`$, the algebra $`𝒜`$ for a single geon. We choose the generators of $`\pi _1(\mathrm{\Sigma })`$ to be the homotopy classes of the loops $`\gamma _1`$ and $`\gamma _2`$ of Fig.1. Each flat connection provides us with a pair of holonomies $`(a,b)=(W(\gamma _1),W(\gamma _2))`$. Since there are no relations among the generators of $`\pi _1(\mathrm{\Sigma })`$, any pair of values $`(a,b)`$ can occur. Therefore $`\stackrel{~}{Q}`$ is $`SO(2,1)\times SO(2,1)`$.
Instead of working with $`(\stackrel{~}{Q})`$ directly, we work with one of its representations. Note that the Haar measure on $`SO(2,1)`$ induces a measure on $`\stackrel{~}{Q}`$. Using this measure we may define an inner product on $`(\stackrel{~}{Q})`$ in the obvious way. The completion of $`(\stackrel{~}{Q})`$ in this norm is a Hilbert space $`_0`$, which is the space of square-integrable functions (with this measure) on $`\stackrel{~}{Q}`$, carrying what we call the defining representation of $`(\stackrel{~}{Q})`$. A function $`f(\stackrel{~}{Q})`$ acts on $`\phi _0`$ as a multiplication operator:
$$(f\phi )(a,b)=f(a,b)\phi (a,b)$$
(2)
With $`gSO(2,1)`$, let $`\widehat{\delta }_g`$ denote the generators of the group algebra $`𝒢`$. These $`\widehat{\delta }_g`$’s are gauge transformations, and act by conjugating holonomies:
$$(\widehat{\delta }_g\phi )(a,b)=\phi (g^1ag,g^1bg)$$
(3)
The mapping class group of $`\mathrm{\Sigma }`$ has two generators $`A`$ and $`B`$, which correspond to Dehn twists along the loops. Their effect on loops $`\gamma _1`$ and $`\gamma _2`$ is given by
$$\begin{array}{c}(A\phi )(a,b)=\phi (a,ba^1),\hfill \\ (B\phi )(a,b)=\phi (ab^1,b)\hfill \end{array}$$
(4)
The generators of $`𝒜^{(1)}`$ are functions $`f(\stackrel{~}{Q})`$, diffeos $`A,B`$ of the mapping class group and gauge transformations $`\delta _g`$.
The mapping class group includes $`C_{2\pi }`$ . Its action on the defining representation is
$$(C_{2\pi }\phi )(a,b)=\phi (cac^1,cbc^1)$$
(5)
where $`c:=aba^1b^1`$. One can verify that $`C_{2\pi }=(AB^1A)^4`$.
These operators can be encoded in what is called a transformation group algebra . Let $`G`$ be a group with a left-invariant measure acting on a space $`X`$. The transformation group algebra is just the set of continuous functions $`(G\times X)`$, with compact support and with the product
$$(F_1F_2)(g,x)=_GF_1(z,x)F_2(z^1g,z^1x)𝑑z.$$
(6)
Here $`xz^1x`$ is the group action on $`X`$, $`z^1g`$ is the group product of $`z^1`$ and $`g`$, and $`dz`$ is the left-invariant measure on $`G`$. The irreducible representations of a transformation group algebra have been worked out in . In our case, $`X=\stackrel{~}{Q}`$ and $`G=SO(2,1)\times M_\mathrm{\Sigma }`$, where $`G`$ can be made into a topological group by giving $`M_\mathrm{\Sigma }`$ the discrete topology. The measure on $`SO(2,1)`$ is the Haar measure and the measure on $`M_\mathrm{\Sigma }`$ is given by
$`{\displaystyle \underset{mM_\mathrm{\Sigma }}{}}f(m)`$
for any function $`f`$ on $`M_\mathrm{\Sigma }`$ with appropriate convergence properties. The measure on $`G`$ is then the product measure. Finally, $`𝒜^{(1)}=(SO(2,1)\times M_\mathrm{\Sigma }\times \stackrel{~}{Q})`$, where we use the bijection
$$\text{ }\mathrm{C}(G)(X)(G\times X)$$
(7)
by interpreting $`\delta _gf`$ as the distribution
$`\delta _gf:(h,x)`$ $``$ $`\delta _g(h)f(x)`$ (8)
$``$ $`\delta (g,h)f(x)`$ (9)
on $`G\times X`$, $`\delta _g`$ being the $`\delta `$-function supported at $`g`$.
Let $`Y=\stackrel{~}{Q}/G`$ be the set of orbits of $`G`$ in $`\stackrel{~}{Q}`$, one such orbit being $`𝒪_\omega `$. Let us choose one representative $`(a_\omega ,b_\omega )\stackrel{~}{Q}`$ for each orbit $`𝒪_\omega `$, and write $`𝒪_\omega =[(a_\omega ,b_\omega )]`$. We define the stabilizer group $`N_\omega G`$ as the set of elements $`(g,\lambda )`$ of $`G`$ such that $`(g,\lambda )(a_\omega ,b_\omega )=(a_\omega ,b_\omega )`$, where the $`G`$ action has been denoted by a dot. Let $`\alpha `$ be a unitary irreducible representation of $`N_\omega `$ on some Hilbert space $`V_\alpha `$. Now consider the space of square-integrable functions $`\varphi :GV_\alpha `$ such that $`\varphi (hg,\xi \lambda )=\alpha (g^1,\lambda ^1)\varphi (h,\xi )`$ for all $`(g,\lambda )N_\omega `$ and $`(h,\xi )G`$. They are called equivariant functions. The set of these functions can be completed into a Hilbert space $`L^2(G,V_\alpha )`$ . The irreducible unitary $``$-representations $`\mathrm{\Pi }_{(\omega ,\alpha )}`$ of $`(G\times \stackrel{~}{Q})`$ can be realized on the Hilbert spaces $`_{(\omega ,\alpha )}=L^2(G,V_\alpha )`$ and, up to unitary equivalence, labeled by $`r=(\omega ,\alpha )`$. This label is a quantum number characterizing a single geon. The action of the operators $`\widehat{F}=\mathrm{\Pi }_r(F)`$, $`F𝒜^{(1)}`$ on a vector $`\varphi ^r_r`$ is given by
$`(\widehat{F}\varphi ^r)(h,\xi )`$ $`={\displaystyle _{SO(2,1)\times M_\mathrm{\Sigma }}}F((h,\xi )(a_\omega ,b_\omega ),(g,\lambda ))\times `$ (10)
$`\varphi ^r(g^1h,\lambda ^1\xi )dz,`$ (11)
for any $`hSO(2,1)`$ and $`\xi M_\mathrm{\Sigma }`$. We find, in particular, that
$$\begin{array}{c}\left(\widehat{\delta _h^{}}\varphi ^r\right)(h,\xi )=\varphi ^r(h^1h,\xi )\hfill \\ \left(\widehat{A}\varphi ^r\right)(h,\xi )=\varphi ^r(h,A^1\xi )\hfill \\ \left(\widehat{B}\varphi ^r\right)(h,\xi )=\varphi ^r(h,B^1\xi )\hfill \\ \left(\widehat{f}\varphi ^r\right)(h,\xi )=f\left(h\xi \stackrel{~}{q}_y\right)\varphi ^r(h,\xi ).\hfill \end{array}$$
(12)
Now, let $`\mathrm{\Sigma }`$ be an orientable surface of genus two with a marked point $`p_{\mathrm{}}`$. It supports a system of two geons. Their algebra $`𝒜^{(2)}`$ can be presented in the defining representation space $`_0_0`$ of $`𝒜^{(1)}𝒜^{(1)}`$. It is generated by elements of $`𝒜^{(1)}𝒜^{(1)}`$ plus the elements of the mapping class group that mix up the geons, with the proviso that we retain only “diagonal”elements of the form $`\delta _g\delta _g`$ from the gauge transformations. There are only two independent generators of $`M_\mathrm{\Sigma }`$ involving both geons. One of them, the diffeo $``$ that exchanges the position of the geons, has already been discussed in connection with the spin-statistics relation. The other one is the so-called handle slide $`H`$. Unlike the exchange $``$, the handle slide $`H`$ has no analogue for particles. Its existence comes from the fact that a geon is an extended object. As the name indicates, it corresponds to the operation of sliding an end of one of the handles through the other handle.
Our description of a pair of geons should be given by an algebra $`𝒜^{(2)}`$ which also includes $`H`$. But since $`H`$ does not enter directly in the spin-statistics relation, we will not include it in $`𝒜^{(2)}`$.
Although $`𝒜^{(1)}`$ is not a Hopf algebra, there is an element $`R𝒜^{(1)}𝒜^{(1)}`$ that plays the role of an $`R`$-matrix. In other words, we can write $`=\sigma R`$ where $`\sigma :_0_0_0_0`$ is the flip automorphism $`\sigma \left(f_1f_2\right)=f_2f_1`$. The $`R`$-matrix turns out to be
$$R=𝑑a𝑑bP_{(a,b)}\delta _{aba^1b^1}^1,$$
(13)
where $`P_{(a,b)}(\stackrel{~}{q},h,\xi )=\delta (\stackrel{~}{q},(a,b))\delta (h,e)\delta (\xi ,e)`$, the $`\delta `$’s being $`\delta `$-functions. The existence of the $`R`$-matrix is essential to establish the connection between spin and statistics. It relates a diffeo performed on a pair of objects with operators acting on each object individually.
Each geon carries a representation $`_r`$ labeled by quantum numbers $`r=(\omega ,\alpha )`$. However, we only need to consider eigenstates of $`\widehat{C}_{2\pi }:=\mathrm{\Pi }^r(C_{2\pi })`$ with spin $`S`$. Let $`\{\varphi _i^{r,S}\}`$ be a basis for the eigenspace of spin $`S`$ in $`_r`$ for some fixed $`r`$. Two geons are said to be identical if they carry the same quantum numbers $`r`$ and $`S`$. We consider identical geons, fix an element $`(a_\omega ,b_\omega )`$ in the corresponding class $`\omega `$ and denote the net flux $`a_\omega b_\omega a_\omega ^1b_\omega ^1`$ by $`c_\omega `$. Consider the characteristic function $`P_c`$ which at $`(a,b)`$ is 1 if $`aba^1b^1=c`$ and zero otherwise. It is clear that a generic vector $`\varphi _i^{r,S}`$ is not an eigenstate of $`\widehat{P}_c`$. A simple computation shows that $`\varphi _i^{r,S}`$ is an eigenstate of $`\widehat{P}_c`$ if and only if it has support only on points $`(h,\xi )`$ such that $`hc_\omega h^1=c_\omega `$.
The quantum state for two identical geons is a linear combination of vectors of the form $`\varphi _i^{r,S}\varphi _j^{r,S}`$. It is enough to show the spin-statistics connection (2) for such decomposable vectors. We must act with the operator $`\widehat{}=(\mathrm{\Pi }_r\mathrm{\Pi }_r)()`$ on these vectors. By using eq. (10), we easily see that
$$\widehat{P}_{(a,b)}\varphi _i^{r,S}(h,\xi )=\delta ((a,b),(h,\xi )(a_\omega ,b_\omega ))\varphi _i^{r,S}(h,\xi )$$
(14)
for every $`(h,\xi )SO(2,1)\times M_\mathrm{\Sigma }`$. Also,
$$\widehat{\delta }_{c^1}\varphi _j^{r,S}(h,\xi )=\varphi _j^{r,S}(ch,\xi ),$$
(15)
where we have put $`c=aba^1b^1`$. Using (13) and the flip automorphism we conclude that
$`\widehat{}\varphi _i^{r,S}(h_1,\xi _1)\varphi _j^{r,S}(h_2,\xi _2)=`$ (16)
$`=\widehat{\delta }_{h_2c_\omega ^1h_2^1}\varphi _j^{r,S}(h_1,\xi _1)\varphi _i^{r,S}(h_2,\xi _2).`$ (17)
At this point we make the assumption that $`\varphi _{i,j}^{r,S}`$ are eigenstates of the net flux $`\widehat{P}_c`$, explaining its physical meaning later. So we can set $`h_2c_\omega h_2^1=c_\omega `$. But we have
$`\widehat{\delta }_{c_\omega ^1}\varphi _j^{r,S}(h_1,\xi _1)`$ $`=`$ $`e^{i2\pi S}\widehat{\delta }_{c_\omega ^1}\widehat{C}_{2\pi }^{}{}_{}{}^{1}\varphi _j^{r,S}(h_1,\xi _1)=`$ (18)
$`=`$ $`e^{i2\pi S}\varphi _j^{r,S}(c_\omega h_1,C_{2\pi }\xi ).`$ (19)
Note that $`\varphi _j^{r,S}(c_\omega h_1,C_{2\pi }\xi )=\varphi _j^{r,S}(h_1c_\omega ,\xi C_{2\pi })`$ because of the above assumption, and because $`c_\omega `$ commutes with $`h_1`$ and $`C_{2\pi }`$ commutes with every element of $`M_\mathrm{\Sigma }`$. On the other hand, $`(c_\omega ,C_{2\pi })N_\omega `$ and hence we can use the equivariance property of $`\varphi _j^{r,S}`$ to rewrite the r.h.s. of the last equality in (18) as
$`\varphi _j^{r,S}(c_\omega h_1,C_{2\pi }\xi )=\alpha (c_\omega ^1,C_{2\pi }^1)\varphi _j^{r,S}(h_1,\xi ).`$
Now, every $`\delta _g`$ commuting with $`a_\omega `$ and $`b_\omega `$ commutes also with $`c_\omega `$, while $`C_{2\pi }`$ is in the center of $`M_\mathrm{\Sigma }`$. Therefore, $`(c_\omega ,C_{2\pi })`$ is in the center of $`N_\omega `$, and by Schur’s lemma we conclude that $`\widehat{\delta _{c_\omega ^1}}\widehat{C}_{2\pi }^{}{}_{}{}^{1}`$ is equal to a phase, say $`e^{i\theta (r)}`$. Eq. (2) then follows:
$$\widehat{}\varphi _i^{r,S}\varphi _j^{r,S}=e^{i[2\pi S\theta (r)]}\varphi _j^{r,S}\varphi _i^{r,S}.$$
(20)
We were able to establish a connection between spin and statistics for all eigenstates of the net flux $`\widehat{P}_c`$. In other words, a spin-statistics exists for states with a definite net flux. Now why are these states special? The answer is that other vectors in the representation space of $`r`$ are not physically allowed as a consequence of a superselection rule, which we will discuss below. As a consequence, only vectors which are in the eigenspace, say $`_c`$, of $`\widehat{P}_c`$ are to be viewed as pure quantum states. Linear combinations of vectors in different $`_c`$’s are not pure, much in the same way as one cannot have pure states of different charges in QED, for example.
This superselection is actually very natural. First, note that the net flux of a geon commutes with all elements of the algebra except the gauge transformations at $`p_{\mathrm{}}`$. Now, the gauge action cannot be viewed as having a local effect from the standpoint of the geons, their effect being limited to performing a transformation on the frame at infinity. The other operators, like those corresponding to the mapping class group operators are “local”, in the sense that they correspond to operations on the geons themselves, i.e., operations which can be taken to leave the region outside some ball surrounding the geons invariant (no other, stronger notion of locality is possible here, since we have no fixed background metric). This is mathematically reflected in the fact that all elements of the geon algebra other than the gauge transformations (which are “local”in the above sense) themselves commute with the gauge action.
Therefore, given some eigenspace $`_c`$ of a net flux operator $`\widehat{P}_c`$, all operators other than gauge transformations preserve $`_c`$. Only the gauge transformation, say corresponding to an element $`gSO(2,1)`$, takes vectors in $`_c`$ into vectors in $`_{gcg^1}`$. That is, gauge transformations do change the net flux, but this change does not correspond to a physical, local operation in the theory; rather it is merely a relabeling of the fluxes. Once one fixes the frame, and considers only local operations, one concludes that the net flux can be regarded as a charge which commutes with all the local operators, and hence is superselected.
## IV Final Remarks
In this paper, we have shown a relation between the actions of the diffeomorphisms $`\widehat{C}_{2\pi }`$ and $``$ on a class of geon staes in $`(2+1)d`$ quantum gravity. An algebra describing the system was identified and its representations were explained in detail.
Our discussion can be viewed a generalization of previous work , where a spin-statistics relation was derived for geonic states arising in a Yang-Mills theory coupled to a Higgs field in the Higgs phase, where the symmetry is spontaneously broken down to a finite gauge group $`H`$. In we showed the existence of a class of “localized” states in quantum gravity arising indirectly from the Yang-Mills theory which did obey the spin-statistics relation derived here. However, those states form a very restricted class. The present paper greatly expands the scope of the original version to a much larger class of geonic states in quantum gravity.
In our version of the spin-statistics relation, there appears an extra phase $`\theta _r`$ for each representation, and a natural question is what is its meaning. It turns out to be a somewhat involved problem, which we are presently tackling .
Acknowledgments
We thank J.C.A. Barata for useful discussions, and S. Carlip, J. Louko, D. Marolf, J. Samuel. R. Sorkin, and D. Witt for important comments regarding the existenc of geons. The work of A.P. Balachandran was supported by the Department of Energy, U.S.A., under contract number DE-FG02-85ERR40231. The work of E. Batista, I.P. Costa e Silva and P. Teotonio-Sobrinho was supported by FAPESP, CAPES and CNPq respectively.
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# THE MERGING SYSTEM AM 2049-691
## 1 Introduction
The observation of the main characteristics of objects with double nucleus and/or distorted morphology indicative of evident gravitational interactions or mergers, allows to improve the knowledge about the involved physical processes. It makes possible to better understand the connections between interactions and different properties such as the separation and size of the components, spatial velocity distributions, global or local star formation activity, infrared emission, etc.
In an attempt to enlarge the available data of this kind of objects we present here spectroscopic and photometric observations of the little studied system AM 2049-691 (ESO 074-IG 020, IRAS 20494-6913) (Fig. 1). Arp & Madore (arp (1987)) classified it as an interacting double object (Category 2) without a definitive statement about its probable merger condition. Lauberts and Valentijn (lauberts (1989)) considered it is composed of a pair of galaxies of Sb-c and Sb morphological types.
A bridge is observed between two very distinct nuclei whose separation resulted of 14”, corresponding to a distance of 10 kpc (H<sub>0</sub> = 75 km s<sup>-1</sup> Mpc<sup>-1</sup>). We also detected the bridge in H$`\alpha `$ emission. The system shows two tails; the tail that emerges from the SW extreme of the central body (0’.6 x 0’.2) is clearly visible up to $``$ 41 kpc from its center, whereas the other one reaches a distance of 58 kpc. At the tip of this last tidal tail it is seen what seems to be a dwarf galaxy as it was observed out of the debris of merging disk galaxies (Mirabel et al. mirabel (1992)).
We have analyzed the main spectral characteristics of AM 2049-691, determined the principal excitation mechanisms, the physical conditions, abundances, the radial velocity distribution of the most important features, as well as its B magnitude and B - V, V - R, and R - I color indexes.
## 2 Observations and Reductions
### 2.1 Spectra
Spectroscopic observations of AM 2049-691 were carried out on 1999 June 15 and 16 with the REOSC spectrograph coupled to the 2.15 m Ritchey-Chrétien telescope of The Complejo Astronómico El Leoncito, San Juan, Argentina, using a Tektronix 1024 x 1024 pixels CCD. The seeing was between 2” and 3” (FWHM). Two set of spectra were obtained through a slit of $``$ 2” wide and $``$ 3’ long, along P.A. = 30. One of them was taken using a 1200 lines mm<sup>-1</sup> grating over a wavelength range of $``$ 6400-7000 Å, and the other one with a 300 lines mm<sup>-1</sup> grating covering the wavelength range of $``$ 3800-7200 Å. The dispersions were 32 and 129 Åmm<sup>-1</sup> and the resolutions 2.5 and 10 Å, respectively. The spectra were corrected for atmospheric and Galactic extinction (A<sub>B</sub> = 0.15; Burstein & Heiles burstein (1984)), and flux calibrated with stars from the catalogue of Stone and Baldwin (stone (1982)).
### 2.2 CCD Photometry
Broadband B, V, R, and I photometric observations were performed on 1999 July 6, 7, and 8 with the CASLEO telescope and the CCD described above. The scale was 0”.27 pixel<sup>-1</sup>. The seeing during observations was also 2”-3” (FWHM). The obtained data were corrected for atmospheric extinction but they were not corrected for Galactic extinction. The photometric calibration was made using standard stars from Graham (graham (1982)) observed through the same filters.
Data reduction of spectra and images was accomplished using the standard methods in IRAF (developed by NOAO) reduction package. The journal of observations of AM 2049-691 is presented in Table 1.
## 3 Spectroscopy
In the spectra of the seven regions listed in Table 2 (regions 3 and 6 correspond to the centers of the NE and SW nuclear regions, respectively), obtained using the 300 lines mm<sup>-1</sup> grating, the lines measured were \[O II\] $`\lambda `$ 3727, H$`\beta `$, \[O III\] $`\lambda `$ 5007, \[O I\] $`\lambda `$ 6300, H$`\alpha `$, \[N II\] $`\lambda `$$`\lambda `$ 6548, 6584, and \[S II\] $`\lambda `$$`\lambda `$ 6717, 6731. The intensities were derived by fitting Gaussians to their profiles. The intensities of the \[O III\] $`\lambda `$ 4959 lines were constrained to the theoretical ratio \[O III\] $`\lambda `$ 4959 = 1/3 \[O III\] $`\lambda `$ 5007, so this line is not listed in Table 2. The internal reddening correction was applied using the interstellar extinction curves given by Seaton (seaton (1979)), assuming that the optical properties of the dust in AM 2049-691 are similar to those of the dust in the Galaxy. Intrinsic ratios H$`\alpha `$/H$`\beta `$ = 2.85 (Osterbrock osterbrock (1989)) were adopted to derive the values of c, the logarithmic extinction at H$`\beta `$.
For the mentioned seven regions the measured and corrected line intensities F<sub>λ</sub> and I<sub>λ</sub>, relative to H$`\beta `$ = 1.00, as well as the errors which were estimated from the noise level around each line, are listed in Table 2. The values of c and the corrected H$`\beta `$ fluxes are given at the bottom of this table. The spectra of all these regions present strong emission lines in the red zones (Fig. 2); their characteristics are typical of H II regions of low excitation, being the excitation considerably lower in region 3 than in region 6. The principal excitation mechanisms would be photoionization by young massive stars. The internal reddening is quite high, especially in region 3; a general decreasing trend is observed from NE to SW.
### 3.1 Abundances, physical conditions and equivalent widths
The abundance ratios N(O)/N(H) and N(N)/N(H), the electron temperatures T<sub>e</sub> and densities N<sub>e</sub> were obtained for regions 1 to 7. For the N(O)/N(H) abundances, the average values of N(O)/NH(H) derived from the empirical calibrations of Edmunds and Pagel (edmunds (1984)) were adopted. The N(N)/N(H) abundances were derived by making the usual assumptions valid for HII regions. Expressions given by Díaz (diaz (1985)) were used for the involved ionic abundances. The electron temperatures were obtained by searching the required values of T<sub>e</sub> for the adopted N(O)/N(H) abundances; the electron densities were derived from the \[S II\] $`\lambda `$6717/$`\lambda `$6731 ratios (osterbrock (1989)). The results are presented in Table 3. The electron temperatures are rather low, but they are in the range of normal values for H II regions; electron densities are also within that range.
The derived nitrogen and oxygen abundances present two maxima corresponding to the region 3 (NE nucleus) and region 6 (SW nucleus), being the abundances of both elements higher in the NE nucleus than in the SW one. In region 3 the N(O)/N(H) and N(N)/N(H) abundance ratios are 2.0 and 1.2 times of the corresponding solar abundances. In region 6 both N(O)/N(H) and N(N)/N(H) ratios are about 1.1 of the respective solar values. The N(N)/N(O) ratios in the region 3 and towards the NE, are practically coincident with those of the galactic emission regions (Shaver et al. shaver (1983)), indicating the same proportions of the involved elements. Towards the SW these ratios increase, being in region 6 about twice of those of galactic regions; this indicates a comparative overabundance of N with respect to the O, which is reflected in the relatively high \[N II\] $`\lambda `$ 6584/H$`\alpha `$ ratios. If this excess is due to an enhancement of nitrogen abundance after a succession of short bursts (Contini et al. contini (1998)), that overabundance suggests that the SW nucleus has undergone previous star formation bursts. Overabundance of N was also detected in the Seyfert component of the interacting pair of galaxies NGC 5953 (Seyfert nucleus) and NGC 5954 (LINER) (González Delgado & Pérez gonzalez (1996)).
The equivalent widths EW(H$`\alpha `$ \+ \[N II\]) also show two maxima at the two nuclei (Fig. 3), being EW(H$`\alpha `$ \+ \[N II\]) = 67 Å and 48 Å for the NE and SW nuclei respectively. All the obtained values indicate enhanced star formation activity compared with isolated galaxies, especially in the NE nucleus. The equivalent width EW(H$`\alpha `$ \+ \[N II\]) = 58 Å derived from the integrated spectrum of AM 2049-691 reflects there is star formation activity in the whole object, which could be favored with the usually large amounts of gas that spiral galaxies have, and is compatible with a merger of two disk galaxies (Liu & Kennicutt liu (1995)). In this system the star formation activity, presumably induced by the interactions, takes place in both nuclei being more significant in the north-eastern one, as detected in some other pairs (Sekiguchi & Wolstencroft sekiguchi (1992)), and in the whole object. This differs from the results of Joseph et al. (joseph (1984)) who found evidence of this activity in only one member of their observed pairs.
The H$`\alpha `$ equivalent widths determined for the NE and SW nuclei are EW(H$`\alpha `$) = 44 and 25 Å; close ages of $``$ 9 x 10<sup>6</sup> yr are derived for their bursts of star formation according to the standard model for instantaneous bursts with metallicities of about 2 Z and 1 Z (Leitherer & Heckman leitherer (1995)) respectively.
The O and N abundances, N(N)/N(O) ratios, electron temperatures and densities, internal reddenings, and equivalent widths are different in the NE and SW nuclei, reflecting the different evolutions they have undergone.
### 3.2 Radial Velocities
Radial velocities were derived from the spectrum obtained with the 1200 lines mm<sup>-1</sup> grating by measuring the centroids of Gaussian curves fitted at the profiles of the strongest emission lines. The resulting heliocentric radial velocities of NE and SW nuclei are V<sub>NE</sub> = (10977 $`\pm `$ 18) km s<sup>-1</sup> and V<sub>SW</sub> = (11144 $`\pm `$ 13) km s<sup>-1</sup> respectively. The average velocity was adopted as the systemic velocity of AM 2049-691, which referred to the Galactic System of Rest is V<sub>GSR</sub> = (10956 $`\pm `$ 30 ) km s<sup>-1</sup>, and the derived distance (H<sub>0</sub> = 75 km s<sup>-1</sup> Mpc<sup>-1</sup>) results 146 Mpc. AM 2049-691 is included in the survey of Sekiguchi & Wolstencroft (sekiguchi (1992)) who reported both nuclear velocities (with errors larger than ours) without any detailed kinematical analysis; the values obtained here are consistent with theirs.
Slight asymmetries were detected in the nuclear emission lines and they can be fitted by secondary components, 100 km s<sup>-1</sup> blueward at SW nucleus and 150 km/s redward at NE nucleus.
The Na I ((5893 Å) absorption line was also detected in the continuum emission of each nuclei, with velocities V<sub>NE</sub> = (10977 $`\pm `$ 20) km s<sup>-1</sup> and V<sub>SW</sub> = (11140 $`\pm `$ 20) km s<sup>-1</sup>. This absorption line appeared with almost the same equivalent width ($``$ 2 Å) at both nuclei but the FWHM of each line were different. Their measured radial velocity dispersions were $`\sigma _{NE}`$ = (280 $`\pm `$ 20) km s<sup>-1</sup> and $`\sigma _{SW}`$ = (330 $`\pm `$ 20) km s<sup>-1</sup> and the deconvolved central radial velocity dispersion of the stellar systems are $`\sigma _{NE}`$ = (225 $`\pm `$ 20) km s<sup>-1</sup> and $`\sigma _{SW}`$ = (280 $`\pm `$ 20) km s<sup>-1</sup>. As the nuclear dynamics of ellipticals and normal bulges in spirals has been found to be indistinguishable (Kormendy & Illingworth kormendy (1982)) we can consider, at first approximation, the central velocity dispersions as indicative of the relative masses of the original systems. Thus the progenitor galaxy of the SW component would be the most massive one. The stellar radial velocity dispersion allows us to estimate the mass and tidal radius (e.g. Bowers & Deeming bowers (1984)) of each nuclear-bulge component, which turns out to be roughly 5.9 x 10<sup>10</sup> M and 6.1 kpc for NE nucleus; and 9.4 x 10<sup>10</sup> M and 8.4 kpc for SW one.
Following the results of Kormendy & Illingworth (kormendy83 (1983)) about the L $``$$`\sigma `$<sup>n</sup> relation for disk-galaxy bulges, the values presented here are roughly consistent with the progenitor systems being spiral galaxies with M<sub>B</sub> $``$ -21. These values are in accordance with the global photometric properties presented in Section 4.2 and the spectrophotometric results discussed in Section 3.1.
The emission lines velocity distribution (along P.A. = 30) is illustrated in Figure 4a, where the open circles correspond to the two distinct nuclei. The radial velocity curve along the line joining both nuclei shows the presence of two different components separated by a velocity discontinuity of $``$ 100 km s<sup>-1</sup>, and a first glance of the curve suggests that each one is associated with a different nucleus. The NE component has an approximate solid body (SB) behavior at all the measured positions and the SW component appears to have a strong asymmetry in the velocity values respect to the nucleus. However, a close inspection of the spectra showed us that the H$`\alpha `$ emission has a minimum between both nuclei at 1/2 of the distance from NE to SW nucleus. Then the three points after the velocity discontinuity, have photometric continuity with the NE emission complex, as shown in Figure 4b. As the tidal radius of NE system is smaller (r<sub>T</sub>$``$ 6 kpc), this feature could be caused by tidal disruption of the NE gaseous system, part of which could have became gravitationally bounded to the SW body, apparently the most massive one. This peculiar kinematic feature and the strong H$`\alpha `$ emission makes this merging galaxy an ideal target for two-dimensional spectroscopy.
The global appearance of the rotation curve is solid body (SB) like at 70% of the observed positions. SB rotation curves appear more frequently in low luminosity galaxies and in interacting disk galaxies (Keel keel (1996)). In the case of AM 2049-691, the SB appearance would not necessary correspond to a spherical halo mass distribution, since recent numerical simulations have shown that an appropriate combination of perturbation and dust obscuration in the disk can explain the SB appearance of an interacting galaxy rotation curve at a wide range of radii (Díaz et al. diaz00 (2000)).
As a whole system, AM 2049-691 shows a velocity amplitude of $``$ 330 km s<sup>-1</sup> (sin i) <sup>-1</sup> within a diameter of $``$ 23 kpc and the kinematical center is possibly located on the line joining both nuclei. The total keplerian mass inside a radius of 11.5 kpc is $``$ 1.4 x 10<sup>11</sup> M ; as the orientation is unknown and this system is far from relaxation, this is only a very rough estimate, but consistent with the luminosities reported in the next section.
## 4 Photometry
### 4.1 Infrared Data
AM 2049-591 is not a luminous infrared object (L<sub>FIR</sub> less than 10<sup>11</sup> L) but it has an appreciable IR emission. Its far infrared flux, FIR = 1.26 x 10<sup>-11</sup>(2.58 S<sub>60</sub> \+ S<sub>100</sub>) (Londsdale et al. londsdale (1985)) = 7.4 x 10<sup>-11</sup> erg cm<sup>-2</sup> s<sup>-1</sup>, calculated using the appropriated data from the IRAS Point Source Catalog (iras (1988)) leads, adopting a mean distance of 146 Mpc (H<sub>0</sub> = 75 km s<sup>-1</sup> Mpc<sup>-1</sup>), to the IR luminosity L<sub>FIR</sub> = 5 x 10<sup>10</sup> L which is compatible with a merger system. Its comparatively low infrared colors $`\alpha `$(60, 25) = -2.5 and $`\alpha `$(100, 60) = -1.5 indicate it is a nonactive object, which is coherent with the spectroscopic results.
In this system, the interaction seems not to have produced very high IR emission when compared with the IR luminosity of typical mergers (L$`{}_{FIR}{}^{}`$ 5 x 10<sup>11</sup> L).
### 4.2 Magnitudes and Colors
B magnitudes and B - V, V - R, and R - I colors of AM 2049-691 were derived using circular apertures with increasing radii (after removing the stars) centered on a point equidistant from the two nuclei. The photometric useful frames were smaller than the total estimated system size, so asymptotical extrapolations of the obtained values were used to estimate the total magnitudes. The obtained results are B = 14.40, B - V = 0.52, V - R = 0.47, and R - I = 0.59. The uncertainties are $`\pm `$0.02 in B, $`\pm `$0.05 in B - V, and $`\pm `$0.06 in V - R and R - I. The magnitudes B and R obtained here are coherent with those of Lauberts & Valentijn (lauberts (1989)), considering as the total magnitude of the whole system that derived from the sum of their individual values for the NE and SW components. The integrated total color B - V would correspond to a Sc-Scd galaxy (Roberts & Haynes roberts (1994)) and indicates that the integrated population is about F7 type.
Indicative magnitudes and colors of both nuclei were derived. Diaphragms with radii of 5” were used for both nuclei. The B magnitudes and B - V colors corresponding to the northeastern and southwestern nuclei are 16.28 and 16.66, and 0.75 and 0.87 respectively. After correcting these values for internal absorption by adopting A<sub>λ</sub> = E<sub>B-V</sub>.X(x) and the extinction curves given by Seaton (seaton (1979)), the B - V colors for the NE and SW nuclei became - 0.25 and 0.19, corresponding respectively to average integrated populations of about B1 and A6 types. These values clearly indicate that the two nuclei are star forming regions, being the observed star formation activity more intense in the NE nucleus than in the other one, as found from the spectroscopic data.
## 5 Summary and Conclusions
We performed CCD spectroscopic and broadband B, V, R, and I photometric observations of AM 2049-691 that is a pair of comparably sized interacting galaxies of morphological types Sb-c and Sb with a separation comparable to their sizes. From the derived information the principal results are:
The spectral characteristics of the all studied regions are typical of H II regions of low excitation; their dominant excitation mechanisms would be the photoionization by young massive stars. The internal reddening is quite high, especially in the northeastern nucleus, and reveals an inhomogeneous obscuration.
All the derived equivalent widths of the H$`\alpha `$ \+ \[N II\] lines indicate enhanced star formation activity compared with isolated galaxies, being this activity more intense in the NE nucleus. The equivalent width corresponding to the integrated spectrum suggests starburst activity in the whole object, and is compatible with a merger of two disk galaxies. For both nuclei the derived starburst ages is about 9 x 10<sup>6</sup> yr suggesting that the merger process triggered the present star formation bursts in both progenitor galaxies at the same time.
The N(N)/N(O) ratio suggests in NE nucleus the same proportion of oxygen and nitrogen as in galactic emission regions; at the SW nucleus this ratio is about twice of those values, indicating there a comparative overabundance of N with respect to O, which is reflected in the relatively high \[N II\] $`\lambda `$6584/H$`\alpha `$ ratios. AM 2049-691 is a merger where overabundance of nitrogen is detected in one of the nuclei, which has the most evolved population and would be the most massive one.
AM 2049-691 is not a very luminous infrared system but has an appreciable IR emission: L<sub>FIR</sub> = 5 x 10<sup>10</sup> L; its comparatively low far infrared colors $`\alpha `$(60,25) and $`\alpha `$(100,60) indicate it is a nonactive object, which is consistent with the derived spectroscopic data.
The integrated total color B - V corresponds to a Sc-Scd galaxy and its average integrated population would be about F7 type. Indicative B - V colors of the nuclei, after correcting for internal extinction, suggest they are regions of star formation activity, specially the NE nucleus as found from the spectroscopic observations.
The central radial velocity dispersions at the nuclei indicate that the most massive galaxy was the progenitor of the SW component. The observed radial velocity curve shows the presence of two components (each one associated to a different nucleus) that undoubtedly correspond to the merging galaxies, what is confirmed by the distinct spectrophotometric and photometric properties shown by the structures associated to each nucleus.
AM 2049-691 is an ongoing merger with a spatially extended star formation activity. It has a high level of disruption and interpenetration and shows a double set of morphological, spectrophotometric and kinematical subsystems.
###### Acknowledgements.
We acknowledge the colaboration of M. Campos at the observing run.
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# Possible Triplet Electron Pairing and an Anisotropic Spin Susceptibility in Organic Superconductors (TMTSF)2X
## Abstract
We argue that (TMTSF)<sub>2</sub>PF<sub>6</sub> compound under pressure is likely a triplet superconductor with a vector order parameter $`𝐝(𝐤)(d_a(𝐤)0,d_c(𝐤)=\mathrm{?},d_b^{}(𝐤)=0)`$; $`|d_a(𝐤)|>|d_c(𝐤)|`$. It corresponds to an anisotropic spin susceptibility at $`T=0`$: $`\chi _b^{}=\chi _0`$, $`\chi _a\chi _0`$, where $`\chi _0`$ is its value in a metallic phase. \[The spin quantization axis, $`𝐳`$, is parallel to a so-called $`𝐛^{}`$-axis\]. We show that the suggested order parameter explains why the upper critical field along the $`𝐛^{}`$-axis exceeds all paramagnetic limiting fields, including that for a nonuniform superconducting state, whereas the upper critical field along the $`𝐚`$-axis ($`𝐚𝐛^{}`$) is limited by the Pauli paramagnetic effects \[I. J. Lee, M. J. Naughton, G. M. Danner and P. M. Chaikin, Phys. Rev. Lett. $`\mathrm{𝟕𝟖}`$, 3555 (1997)\]. The triplet order parameter is in agreement with the recent Knight shift measurements by I. J. Lee et al. as well as with the early results on a destruction of superconductivity by nonmagnetic impurities and on the absence of the Hebel-Slichter peak in the NMR relaxation rate.
PACS numbers: 74.70.Kn, 74.20.-z, 74.60.Ec
Quasi-one-dimensional (Q1D) organic compounds
(TMTSF)<sub>2</sub>X (X = PF<sub>6</sub>, ClO<sub>4</sub>, etc.) have been intensively investigated since the discovery of superconductivity<sup>1,2</sup> in the first organic superconductor (TMTSF)<sub>2</sub>PF<sub>6</sub>. From the beginning, it was clear that their properties are unusual. It was found<sup>3-8</sup> that superconductivity in (TMTSF)<sub>2</sub>X (X = PF<sub>6</sub>, ClO<sub>4</sub>) is destroyed by nonmagnetic impurities. This was interpreted in terms of a possible triplet pairing of electrons<sup>9</sup>. Another unusual feature, the absence of the Hebel-Slichter peak in the $`1/T_1`$ NMR data in (TMTSF)<sub>2</sub>X (X = PF<sub>6</sub>, ClO<sub>4</sub>)<sup>10,11,12</sup>, was prescribed<sup>13</sup> to the existence of zeros of a superconducting order parameter on the Q1D Fermi surfaces (FS). As was stressed<sup>13</sup>, the early experiments<sup>3-8,10,11</sup> provided information only about an orbital part of the order parameter and could not distinguish between some triplet and singlet pairings<sup>2,13</sup>.
To reveal triplet superconductivity, experimental tests which probe a spin part of an order parameter are essential. Among them, are: a surviving of triplet superconductivity in Q1D case<sup>14-17</sup> at magnetic fields higher than both the upper orbital critical field and the Clogston paramagnetic limit<sup>18</sup>, observation of spin-wave exitations<sup>15</sup>, the Knight shift measurements<sup>12</sup> and some others. Nowadays, interest in a possible triplet pairing has been renewed due to remarkable measurements of the upper critical fields (which are sensitive to a spin part of the order parameter) in (TMTSF)<sub>2</sub>ClO<sub>4</sub> and in (TMTSF)<sub>2</sub>PF<sub>6</sub> at $`P6kbar`$ by Naughton, Lee, Chaikin and Danner<sup>19-21</sup> and due to the theoretical analysis<sup>16</sup> of these experiments. The experimental fields along $`𝐛^{}`$-axis (which are 3 times bigger<sup>20,21</sup> than the Clogston paramagnetic limit) were shown<sup>16</sup> to be even bigger than the paramagnetic limit<sup>16,22</sup> for the Larkin-Ovchinnikov-Fulde-Ferrell (LOFF) phase<sup>23</sup>. Therefore, measurements<sup>19-21</sup> were interpreted<sup>16,19-21</sup> in term of triplet superconductivity. Recently, Lee et al.<sup>12</sup> have found no change of the Knight shift for $`𝐇𝐛^{}`$ in a superconducting phase of (TMTSF)<sub>2</sub>PF<sub>6</sub> at $`P6kbar`$. This is consistent with the results<sup>16,19-21</sup> and strongly supports the triplet scenario<sup>9,16,19-21</sup> of superconductivity.
The goals of our paper are: 1) To calculate the paramagnetic limited field along $`𝐛^{}`$-axis, $`H_p^b^{}`$, for the LOFF phase in a Q1D superconductor, taking account of both the paramagnetic<sup>16</sup> and orbital destructive effects against superconductivity. (We show that the calculated value of $`H_p^b^{}`$ is 4-5 times less than the experimental fields<sup>20,21</sup> in (TMTSF)<sub>2</sub>PF<sub>6</sub>); 2) To demonstrate that the value of $`H^b^{}`$ becomes consistent with<sup>20,21</sup> if we switch off the paramagnetic effects. (These indicate that an electron spin susceptibility along $`𝐛^{}`$-axis, $`\chi _b^{}`$, at $`T=0`$ is equal to its value in a metallic state, $`\chi _0`$, which is a distinct feature of triplet superconductivity<sup>24,27</sup>); 3) To stress that the experimental critical fields<sup>20,21</sup> along the conducting chains (i.e., along $`𝐚`$-axis), $`H_p^a`$, are strongly paramagnetically limited and thus the corresponding electron spin susceptibility $`\chi _a\chi _0`$ at $`T=0`$; 4) To show that the above described properties are naturally explained within the framework of a triplet superconductivity scenario with the following vector order parameter frozen into the crystalline lattice (i.e., the case of strong spin-orbit coupling<sup>27</sup>):
$$𝐝(𝐤)=(d_a(𝐤)0,d_c(𝐤)=\mathrm{?},d_b^{}(𝐤)=0);|d_a(𝐤)|>|d_c(𝐤)|$$
(1)
corresponding to the BCS-pair’s wave function
$$\mathrm{\Psi }(𝐤)=[d_a(𝐤)+id_c(𝐤)]|>+[d_a(𝐤)+id_c(𝐤)]|>$$
(2)
and to the anisotropic spin susceptibility at $`T=0`$:
$$\chi _b^{}=\chi _0,\chi _a\chi _0,$$
(3)
where $`|>`$ $`(|>)`$ stands for a spin-up (spin-down) electron with respect to the quantization axis $`𝐳𝐛^{}`$ \[$`𝐚(𝐱)𝐛^{}(𝐳)𝐜^{}(𝐲)`$\], the momentum $`𝐤`$ defines the position on the FS. \[We stress that $`𝐛^{}`$ is the easy axis for a spin direction in a spin-density-wave (SDW) phase of (TMTSF)<sub>2</sub>PF<sub>6</sub>. Thus, one may expect that the order parameter (1) is the most stable since it corresponds to the BCS pairs (2) only with $`S_b^{}S_z=\pm 1`$\]. At the end of the paper, we discuss some consequences of a group theory classification of the possible triplet phases, including the most probable orbital part of the order parameter and a possibility to break the time reversal symmetry.
Q1D electron spectrum corresponds to two open sheets of the FS<sup>1,2</sup>:
$$ϵ^\pm (𝐩)=\pm v_F(p_ap_F)2t_b\mathrm{cos}(p_bb^{})2t_c\mathrm{cos}(p_cc^{}),$$
(4)
where $`+()`$ stands for the right (left) sheet of the FS; $`v_F=t_aa/\sqrt{2}`$ and $`p_F`$ are the Fermi velocity and Fermi momentum, respectively; $`t_a1600K`$, $`t_b200K`$ and $`t_c5K`$; ($`\mathrm{}=1`$).
Singlet (S = 0) and triplet (S = 1) phases are characterized by the following wave functions of the BCS pair’s<sup>27</sup>:
$$\psi _s(𝐤,𝐫)=(|>|>)\psi (𝐤,𝐫),S=0;$$
(5)
$`\psi _t(𝐤,𝐫)=`$ $`|>[d_x(𝐤,𝐫)+id_y(𝐤,𝐫)]+(|>+|>)`$ (7)
$`d_z(𝐤,𝐫)+|>[d_x(𝐤,𝐫)+id_y(𝐤,𝐫)],S=1.`$
\[In Eqs. (5,6), $`S`$ is the total spin of the BCS pair, $`𝐫`$ is its coordinate of a center of masses; $`\psi _s(𝐤,𝐫)=\psi _s(𝐤,𝐫)`$, $`𝐝(𝐤,𝐫)=𝐝(𝐤,𝐫)]`$. At $`H0`$, $`\psi (𝐤,𝐫)`$ and $`𝐝(𝐤,𝐫)`$ do not depend on $`𝐫`$. Electron spin susceptibility tensor, $`\chi _{i,j}`$, at $`T=0`$ for a singlet phase is $`\chi _{i,j}=0`$ whereas for a triplet phase is given by<sup>27</sup> :
$$\chi _{i,j}=\chi _0<\delta _{i,j}\frac{d_i^{}(𝐤)d_j(𝐤)}{𝐝^{}(𝐤)𝐝(𝐤)}>_𝐤,$$
(8)
where $`\delta _{i,j}=1`$ if $`i=j`$ and $`\delta _{i,j}=0`$ if $`ij`$; $`<|𝐝(𝐤)|^2>_𝐤=1`$, $`<\mathrm{}>_𝐤`$ means an averaging over the FS. \[Here, we consider only unitary triplet phases<sup>27</sup> (i.e., $`d_a(𝐤)d_c^{}(𝐤)=d_a^{}(𝐤)d_c(𝐤)`$)\].
At first we consider the case $`𝐇𝐛^{}(𝐳)`$. In singlet phase (5), superconductivity is destroyed by paramagnetic effects in arbitrary directed magnetic field. In a triplet phase (6), as it follows from Eq. (7), $`d_b^{}(𝐤)d_z(𝐤)`$ component is responsible for the deviation of the spin susceptibility $`\chi _b^{}\chi _{zz}`$ from $`\chi _0`$. If $`d_b^{}(𝐤)0`$ there exist two related phenomena: the paramagnetic destructive mechanism against superconductivity and a change of the Knight shift at $`T<T_c(H)`$. Let us calculate the upper critical field for $`𝐇𝐛^{}`$. By using a common approach<sup>28</sup> to the upper critical field of a clean superconductor<sup>25</sup> with open electron orbits and with one-component order parameter, it is possible to prove that Eq. (5) of Ref. ,
$`\mathrm{\Delta }(x)=`$ $`{\displaystyle \frac{g}{2}}{\displaystyle _{|xx_1|>d}}{\displaystyle \frac{2\pi Tdx_1}{v_F\mathrm{sinh}(\frac{2\pi T|xx_1|}{v_F})}}J_0\left[{\displaystyle \frac{2\alpha \mu _BH(xx_1)S_z}{v_F}}\right]`$ (11)
$`\times J_0(2\lambda \mathrm{sin}[{\displaystyle \frac{\omega _c(xx_1)}{2v_F}}]\mathrm{sin}[{\displaystyle \frac{\omega _c(x+x_1)}{2v_F}}\left]\right)`$
$`\times \mathrm{cos}[{\displaystyle \frac{2\mu _BH(xx_1)S_z}{v_F}}]\mathrm{\Delta }(x_1),`$
is extended to a singlet phase $`\psi _s(𝐤,𝐫)f(𝐤)\mathrm{\Delta }(x)`$ as well as to the triplet phases $`𝐝_\mathrm{𝟏}(𝐤,𝐫)(d_a=1,d_c=0,d_b^{}=0)f(𝐤)\mathrm{\Delta }(x)`$ and $`𝐝_\mathrm{𝟐}(𝐤,𝐫)(d_a=0,d_c=0,d_b^{}=1)f(𝐤)\mathrm{\Delta }(x)`$. \[Here, $`<|f(𝐤)|^2>_𝐤=1`$; $`g`$ is an effective electron interaction constant, $`d`$ is a cutoff distance; $`\alpha =\sqrt{2}t_b/t_a`$, $`\omega _c=ev_FHc^{}/c`$, $`\lambda =4t_c/\omega _c`$; $`\mu _B`$ is a Borh magneton, $`e`$ and $`c`$ are the electron charge and the velocity of light, correspondingly; $`S_z=1`$ for singlet and for $`𝐝_2`$-triplet phases whereas $`S_z=0`$ for $`𝐝_1`$-triplet phase. By solving Eq. (8) numerically for $`S_z=1`$, $`\alpha =0.17`$, $`|dH^b^{}/dT|_{Tc}2T/K`$, $`v_F=10^7cm/sec,t_c3K,T_c(0)=1.14K,c^{}=13.6A`$ (see Refs. \[1,2,19-21,29\]), we found that the calculated value of the paramagnetic limited critical field, $`H_p^b^{}1.31.4T`$, is 4-5 times less than the experimental ones<sup>20-21</sup> (see Fig. 1). A similar analysis for $`𝐝_\mathrm{𝟏}`$-triplet phase (which is not paramagnetically limited) shows that superconductivity survives at $`H^b^{}6T`$ and $`T0.20.25K`$ in qualitative agreement with experiments<sup>20,21</sup> (see Fig. 1). On the basis of the calculation of $`H_p^b^{}`$ and $`H^b^{}`$, we can conclude that $`|d_b^{}(𝐤)||d_z(𝐤)|0`$ in Eq. (7) and thus $`\chi _b^{}\chi _{zz}\chi _0`$. Note that the recent Knight shift measurements<sup>12</sup> are also in favor of $`\chi _b^{}=\chi _0`$ below $`T_c(H)`$.
If we consider the case $`𝐇𝐚(𝐱)`$ then $`d_ad_x`$-component of the order parameter (6) is responsible for the destructive paramagnetic effects against superconductivity and for the change of the Knight shift at $`T<T_c(H)`$ (see Eq. (7)). Let us calculate the critical field for $`𝐇𝐚`$ in $`𝐝_\mathrm{𝟏}(𝐫)(d_a0,d_c0,d_b^{}=0)\mathrm{\Delta }(x)`$ triplet phase (which is paramagnetically limited). The corresponding linearized gap equation can be obtained from the common Eq. (5) of Ref. :
$`\mathrm{\Delta }(x)=`$ $`{\displaystyle \frac{g}{2}}{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi }{2\pi }}{\displaystyle _{|xx_1|>\sqrt{2}d|\mathrm{sin}\varphi |/\gamma }^{\mathrm{}}}`$ (15)
$`\mathrm{\Delta }(x_1){\displaystyle \frac{\sqrt{2}\gamma \pi Tdx_1}{v_F\mathrm{sin}\varphi \mathrm{sinh}[\frac{\sqrt{2}\gamma \pi T|xx_1|}{v_F\mathrm{sin}\varphi }]}}`$
$`\times J_0\left({\displaystyle \frac{\sqrt{2}\lambda \gamma }{\mathrm{sin}\varphi }}\mathrm{sin}[{\displaystyle \frac{\omega _c(xx_1)}{2v_F}}]\mathrm{sin}[{\displaystyle \frac{\omega _c(x+x_1)}{2v_F}}]\right)`$
$`\times \mathrm{cos}[{\displaystyle \frac{\sqrt{2}\gamma \mu _BHS_z(xx_1)}{v_F\mathrm{sin}\varphi }}],`$
where $`\gamma =t_aa/(2t_bb)`$. Numerical solution of Eq. (9) (with the same values of parameters as Eq. (8)) shows that the best fitting of the data<sup>20,21</sup> at $`H1.5T`$ (see Fig. 1) corresponds to $`S_z0.9`$ (i.e., $`d_a0.9`$, $`\chi _a0.2\chi _0\chi _0`$) and $`|dH^a/dT|_{Tc}8T/K`$. The latter is in a good agreement with the experimental slopes<sup>20,21</sup> $`|dH^b^{}/dT|_{Tc}2T/K`$ since the value of $`t_b/t_a8.5`$ is known<sup>29</sup>. Note that the accuracy of our calculations does not allow us to distinguish between the triplet phases with $`d_c=0`$ and $`|d_a|>|d_c|`$.
Summarizing, our analysis of the experimental critical fields<sup>20,21</sup> measured in (TMTSF)<sub>2</sub>PF<sub>6</sub> at $`P6kbar`$ has shown that paramagnetic destructive effects against superconductivity do not affect $`H^b^{}`$ whereas $`H^a`$ is paramagnetically limited at $`H1.5T`$. These are naturally explained within a triplet scenario of superconductivity<sup>9,16,19-21</sup> with the triplet order parameter (1). We suggest to measure the Knight shift along the $`𝐚`$-axis at $`H1.5T`$ and $`T<T_c(H)`$ to prove the order parameter (1). Note that temperature dependence of the critical field along a-axis, $`H^a(T)`$, changes drastically<sup>20,21</sup> at $`H1.5T`$. We speculate that at $`H1.5T`$ there may appear a triplet phase with $`𝐝(𝐤)𝐇`$, which minimizes the magnetic contribution to the free energy<sup>30</sup>. Nevertheless, we cannot completely exclude another possibility - the appearance of the LOFF state at $`H1.5T`$ for $`𝐇𝐚`$. Note that our theoretical analysis of the critical fields is based on the the Fermi-liquid picture<sup>29</sup> proved at $`P6kbar`$ in (TMTSF)<sub>2</sub>PF<sub>6</sub>. At higher pressures, $`P9.8kbar`$, the behavior of (TMTSF)<sub>2</sub>PF<sub>6</sub> may deviate from the Fermi liquid one<sup>31</sup>.
At the end of the paper, we would like to make a few comments based on symmetry arguments. We classify the possible triplet phases in the case of strong spin-orbit coupling for orthorhombic (D<sub>2h</sub>) and triclinic (C<sub>i</sub>) point group symmetries (see Table I), where the matrix order parameter $`\widehat{\mathrm{\Delta }}(𝐤)=d_i(𝐤)\widehat{\tau }_i`$, ($`\widehat{\tau }_i=i\widehat{\sigma }_i\widehat{\sigma }_y`$; $`\widehat{\sigma }_i`$ are the Pauli matrices). As it seen from Table I, there are no degenerated orbital states, thus a time reversal symmetry is broken only if a nonunitary triplet phase appears<sup>27</sup>. In our particular case, this happens when $`d_a(𝐤)d_c^{}(𝐤)d_a^{}(𝐤)d_c(𝐤)`$. Using the expression for a gap in a quasi-particle spectrum<sup>27</sup>, $`\delta (k)`$ = $`|𝐝(𝐤)|`$ (the unitary case), it is possible to make sure that there are no generic phases with the lines of zeros on the FS in accordance with a common theorem<sup>32</sup>. This is in agreement with the experimental data<sup>26,33</sup> which seem to be in favor of fully gapped FS and against the existence of isolated zeros on the FS<sup>32</sup>. Therefore, we speculate that the orbital part of the order parameter is likely $`d_a(𝐤)d_c(𝐤)sgn(k_a)`$ which corresponds to a fully gapped Q1D sheets of the FS. From Table I, it is possible to conclude that, for a triclinic space group of (TMTSF)<sub>2</sub>PF<sub>6</sub>, the most generic case is $`d_a0`$, $`d_c0`$ and $`d_b^{}0`$. However, it is known<sup>1,2,34</sup> that the spin dependent interactions in a SDW phase of (TMTSF)<sub>2</sub>PF<sub>6</sub> (which has a common boundary with the superconducting phase) result in an alignment of spins along $`𝐛^{}`$-axis. Therefore, it is natural to expect the form (1) for the superconducting order parameter corresponding to the absence of the BCS pairs with $`S_b^{}=0`$ (see Eq. (2)).
One of us (A.G.L.) is thankful to Agterberg, N. N. Bagmet, K. Behnia, S. Brown, E. V. Brusse, P. M. Chaikin, T. Ishiguro, H. Fukuyama, I.J. Lee, P. Lee, Y. Maeno, V. P. Mineev, M. J. Naughton, K. Oshima, M. Sigrist, V. M. Yakovenko for useful discussions. A.G.L. is especially thankful to S. Brown, P. M. Chaikin, I. J. Lee and M. J. Naughton for fruitful and numerous discussions during a workshop organized by M. J. Naughton.
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# Killing Spinors and SYM on Curved Spaces
## 1 Introduction
It is well known that Poincaré supersymmetric gauge theories retain a certain fraction of their supersymmetry when placed on Ricci flat manifolds $`M`$ admitting covariantly constant spinors, simply by using these parallel spinors as the supersymmetry parameters. For the same reason string theory compactifications on such manifolds lead to space-time supersymmetry.
From the string or supergravity theory point of view it is almost equally natural to consider (maximally) supersymmetric compactifications of the form $`M_1\times M_2`$ where this time the $`M_i`$ are required to be Einstein manifolds admitting Killing spinors rather than covariantly constant spinors.
It is therefore natural to ask if super-Yang-Mills (SYM) theories retain some global supersymmetry when placed on backgrounds admitting Killing spinors. For instances, this question arises in the context of the AdS/CFT correspondence when considering curved wrapped D-branes, as e.g. in . It also ought to arise, for the same reason as in the case of branes wrapped over supersymmetric cycles of manifods admitting parallel spinors (see e.g. ), in the context of AdS-calibrations studied in .
Morally speaking, by virtue of the existence of Killing spinors, globally supersymmetric SYM theories should exist on such manifolds, and it should be possible to deduce their existence and properties directly, i.e. without having to pass through supergravity and the possibly arduous task of studying fluctuations around a given (perhaps not even maximally) supersymmetric background.
It appears to be almost folklore knowledge that for the four-dimensional SYM theories addition of a suitable mass term for the scalars in the vector multiplet is sufficient to ensure supersymmetry on a background with Killing spinors. However, I am not aware of any general and systematic, i.e. not tied to a particular dimension, discussion of these matters.
Here, in addition to reproducing these results for $`n=4`$, we will find two families of Killing SYM theories for $`n4`$, both of them with the same field content as their Poincaré supersymmetric counterparts but with different actions and (generically) different supersymmetry transformation laws. From the results one can see in retrospect that the four-dimensional case (with equal masses for all the scalars, no other scalar potential terms, no mass term for the fermions) is sufficiently special to preclude a straightforward extrapolation to other dimensions.
One of these families of theories, given in (35), has the presumably unsurprising property of differing from the flat space theory by mass terms for the scalar fields and (unless the space-time dimension is $`n=4`$) fermions. I would suspect that these theories can be readily extracted from the supergravity literature. However, even one of the simplest members of this family of theories we will find, namely the $`N=2`$ theory on $`AdS_5`$, was only constructed very recently in , so perhaps these theories are not so well known after all.
The other family, given in (41), existing in all dimensions $`n7`$, has the more curious feature of requiring Chern-Simons-like cubic couplings of the scalar fields for supersymmetry and appears to be new.
One unexpected consequence of this is the existence of two inequivalent supersymmetric curved space counterparts of the three-dimensional $`N=4`$ SYM theory on locally AdS spaces: one with with fermionic and bosonic mass terms and modified supersymmetry transformation rules, the other with the same supersymmetry transformation rules as in flat space but with a cubic interaction term for the scalars instead of a mass term.
If these Killing SYM theories are realized as world volume theories of certain curved D-branes - and the role of wrapped branes e.g. in studies of the AdS/CFT correspondence certainly suggests that they should be thought of as being equipped with a supersymmetric world volume dynamics - then certainly the fundamental properties of these theories, supersymmetric vacua, BPS configurations etc., need to be understood. Here we will just discuss one simple but intriguing aspect of these theories, namely the counterpart of what is usually called the Coulomb branch. What we will find is that the structure of the vacua with unbroken supersymmetries in these theories differs quite markedly from that in the Poincaré supersymmetric theories - e.g. in the sense that generically there is no continuous family of maximally supersymmetric vacua, i.e. all the flat directions of the potential are lifted by a contribution to the potential induced by the curvature.
This in itself may not be terribly surprising, given the known results about other quantum field theories in AdS space-times. However, it certainly calls for a reappreciation of these issues in the context of brane dynamics.
As signs are crucial when it comes to checking supersymmetry, section 2 and an appendix serve to establish the conventions and notation and to provide some background information regarding supersymmetry variations in curved backgrounds and Killing spinors. In section 3, the two classes of theories mentioned above are described, and in section 4 the supersymmetry algebra in these models is (partially) calculated. Section 5 contains some sample calculations in these models, dealing mainly with the absence of a maximally supersymmetric Coulomb branch and the existence of a half-BPS Coulomb branch.
There are a large number of open issues, e.g. a more conceptual understanding of the existence of these theories (which here have been constructed more or less by brute force), and their superalgebraic underpinning, the study of the corresponding quantum theories, spaces of vacua, BPS configurations, application to worldvolume theories of curved D-branes, etc. Work on these and related issues (the original motivation for looking at (and hence first for) these theories was part of an attempt to find a topological counterpart of the AdS/CFT correspondence) is in progress, and I will briefly come back to these issues in the concluding section 6.
###### Acknowledgments.
I am grateful to Jose Figueroa-O’Farrill, Edi Gava, K.S. Narain, Martin O’Loughlin, Seif Randjbar-Daemi and George Thompson for discussions and suggestions at various stages of this work and for encouraging me to finally write up these results. This work was supported in part by the EC under the TMR contract ERBFMRX-CT96-0090.
## 2 Background
### 2.1 SYM Theories in Flat Space
We will consider the $`N=1`$ SYM theories in $`d=2+1,3+1,5+1`$ and $`9+1`$ dimensions as well as their dimensional reductions to $`nd`$ dimensions. This dimensional reduction could be along space-like directions to produce the standard Minkowski signature SYM theories, but it could also involve the time-direction to give rise to hermitian SYM actions in Euclidean signature . Thus in particular these theories include the $`N=2`$ and $`N=4`$ theories in $`n=3+1`$ as well as their Euclidean counterparts.
Quite generally, for all these theories the Lagrangian in $`d`$ or $`n`$ dimensional flat space can be written in the compact form
$$L_{SYM}=\frac{1}{2}F_{MN}F^{MN}+\overline{\mathrm{\Psi }}\mathrm{\Gamma }^MD_M\mathrm{\Psi }.$$
(1)
Here the following conventions have been used:
* Capital indices $`L,M,N,\mathrm{}`$ run from $`0`$ to $`d1`$.
* The gauge fields $`A_M`$ and $`\mathrm{\Psi }`$ only depend on the coordinates $`x^\mu `$, $`\mu =0,\mathrm{},n1`$ or $`\mu =1,\mathrm{},n`$ depending on whether one performs a space or time reduction. Thus $`A_\mu `$ is an $`n`$-dimensional gauge field and the remaining $`(dn)`$ components $`A_m\varphi _m`$ are scalar fields transforming as a vector under the manifest R-symmetry group $`SO(dn)`$ or $`SO(dn1,1)`$.
* A trace is implicit in (1) for the interacting (non-Abelian) theories, the fields transforming in the adjoint representation of the gauge group $`G`$,
$$A_M=A_M^iT_i,\mathrm{\Psi }=\mathrm{\Psi }^iT_i.$$
(2)
These Lie algebra indices will usually be suppressed in the following.
* $`A_M`$ will be taken to be anti-hermitian, so that the field strength tensor is
$$F_{MN}=_MA_N_NA_M+[A_M,A_N]$$
(3)
(no factors of $`i`$).
* The $`\mathrm{\Gamma }^M`$ are $`d`$-dimensional unitary gamma matrices and satisfy
$$\{\mathrm{\Gamma }_M,\mathrm{\Gamma }_N\}=\eta _{MN}$$
(4)
with
$$\eta _{MN}=\text{diag}(1,\underset{d1}{\underset{}{+1,\mathrm{},+1}}).$$
(5)
* $`\mathrm{\Psi }`$ is an anticommuting spinor in $`d`$ dimensions satisfying the condition
$$\begin{array}{cc}d=2+1:& \text{Majorana}\hfill \\ d=3+1:& \text{Majorana or Weyl}\hfill \\ d=5+1:& \text{Weyl}\hfill \\ d=9+1:& \text{Majorana-Weyl}\hfill \end{array}$$
(6)
* $`\overline{\mathrm{\Psi }}`$ is the Dirac adjoint of $`\mathrm{\Psi }`$ defined by
$$\overline{\mathrm{\Psi }}=\mathrm{\Psi }^{}A_{},$$
(7)
where $`A_{}=\mathrm{\Gamma }_0`$ satisfies
$$\mathrm{\Gamma }_M^{}=A_{}\mathrm{\Gamma }_MA_{}^1.$$
(8)
* $`D_M`$ is the gauge covariant derivative,
$`D_\mu \mathrm{\Psi }`$ $`=`$ $`_\mu \mathrm{\Psi }+[A_\mu ,\mathrm{\Psi }]`$
$`D_m\mathrm{\Psi }`$ $`=`$ $`[\varphi _m,\mathrm{\Psi }].`$ (9)
With these conventions, and the rule
$$(\chi ^{}\psi )^{}=\psi ^{}\chi $$
(10)
for anticommuting spinors $`\chi ,\psi `$, the above action is hermitian. Explicitly it reads
$`L_{SYM}`$ $`=`$ $`\frac{1}{2}F_{\mu \nu }F^{\mu \nu }D_\mu \varphi _mD^\mu \varphi ^m\frac{1}{2}[\varphi _m,\varphi _n][\varphi ^m,\varphi ^n]`$ (11)
$`+`$ $`\overline{\mathrm{\Psi }}\mathrm{\Gamma }^\mu D_\mu \mathrm{\Psi }+\overline{\mathrm{\Psi }}\mathrm{\Gamma }^m[\varphi _m,\mathrm{\Psi }].`$
In flat space it is invariant under the supersymmetry transformations
$`\delta A_M^i`$ $`=`$ $`(\overline{\epsilon }\mathrm{\Gamma }_M\mathrm{\Psi }^i\overline{\mathrm{\Psi }}^i\mathrm{\Gamma }_M\epsilon )`$
$`\delta \mathrm{\Psi }^i`$ $`=`$ $`\mathrm{\Gamma }^{MN}F_{MN}^i\epsilon `$
$`\delta \overline{\mathrm{\Psi }}^i`$ $`=`$ $`\overline{\epsilon }\mathrm{\Gamma }^{MN}F_{MN}^i`$ (12)
(modulo total derivatives) when $`\epsilon `$ is a constant spinor also satisfying the condition (6). Here
$$\mathrm{\Gamma }^{MN}=\frac{1}{2}[\mathrm{\Gamma }^M,\mathrm{\Gamma }^N].$$
(13)
In the non-Abelian case, vanishing of the quartic fermionic terms arising from the variation of the gauge field in the fermion kinetic term requires a Fierz identity to hold, which is satisfied by virtue of the conditions (6). The free theories are invariant under (12) without this requirement.
For brevity we will frequently refer to the dimensional reduction of the $`d`$-dimensional $`N=1`$ theory to $`n`$ dimensions as the $`(d,n)`$ theory. Thus the $`(10,4)`$ theory is $`N=4`$ SYM in four dimensions and e.g. $`(6,5)`$ refers to the five dimensional $`N=2`$ theory with one Dirac spinor (actually two symplectic Majorana spinors, hence $`N=2`$) and one real scalar in addition to the five-dimensional gauge field. We will mostly consider standard space-like reductions, but following the procedure outlined in one can also obtain Euclidean SYM theories by performing the dimensional reduction along the time-direction. These will be discussed seperately in section 3.4.
### 2.2 Supersymmetry Variations in Curved Space
Let us now consider what happens when one tries to place these theories (after the appropriate dimensional reduction) on a curved background. To be specific, denote by $`(M,g)`$ a (pseudo-)Riemannian $`n`$-dimensional spin manifold with metric $`g_{\mu \nu }`$.
There is of course no problem with writing down the action (1) on $`M`$ by introducing a vielbein $`e_\mu ^a`$, a spin connection $`\omega _\mu ^{ab}`$, etc. Just to further pin down the conventions, the spin connection part of the covariant derivative is
$$_\mu \mathrm{\Psi }=_\mu \mathrm{\Psi }+\frac{1}{4}\mathrm{\Gamma }_{ab}\omega _\mu ^{ab}\mathrm{\Psi }.$$
(14)
The real issue is whether this theory will have any supersymmetry, the point being that constant spinors $`\epsilon `$ will in general not exist on $`M`$ while using non-constant supersymmetry parameters in (12) will lead to a non-zero variation of the action through terms depending on the derivatives of $`\epsilon `$.
By just keeping track of the terms that depend on the (covariant) derivatives of $`\epsilon `$, it is straightforward to compute the supersymmetry variation of the action on $`M`$ and the result is (once again modulo total derivatives)
$$\delta L_{SYM}=\left[(_\mu \overline{\epsilon })\mathrm{\Gamma }^{NL}\mathrm{\Gamma }^\mu \mathrm{\Psi }+\overline{\mathrm{\Psi }}\mathrm{\Gamma }^\mu \mathrm{\Gamma }^{NL}(_\mu \epsilon )\right]F_{NL}.$$
The $`F_{NL}`$-terms encapsulate the curvature terms $`F_{\mu \nu }`$ as well as derivative terms of the scalars and scalar commutator terms. Note that the supersymmetry parameters $`\epsilon `$ are gauge singlets so that the covariant derivative $`_\mu \epsilon `$ includes only the spin connection but not the gauge field. The gauge and gravitational covariant derivative will be denoted by $`D_M`$.
### 2.3 Killing Spinor Equations
The most immediate non-trivial solutions to $`\delta L_{SYM}=0`$ (2.2) are of course provided by parallel spinors,
$$_\mu \epsilon =0.$$
(15)
The resulting supersymmetric theories and their Euclidean/topological counterparts on Ricci-flat special holonomy manifolds are reasonably well understood (see e.g. and references therein) and will not be considered further in this paper.
A natural generalization of a parallel spinors is a Killing spinor, i.e. a Dirac spinor $`\eta `$ in $`n`$ dimensions satisfying an equation of the form
$$_\mu \eta =\alpha \gamma _\mu \eta $$
(16)
where the $`\gamma _\mu `$ are $`n`$-dimensional $`\gamma `$-matrices and $`\alpha `$ is some real or imaginary constant.<sup>1</sup><sup>1</sup>1Actually, while in the mathematics literature the name Killing spinor is usually reserved for spinors satisfying (16), in the supergravity literature any equation of the form $`_\mu \eta =M_\mu (x)\eta `$ arising from setting to zero the gravitino variation in a bosonic background is called a Killing spinor equation. Here $`M_\mu (x)`$ is typically made up from contractions of supergravity antisymmetric tensor background fields with gamma matrices, hence the explicit $`x`$-dependence. Here we have no such background fields, and thus we are left with (16). These equations have been thoroughly investigated in the supergravity and mathematics literature, at least in the case when $`M`$ is compact and Riemannian - see e.g. and and the references therein for the mathematical and Kaluza-Klein supergravity aspects respectively. For recent work on the pseudo-Riemannian case see .
To write this back in $`d`$-dimensional terms, it is not correct to just consider an equation like $`_\mu \epsilon =\alpha \mathrm{\Gamma }_\mu \epsilon `$ as this would for instance be incompatible with a chirality condition on $`\epsilon `$. Instead, we postulate the slightly more general Killing spinor equation
$$_\mu \epsilon =\alpha \mathrm{\Gamma }_\mu \mathrm{\Gamma }\epsilon ,$$
(17)
where $`\epsilon `$ denotes the $`d`$-dimensional (chiral, Majorana, …) spinor and where $`\mathrm{\Gamma }`$ could be an arbitrary element of the Clifford algebra generated by the $`\mathrm{\Gamma }^M`$. In fact we will be more specific than that and consider the case in which $`\mathrm{\Gamma }`$ is a monomial constructed from the ‘internal’ gamma matrices $`\mathrm{\Gamma }^m`$, i.e. a completely anti-symmetrized product of $`0pdn`$ gamma matrices. When it is necessary to indicate the degree $`p`$, we will write $`\mathrm{\Gamma }^{[p]}`$ instead of $`\mathrm{\Gamma }`$. Then one in particular has $`(\mathrm{\Gamma })^2=\pm I`$. Generalizations of this are certainly possible but will not be explored here.
This equation now preserves chirality when $`p`$ is odd, and so it can also be used in the theories arising upon dimensional reduction of the chiral $`N=1`$ theories. Moreover, the freedom in the choice of $`\mathrm{\Gamma }`$ may allow one to find different supersymmetric theories for a given field content (on manifolds satisfying either the same or different integrability conditions of the Killing spinor equation). We will see examples of this below.
Finally, this generalized Killing spinor equation, when written out in $`n`$-dimensional terms, will always reduce to the standard Killing spinor equation of the type (16) for (appropriate linear combinations of) $`n`$-dimensional Dirac spinors<sup>2</sup><sup>2</sup>2or perhaps to some simple variant thereof when $`n`$ is even,
$$_\mu \eta =i\alpha \gamma _\mu \gamma ^{(n+1)}\eta $$
(here $`\gamma ^{(n+1)}`$ is the chirality operator). This equation can be mapped to the standard equation (16) by passing to the unitarily equivalent representation $`\stackrel{~}{\gamma }_\mu =i\gamma _\mu \gamma ^{(n+1)}`$., and therefore the standard existence criteria for ordinary Killing spinors can be applied to (17).
### 2.4 Integrability Conditions
The (first) integrability condition arising from the Killing spinor equation (17) is, taking commutators and recalling (14),
$$\frac{1}{4}\mathrm{\Gamma }_{ab}\mathrm{\Omega }_{\mu \nu }^{ab}\epsilon =\alpha ^2[\mathrm{\Gamma }_\nu \mathrm{\Gamma },\mathrm{\Gamma }_\mu \mathrm{\Gamma }]\epsilon ,$$
(18)
where $`\mathrm{\Omega }_{\mu \nu }^{ab}`$ denotes the curvature tensor of the spin connection $`\omega _\mu ^{ab}`$. Upon contraction with $`\mathrm{\Gamma }^\nu `$ this leads to
$$R_{\mu \nu }\mathrm{\Gamma }^\nu \epsilon =2\alpha ^2g_{\mu \nu }[(n2)\mathrm{\Gamma }\mathrm{\Gamma }^\nu \mathrm{\Gamma }+\mathrm{\Gamma }^\nu \mathrm{\Gamma }^\lambda \mathrm{\Gamma }\mathrm{\Gamma }_\lambda \mathrm{\Gamma }]\epsilon .$$
(19)
For $`\mathrm{\Gamma }=\mathrm{\Gamma }^{[p]}`$ ‘internal’ in the sense described before, so that $`\mathrm{\Gamma }^{[p]}`$ commutes (anticommutes) with all the $`\mathrm{\Gamma }^\mu `$ if $`p`$ is even (odd), one finds
$$R_{\mu \nu }\mathrm{\Gamma }^\nu \epsilon =4\alpha ^2(1)^p(\mathrm{\Gamma }^{[p]})^2(n1)g_{\mu \nu }\mathrm{\Gamma }^\nu \epsilon .$$
(20)
In the Riemannian case, an equation of the form $`A_{\mu \nu }\mathrm{\Gamma }^\nu \epsilon =0`$ implies $`A_{\mu \nu }=0`$. This can be seen by multiplying by $`A_\lambda ^\mu \mathrm{\Gamma }^\lambda `$. Thus in this case (19) implies that
$$R_{\mu \nu }=4\alpha ^2(1)^p(\mathrm{\Gamma }^{[p]})^2(n1)g_{\mu \nu }.$$
(21)
and hence that $`(M,g)`$ is an Einstein manifold. In particular, for $`\mathrm{\Gamma }=I`$ or, equivalently, for the ordinary Killing spinor equation (16), one obtains
$$R_{\mu \nu }=4\alpha ^2(n1)g_{\mu \nu }.$$
(22)
Thus Killing spinors (16) for imaginary $`\alpha `$ (referred to as real Killing spinors in the mathematics literature) lead to positive curvature, and spinors with real $`\alpha `$ (imaginary Killing spinors) lead to negative curvature.
This unfortunate clash in terminology is due to the fact that typically in the mathematics literature the conventions for Clifford algebras are such that $`\{\mathrm{\Gamma }_\mu ,\mathrm{\Gamma }_\nu \}=2g_{\mu \nu }`$, the opposite of the convention used here. Perhaps a more invariant and informative terminology would have been to call a Killing spinor positive or negative according to whether the integrability condition leads to positive or negative curvature, and we will adopt this terminology from now on.
In general, the sign of the curvature depends on $`\alpha `$, $`p`$ and on $`(\mathrm{\Gamma }^{[p]})^2=\pm I`$. For the chiral $`N=1`$ theories and their descendants, $`p`$ has to be odd in order for the Killing spinor equation to be compatible with the chirality of $`\epsilon `$.
The integrability condition (21) is not sufficient for the existence of Killing spinors (not every Einstein manifold admits Killing spinors) but fortunately an analysis of the higher integrability conditions can be side-stepped by relating Killing spinors on $`M`$ to parallel spinors on another Ricci flat manifold and hence establishing existence of Killing spinors directly - see for positive Killing spinors and for negative Killing spinors.
In the pseudo-Riemannian case, (21) is neither necessary nor sufficient. An argument like the above only leads to the conclusion that for each value of $`\mu `$ the vector $`V_{(\mu )}`$ with components $`V_{(\mu )}^\nu =A_\mu ^\nu `$ is null, with the additional constraint $`g^{\mu \nu }A_{\mu \nu }=0`$. In the case of parallel spinors, the resulting Ricci-null Lorentzian manifolds which are not Ricci flat were recently investigated in detail in (see also ). By the same token, one might suspect that there are non-Einstein Lorentzian manifolds admitting Killing spinors. There are indeed such examples for negative pseudo-Riemannian Killing spinors whereas a pseudo-Riemannian manifold admitting a positive Killing spinor is necessarily Einstein . Nevertheless, in the following we will simply assume that (21) holds. In this way we will certainly miss some solutions (in the negative curvature case), but as a first orientation this is good enough.
### 2.5 The Supersymmetry Variation for Killing Spinors
In order to plug (17) into the formula (2.2) for $`\delta L_{SYM}`$, one first needs an expression for $`_\mu \overline{\epsilon }`$. By using the fact that
$$(\mathrm{\Gamma }^{[p]})^{}=\eta _pA_{}\mathrm{\Gamma }^{[p]}A_{}^1,$$
(23)
where
$$\eta _p=(1)^{\left(\genfrac{}{}{0pt}{}{p+1}{2}\right)},$$
(24)
one obtains
$$_\mu \overline{\epsilon }=\eta _p\alpha ^{}\overline{\epsilon }\mathrm{\Gamma }\mathrm{\Gamma }_\mu .$$
(25)
Thus
$`\delta L_{SYM}`$ $`=`$ $`[\eta _p\alpha ^{}\overline{\epsilon }\mathrm{\Gamma }\mathrm{\Gamma }_\mu \mathrm{\Gamma }^{NL}\mathrm{\Gamma }^\mu \mathrm{\Psi }+\alpha \overline{\mathrm{\Psi }}\mathrm{\Gamma }^\mu \mathrm{\Gamma }^{NL}\mathrm{\Gamma }_\mu \mathrm{\Gamma }\epsilon ]F_{NL}`$ (26)
$`=`$ $`2\text{Re}(\overline{\mathrm{\Psi }}\mathrm{\Gamma }^\mu \mathrm{\Gamma }^{NL}\mathrm{\Gamma }_\mu \mathrm{\Gamma }\epsilon ))F_{NL}.`$
Splitting the gamma matrices $`\{\mathrm{\Gamma }^M\}=\{\mathrm{\Gamma }^\mu ,\mathrm{\Gamma }^m\}`$ and using the standard identities
$`\mathrm{\Gamma }_\mu \mathrm{\Gamma }^{\nu \lambda }\mathrm{\Gamma }^\mu `$ $`=`$ $`(n4)\mathrm{\Gamma }^{\nu \lambda }`$
$`\mathrm{\Gamma }_\mu \mathrm{\Gamma }^{\nu m}\mathrm{\Gamma }^\mu `$ $`=`$ $`(n2)\mathrm{\Gamma }^{\nu m}`$
$`\mathrm{\Gamma }_\mu \mathrm{\Gamma }^{lm}\mathrm{\Gamma }^\mu `$ $`=`$ $`n\mathrm{\Gamma }^{lm},`$ (27)
one can evaluate this to find
$$\begin{array}{ccccc}\delta L_{SYM}& =& (n4)& [\eta _p\alpha ^{}\overline{\epsilon }\mathrm{\Gamma }\mathrm{\Gamma }^{\nu \lambda }\mathrm{\Psi }+\alpha \overline{\mathrm{\Psi }}\mathrm{\Gamma }^{\nu \lambda }\mathrm{\Gamma }\epsilon ]\hfill & F_{\nu \lambda }\hfill \\ & +& 2(n2)& [\eta _p\alpha ^{}\overline{\epsilon }\mathrm{\Gamma }\mathrm{\Gamma }^{\nu m}\mathrm{\Psi }+\alpha \overline{\mathrm{\Psi }}\mathrm{\Gamma }^{\nu m}\mathrm{\Gamma }\epsilon ]\hfill & D_\nu \varphi _m\hfill \\ & +& n& [\eta _p\alpha ^{}\overline{\epsilon }\mathrm{\Gamma }\mathrm{\Gamma }^{lm}\mathrm{\Psi }+\alpha \overline{\mathrm{\Psi }}\mathrm{\Gamma }^{lm}\mathrm{\Gamma }\epsilon ]\hfill & [\varphi _l,\varphi _m].\hfill \end{array}$$
Barring numerical coincidences, it is clear that this expression can only vanish when the expression in brackets vanishes all by itself, i.e. when
$$\text{Re}(\alpha \overline{\mathrm{\Psi }}\mathrm{\Gamma }^{NL}\mathrm{\Gamma }\epsilon )=0N,L.$$
(28)
This is only possible if $`\alpha =0`$ so that one is dealing with ordinary parallel spinors (and hence Ricci flat geometries in the Euclidean case and a few more possiblities for Lorentzian signature).
However, there is one numerical coincidence which occurs when $`d=n=4`$. In that case only the first line of (2.5) is present, but multiplied by $`n4=0`$. Thus e.g. for any solution to the ordinary Killing spinor equations (16) the $`N=1`$ theory in $`d=3+1`$ has a supersymmetry. The relevant gamma matrix identity shows that this is due to the fact that SYM theory is a theory of (non-Abelian) one-forms, and one might want to speculate about an analogous result for (non-Abelian?) two-form theories in $`d=5+1`$
## 3 Supersymmetric SYM Theories in Curved Space
On the basis of these preliminaries we can now write down two families of Dirac-Yang-Mills theories in curved space which are globally supersymmetric courtesy of the existence of solutions to a suitable Killing spinor equation. These theories generically differ from the simple SYM action $`L_{SYM}`$ by mass terms for both the scalars and the fermions and by a modified supersymmetry transformation rule for $`\mathrm{\Psi }`$. In addition, one class of these theories curiously has Chern-Simons-like cubic couplings for the scalar fields. Both of these families of theories turn out (a priori for no good reason) to be particularly simple in four dimenions, $`n=4`$, and we will start with that particular case.
### 3.1 Theories for $`n=4`$
Let $`L_{SYM}`$ be the $`(d,4)`$ Lagrangian, that is the dimensional reduction of the $`d`$-dimensional theory to $`4=3+1`$ dimensions, suitably covariantized, of course. Consider the Lagrangian
$`L`$ $`=`$ $`L_{SYM}8\alpha ^2{\displaystyle \underset{m=1}{\overset{dn}{}}}\varphi _m^2`$ (29)
$`=`$ $`\frac{1}{2}F_{MN}F^{MN}+\overline{\mathrm{\Psi }}\mathrm{\Gamma }^MD_M\mathrm{\Psi }8\alpha ^2{\displaystyle \underset{m=1}{\overset{dn}{}}}\varphi _m^2.`$
This action is invariant under the supersymmetry transformations (suppressing the Lie algebra labels on the fields)
$`\delta A_M`$ $`=`$ $`(\overline{\epsilon }\mathrm{\Gamma }_M\mathrm{\Psi }\overline{\mathrm{\Psi }}\mathrm{\Gamma }_M\epsilon )`$
$`\delta \mathrm{\Psi }`$ $`=`$ $`\mathrm{\Gamma }^{MN}\epsilon F_{MN}4\alpha {\displaystyle \underset{m=1}{\overset{dn}{}}}\varphi _m\mathrm{\Gamma }^m\mathrm{\Gamma }\epsilon ,`$ (30)
provided that $`\epsilon `$ satisfies the Killing spinor equation
$$_\mu \epsilon =\alpha \mathrm{\Gamma }_\mu \mathrm{\Gamma }\epsilon $$
(31)
where $`\mathrm{\Gamma }`$ is any odd, internal matrix with $`\mathrm{\Gamma }^2=\pm I`$. Here $`\alpha `$ has to be real for the $`d=4`$ and $`d=10`$ Majorana(-Weyl) theories, but can be either real or imaginary for the $`d=4,6`$ Weyl theories.
Indeed it is easy to see that due to the modification of the $`\mathrm{\Psi }`$-transformation the standard variation of $`L_{SYM}`$ given in (2.5) is exactly cancelled. But now one picks up terms linear in the scalar fields $`\varphi _m`$ from the Killing spinor equation, namely when the derivative $`D_M`$ in the fermionic kinetic term hits $`\epsilon `$ in the second term of $`\delta \mathrm{\Psi }`$. This gives a term proportional to $`\alpha ^2`$,
$$\delta L_{SYM}=\pm 16\alpha ^2[\overline{\epsilon }\mathrm{\Gamma }^m\mathrm{\Psi }\overline{\mathrm{\Psi }}\mathrm{\Gamma }^m\epsilon ]\varphi _m,$$
(32)
which is of course cancelled precisely by the variation of the mass term for the scalars.
Remarks:
1. We have just recovered the folklore statement that addition of mass terms for the scalars is sufficient to render four-dimensional SYM theories supersymmetric in a background admitting Killing spinors, provided that also the supersymmetry transformation rules of the fermions are changed appropriately.
2. In particular, the mass term is precisely the conformally invariant mass term arising in the conformally invariant wave operator
$$\mathrm{}\frac{1}{4}\frac{n2}{n1}R,$$
(33)
where $`R`$ is the scalar curvature
$$R=\pm 4\alpha ^2n(n1).$$
(34)
3. Note the striking similarity of the supersymmetry transformations with those of the special (i.e. superconformal) supersymmetry transformations as given e.g. in .
4. Similar linear terms in the transformations of the fermions also appear e.g. in the Wess-Zumino model in a curved background and are a rather generic feature of AdS supersymmetry \- for a recent review see .
5. Looking at the integrability conditions deduced before we learn that in particular the counterpart of the four-dimensional $`N=2`$ theory can be supersymmetric on Einstein manifolds of either positive or negative curvature admitting solutions of the Killing spinor equation, depending on whether $`\alpha `$ is chosen to be real or imaginary.
6. Likewise, the $`N=4`$ theory can be supersymmetric in both cases, depending on whether one chooses $`\mathrm{\Gamma }=\mathrm{\Gamma }^{[1]}`$ or $`\mathrm{\Gamma }=\mathrm{\Gamma }^{[3]}`$, with $`\alpha `$ real in both cases.
7. Even though the choice of $`\mathrm{\Gamma }`$ singles out one (or three) ‘internal’ directions, all the scalars have the same mass. This is a feature that will not persist in $`n4`$.
8. There is no mass term for the fermions. Once again, this is a feature peculiar to the $`n=4`$ theories.
9. Finally, it may be possible to construct this theory as a rigid limit of conformal supergravity in four dimensions.<sup>3</sup><sup>3</sup>3I thank Ergin Sezgin for this suggestion.
### 3.2 Family A: Theories for $`n5`$ with $`\mathrm{\Gamma }=\mathrm{\Gamma }^{[1]}`$
We will now consider the case where $`\mathrm{\Gamma }`$ is just a single internal gamma matrix which we will call $`\mathrm{\Gamma }^1`$. In particular, $`(\mathrm{\Gamma })^2=+I`$. Now consider the following action
$$L=L_{SYM}4\alpha ^2[(n2)\underset{m=1}{\overset{dn}{}}\varphi _m^2+(n4)\varphi _1^2](n4)\alpha \overline{\mathrm{\Psi }}\mathrm{\Gamma }^1\mathrm{\Psi }.$$
(35)
As it stands this action makes sense for the $`(d=4,n<4)`$, $`(d=6,n<6)`$, and the fermionic mass term is hermitian provided that $`\alpha `$ is imaginary (cf. the Appendix) and this rules out the $`(d=10,n4)`$ theories. We could also allow $`(d=10,n=4)`$ and $`\alpha `$ real, but in that case the action reduces to the one discussed above.
This action is invariant under the supersymmetry transformations
$`\delta A_M`$ $`=`$ $`(\overline{\epsilon }\mathrm{\Gamma }_M\mathrm{\Psi }\overline{\mathrm{\Psi }}\mathrm{\Gamma }_M\epsilon )`$
$`\delta \mathrm{\Psi }`$ $`=`$ $`\mathrm{\Gamma }^{MN}\epsilon F_{MN}4\alpha [{\displaystyle \underset{m=1}{\overset{dn}{}}}\varphi _m\mathrm{\Gamma }^m\mathrm{\Gamma }^1\epsilon +(n4)\varphi _1\epsilon ]`$ (36)
provided that $`\epsilon `$ satisfies the Killing spinor equation
$$_\mu \epsilon =\alpha \mathrm{\Gamma }_\mu \mathrm{\Gamma }^1\epsilon $$
(37)
Remarks:
1. We now have mass terms both for the scalars and the fermions. The masses depend only on the space-time dimension $`n`$, not on the parent dimension $`d`$.
2. The mass of $`\varphi _1`$ differs from that of the $`\varphi _{m1}`$, but neither is the conformally invariant value unless $`n=2`$ when $`dn1`$ of the scalars are massless.
3. The integrability conditions tell us that these theories can only exist on Einstein manifolds of negative curvature - in particular locally AdS space-times.
4. The $`(6,5)`$-theory on AdS<sub>5</sub> has been constructed recently by Shuster in terms of symplectic Majorana spinors. It can be checked that, when these are reassembled into a Dirac spinor, his action and supersymmetries agree precisely with those given above when one sets $`d=6,n=5`$.
5. The R-symmetry of the action has been reduced from $`SO(dn)`$ (which is the manifest R-symmetry group of the Poincaré supersymmetric theory) to $`SO(dn1)`$.
The simplest of these theories is the $`(4,3)`$ theory, i.e. the $`N=2`$ theory in $`n=3`$. It differs from $`L_{SYM}`$ only by the fermionic mass term, and the supersymmetry transformation rules are the standard ones, i.e. we have
$`(d=4,n=3)`$ $`L=L_{SYM}+\alpha \overline{\mathrm{\Psi }}\mathrm{\Gamma }^1\mathrm{\Psi }`$ (38)
$`\delta \mathrm{\Psi }=\mathrm{\Gamma }^{MN}\epsilon F_{MN}.`$
This theory is supersymmetric almost by inspection. For $`n=3`$, the two first lines of (2.5) enter with opposite signs and the third line is absent. As $`\mathrm{\Gamma }_1`$ anticommutes with the $`\mathrm{\Gamma }^\mu `$ but commutes with the $`\mathrm{\Gamma }^{\mu \nu }`$, this is cancelled by the variation of the above fermionic mass term.
Let us also write down explictly the $`(6,3)`$-theory. It is given by
$`(d=6,n=3)`$ $`L=L_{SYM}4\alpha ^2(\varphi _2^2+\varphi _3^2)+\alpha \overline{\mathrm{\Psi }}\mathrm{\Gamma }^1\mathrm{\Psi }`$ (39)
$`\delta \mathrm{\Psi }=\mathrm{\Gamma }^{MN}\epsilon F_{MN}4\alpha (\varphi _2\mathrm{\Gamma }^2+\varphi _3\mathrm{\Gamma }^3)\mathrm{\Gamma }^1\epsilon .`$
### 3.3 Family B: Theories for $`n7`$ with $`\mathrm{\Gamma }=\mathrm{\Gamma }^{[3]}`$
If we want to use $`\mathrm{\Gamma }=\mathrm{\Gamma }^{[3]}`$ and still insist on this being an element of the ‘internal’ Clifford algebra, we obviously need $`nd3`$. Let us choose
$$\mathrm{\Gamma }=\mathrm{\Gamma }^{123}=\frac{1}{3!}\epsilon _{abc}\mathrm{\Gamma }^{abc},$$
(40)
so that $`(\mathrm{\Gamma })^2=I`$. Consider the action
$`L`$ $`=`$ $`L_{SYM}+4\alpha ^2[(n2){\displaystyle \underset{m=1}{\overset{dn}{}}}\varphi _m^2+(n4){\displaystyle \underset{a=1}{\overset{3}{}}}\varphi _a^2]`$ (41)
$``$ $`{\displaystyle \frac{(n4)\alpha }{3!}}\epsilon _{abc}\left[\overline{\mathrm{\Psi }}\mathrm{\Gamma }^{abc}\mathrm{\Psi }8\varphi ^a[\varphi ^b,\varphi ^c]\right].`$
Hermiticity of the mass term requires $`\alpha R`$.
This action is invariant under the supersymmetry transformations
$`\delta A_M`$ $`=`$ $`(\overline{\epsilon }\mathrm{\Gamma }_M\mathrm{\Psi }\overline{\mathrm{\Psi }}\mathrm{\Gamma }_M\epsilon )`$
$`\delta \mathrm{\Psi }`$ $`=`$ $`\mathrm{\Gamma }^{MN}\epsilon F_{MN}4\alpha [{\displaystyle \underset{m=1}{\overset{dn}{}}}\varphi _m\mathrm{\Gamma }^m+(n4){\displaystyle \underset{a=1}{\overset{3}{}}}\varphi _a\mathrm{\Gamma }^a]\mathrm{\Gamma }^{123}\epsilon `$ (42)
provided that $`\epsilon `$ satisfies the Killing spinor equation
$$_\mu \epsilon =\alpha \mathrm{\Gamma }_\mu \mathrm{\Gamma }^{123}\epsilon $$
(43)
Remarks:
1. The most striking property of this action is perhaps the appearance of the cubic term for the scalar fields. It looks like the dimensional reduction of a standard Chern-Simons term living in the three internal directions singled out by $`\mathrm{\Gamma }^{123}`$.
2. It is certainly suggestive of a supergravity origin of this term, but it would be desirable to find a pure gauge theory explanation for it as well.
3. Such terms can appear in the completely T-duality invariant D-brane world-volume actions discussed by Myers in , where they arise due to the coupling to non-trivial background antisymmetric tensor fields.
4. Some such term also appears in the off-shell rheonomic formulation of $`N=1`$ $`d=10`$ SYM in flat space - see \[30, (II.9.41)\]. The relation to the appeareance of such a term in the on-shell space-time action here is not clear (to the author) but may be worth understanding.
5. The integrability conditions once again lead to negative curvature because even though $`\alpha `$ is now real, one also has $`\mathrm{\Gamma }^2=I`$.
Apart from the $`(d=10,n=4)`$ theory already discussed above, for which there are neither fermionic mass terms nor Chern-Simons like couplings, the simplest theory is once again the three-dimensional $`(6,3)`$-theory with Lagrangian and supersymmetry transformation
$`(d=6,n=3)`$ $`L=L_{SYM}+\alpha (\overline{\mathrm{\Psi }}\mathrm{\Gamma }^{123}\mathrm{\Psi }8\varphi ^1[\varphi ^2,\varphi ^3])`$ (44)
$`\delta \mathrm{\Psi }=\mathrm{\Gamma }^{MN}\epsilon F_{MN}.`$
It is straightforward to check directly in this case that the action is supersymmetric: upon performing the supersymmetry variation, the terms involving $`F_{\mu \nu }`$ and $`D_\mu \varphi _n`$ arising from the variation of $`L_{SYM}`$ and the fermionic mass term cancel whereas those involving the commutator $`[\varphi _m,\varphi _n]`$ add up. The latter are then precisely cancelled by the variation of the cubic scalar term.
Note that we now have two obviously inequivalent curved space versions of the $`(6,3)`$-theory, i.e. of what in standard parlance be called the three-dimensional $`N=4`$ SYM theory ($`N=4`$ because in $`2+1`$ dimensions spinors are two-component real: $`so(2,1)sl(2,R)`$), one of them with a standard mass term for two of the three scalars (39), the other one instead with a Chern-Simons like term (44). Is there any interesting (duality?) relationship between these theories?
### 3.4 Euclidean Supersymmetric SYM Theories in Curved Space
Euclidean (or better perhaps: Riemannian) versions of the theories described above may be of interest for a variety of reasons, e.g. for D-brane instantons, within the Euclidean approach to the AdS/CFT correspondence, and in connection with Hull’s E-branes and an eye towards cohomological versions of these theories.
As explained in (see also ), a convenient way to obtain manifestly hermitian Euclidean SYM theories is by time-like dimensional reduction of any one of the standard Minkowskian SYM theories to a Lagrangian $`L_{ESYM}`$.
This construction naturally explains the features one has in the past come to expect of Euclidean supersymmetric theories, such as non-compact R-symmetry groups (namely the internal roation group which is now the Lorentz group $`SO(dn1,1)`$) and kinetic terms with the ‘wrong’ sign (namely the time-compoent of the gauge field, now a scalar from the point of view of the Euclidean space-time).
These theories then also automatically make sense on Riemannian manifolds and retain some fraction of their supersymmetry when this manifold admits parallel spinors. In this way one obtains cohomological theories on special holonomy manifolds with many beautiful features, studied for example from this point of view in .
Now let us, in analogy with what we did before, discuss the extension of these Euclidean SYM theories to supersymmetric theories on Riemannian manifolds admitting Killing spinors. Let us start with the $`n=4`$ theories of section 3.1. It is readily seen that the theory as it stands is supersymmetric also for the Euclidean theory provided that the mass term is chosen to be $`\eta ^{mn}\varphi _m\varphi _n`$, i.e.
$$L=L_{ESYM}8\alpha ^2\eta ^{mn}\varphi _m\varphi _n,$$
(45)
for any choice of (internal, odd) $`\mathrm{\Gamma }`$. In particular, $`\mathrm{\Gamma }`$ could be chosen to be equal to (or include) $`\mathrm{\Gamma }^0`$. The interesting thing about this is that according to (21) this changes the sign of the integrability condition. Whereas for $`\mathrm{\Gamma }=\mathrm{\Gamma }^1`$, say, the sign of the curvature is the sign of $`\alpha ^2`$, for $`\mathrm{\Gamma }=\mathrm{\Gamma }^0`$ it is minus the sign of $`\alpha ^2`$.
This may not be of great consequence in the present example since, as we saw before, we could anyhow obtain both signs by either choosing $`\alpha `$ to be real or imaginary (for the $`(6,4)`$ theory) or by choosing $`\mathrm{\Gamma }=\mathrm{\Gamma }^{[1]}`$ or $`\mathrm{\Gamma }=\mathrm{\Gamma }^{[3]}`$ (for the $`(10,4)`$ theory) - the integrability condition only depends on the square of $`\alpha \mathrm{\Gamma }`$.
Moreover, for $`n=4`$, but only for $`n=4`$, there is practically no dependence of the action on $`\mathrm{\Gamma }`$ (apart from the sign of the mass term) so that we do not get any essentially new theories in this way. But we will see below that in the other theories the freedom to choose $`\mathrm{\Gamma }`$ to include or not to include $`\mathrm{\Gamma }^0`$ gives us an added flexibility not present in the pseudo-Riemannian theories.
More or less the same modifications as above are required for the Family A theories of section 3.2. Provided that we define the mass terms as above and reintroduce the dependence of the sign of the mass term on $`\mathrm{\Gamma }^2=\pm I`$, as above, we obtain a supersymmetric Lagrangian. Thus essentially the only two different possibilities are
$`\mathrm{\Gamma }=\mathrm{\Gamma }^1`$ $`L=L_{ESYM}4\alpha ^2[(n2)\eta ^{mn}\varphi _m\varphi _n+(n4)\varphi _1^2](n4)\alpha \overline{\mathrm{\Psi }}\mathrm{\Gamma }^1\mathrm{\Psi }`$
$`\mathrm{\Gamma }=\mathrm{\Gamma }^0`$ $`L=L_{ESYM}+4\alpha ^2[(n2)\eta ^{mn}\varphi _m\varphi _n(n4)\varphi _0^2](n4)\alpha \overline{\mathrm{\Psi }}\mathrm{\Gamma }^0\mathrm{\Psi }`$
We know that $`\alpha `$ has to be imaginary for hermiticity of the fermionic mass term (this condition is the same for $`\overline{\mathrm{\Psi }}\mathrm{\Gamma }^1\mathrm{\Psi }`$ and $`\overline{\mathrm{\Psi }}\mathrm{\Gamma }^0\mathrm{\Psi }`$), and previously this forced the manifold to have negative curvature. However, now we actually gain something by being able to choose $`\mathrm{\Gamma }=\mathrm{\Gamma }^0`$ or $`\mathrm{\Gamma }=\mathrm{\Gamma }^1`$ (of course, in order to have this choice one needs $`nd2`$). Namely, the Euclidean theory now has a supersymmetric version for negative curvature ($`\mathrm{\Gamma }=\mathrm{\Gamma }^1`$) and another supersymmetric version for positive curvature, when $`\mathrm{\Gamma }=\mathrm{\Gamma }^0`$.
Mutatis mutandis one can draw the same conclusions for the theories of section 3.3. The mass terms require the same treatment as before, and the only novelty is the Chern-Simons-like cubic coupling for the scalar field. If one chooses $`\mathrm{\Gamma }=\mathrm{\Gamma }^{123}`$, no further explanation is required. On the other hand, if one chooses, say, $`\mathrm{\Gamma }=\mathrm{\Gamma }^{012}`$, then one obviously has to take into account the minus sign implicit in using $`\varphi ^a=\eta ^{ab}\varphi _b`$. Thus explicitly the Chern-Simons term reads
$$\frac{1}{3!}ϵ_{abc}\varphi ^a[\varphi ^b,\varphi ^c]=\varphi _0[\varphi _1,\varphi _2].$$
(47)
The only thing worth noting here is perhaps that, unlike an ordinary Chern-Simons term, which contains a first order time derivative, this algebraic term remains real in Euclidean signature. The payoff from using $`\mathrm{\Gamma }^{012}`$ is that this theory exists on manifolds of positive curvature (admitting solutions of the corresponding Killing spinor equation, of course). Thus we have essentially the following two Lagrangians:
$`\mathrm{\Gamma }=\mathrm{\Gamma }^{123}`$ $`L=L_{ESYM}+4\alpha ^2[(n2)\eta ^{mn}\varphi _m\varphi _n+(n4)\delta ^{ab}\varphi _a\varphi _b]`$
$`(n4)\alpha \left[\overline{\mathrm{\Psi }}\mathrm{\Gamma }_{123}\mathrm{\Psi }8\varphi _1[\varphi _2,\varphi _3]\right]`$
$`\mathrm{\Gamma }=\mathrm{\Gamma }^{012}`$ $`L=L_{ESYM}4\alpha ^2[(n2)\eta ^{mn}\varphi _m\varphi _n+(n4)\eta ^{ab}\varphi _a\varphi _b]`$ (48)
$`+(n4)\alpha \left[\overline{\mathrm{\Psi }}\mathrm{\Gamma }_{012}\mathrm{\Psi }8\varphi _0[\varphi _1,\varphi _2]\right].`$
We see that whereas in the pseudo-Riemannian case we had the freedom to choose either positive or negative curvature space-times only for $`n=4`$, in the Riemannian case the theories have this property for all $`n`$, subject to the restrictions $`nd2`$ for the A theories and $`nd4`$ for the B theories. In $`d1`$ (respectively $`d3`$) dimensions, there is no choice, $`\mathrm{\Gamma }`$ is dictated by whether one makes a purely spaceklike or a (space-)time-reduction.
## 4 Aspects of the Supersymmetry Algebra
In order to gain some insight into the structure of the theories introduced above, and to attempt to understand them from the (A)dS superalgebra point of view, in the following we will now (partially) calculate the supersymmetry algebras in these models.
### 4.1 The Superalgebra for Family A
Using (30), it is straightforward to calculate the commutator of two supersymmetry transformations $`\delta _i`$, associated with Killing spinors $`\epsilon _1,\epsilon _2`$ satisfying $`_\mu \epsilon _i=\alpha \mathrm{\Gamma }_\mu \mathrm{\Gamma }_1\epsilon _i`$, acting on the bosonic fields $`A_\mu `$ and $`\varphi _m`$. One finds
$`\frac{1}{4}[\delta _1,\delta _2]A_\mu `$ $`=`$ $`V^NF_{N\mu }+(n3)(\alpha +\alpha ^{})V_\mu \varphi _1+(\alpha \alpha ^{})V_{\mu i1}\varphi ^i`$
$`\frac{1}{4}[\delta _1,\delta _2]\varphi _1`$ $`=`$ $`V^NF_{N1}+(n3)(\alpha +\alpha ^{})V_1\varphi _1(\alpha +\alpha ^{})V_i\varphi ^i`$
$`\frac{1}{4}[\delta _1,\delta _2]\varphi _j`$ $`=`$ $`V^NF_{Nj}+(n3)(\alpha +\alpha ^{})V_j\varphi _1+(\alpha +\alpha ^{})V_1\varphi _j+(\alpha \alpha ^{})V_{ij1}\varphi ^i`$
Here we have introduced the notation
$`V_M`$ $`=`$ $`\overline{\epsilon }_1\mathrm{\Gamma }_M\epsilon _2\overline{\epsilon }_2\mathrm{\Gamma }_M\epsilon _1`$
$`V_{MNP}`$ $`=`$ $`\overline{\epsilon }_1\mathrm{\Gamma }_{MNP}\epsilon _2\overline{\epsilon }_2\mathrm{\Gamma }_{MNP}\epsilon _1.`$ (50)
Ordinarily, i.e. in Poincaré supersymmetry, one would just find the first term on the right hand side. Acting on the scalar fields, this is just the Lie derivative (diffeomorphism) with resepect to $`V^\mu `$ plus a field dependent gauge transformation,
$`V^NF_{Nm}`$ $`=`$ $`L_V\varphi _m+\delta _V\varphi _m`$
$`\delta _V\varphi _m`$ $`=`$ $`[V^NA_N,\varphi _m].`$ (51)
Here and in the following it should be understood that the $`V`$ in the Lie derivative refers only to the space-time components $`V^\mu `$ whereas all components $`V^M`$ enter in $`\delta _V`$.
The same is true for the gauge field provided that $`V_m`$ is constant, as is the case for parallel spinors. In that case, one has
$`_\mu V_m=0`$ $`V^NF_{N\mu }=L_VA_\mu +\delta _VA_\mu `$ (52)
$`\delta _VA_\mu =D_\mu (V^NA_N).`$
However, when the $`V_m`$ are not constant, then one has instaed
$$_\mu V_m0V^NF_{N\mu }=L_VA_\mu +\delta _VA_\mu +(_\mu V_m)\varphi ^m.$$
(53)
In order to understand how the right hand side of the supersymmetry algebra, including also all the other new terms, nevertheless manages to be an invariance of the Lagrangian in this case, we need to know some properties of the objects $`V_M`$ and $`V_{MNP}`$. The following identities are easily verified:
$`_\mu V_1`$ $`=`$ $`(\alpha +\alpha ^{})V_\mu `$
$`_\mu V_i`$ $`=`$ $`(\alpha \alpha ^{})V_{1i\mu }`$
$`_\mu V_{ij1}`$ $`=`$ $`(\alpha +\alpha ^{})V_{\mu ij}`$
$`_\mu V_\nu `$ $`=`$ $`(\alpha +\alpha ^{})g_{\mu \nu }V_1+(\alpha \alpha ^{})V_{1\nu \mu }.`$ (54)
In particular, therefore, $`V_\mu `$ is a Killing vector if $`\alpha ^{}=\alpha `$, and a conformal Killing vector (and a gradient vector) if $`\alpha ^{}=\alpha `$. In the former case, $`V_1`$ and the antisymmetric matrices $`V_{1ij}`$ are constant, whereas the other space-time scalars $`V_i`$ are not (and vice-versa for $`\alpha `$ real). Moreover, note that the above equations imply that for $`\alpha `$ real the function $`V_1^2+V_\mu V^\mu `$ is constant.
Using these results, we learn that the commutator of supersymmetry transformations on the gauge field can be written as
$$\frac{1}{4}[\delta _1,\delta _2]A_\mu =L_VA_\mu +\delta _VA_\mu (n4)(_\mu V_1)\varphi _1.$$
(55)
But since $`V_1`$ is constant for imaginary $`\alpha `$ and real $`\alpha `$ is only allowed for $`n=4`$, we see that in all cases the last term actually disappears and the commutator takes the standard form
$$\frac{1}{4}[\delta _1,\delta _2]A_\mu =L_VA_\mu +\delta _VA_\mu .$$
(56)
If $`\alpha `$ is imaginary, then the commutator on the scalars takes the form
$`\alpha ^{}=\alpha `$ $`\frac{1}{4}[\delta _1,\delta _2]\varphi _1=L_V\varphi _1+\delta _V\varphi _1`$ (57)
$`\frac{1}{4}[\delta _1,\delta _2]\varphi _j=L_V\varphi _j+\delta _V\varphi _j+2\alpha V_{ij1}\varphi ^i.`$
We see that in addition to diffeomorphisms (along a Killing vector) and gauge transformations, the algebra now also includes a rotation of the scalar fields by the constant matrix $`V_{ij1}`$ \- this is (a subgroup of) the R-symmetry algebra of the theory and, combined with an appropriate transformation of the fermions, a separate invariance of the Lagrangian. The appearance of the R-symmetry algebra in the commutator of supersymmetries is of course a well known feature of AdS superalgebras (for a recent review of AdS supersymmetry see ) which we have recovered here somewhat experimentally. Note that this extra rotation only appears for $`nd3`$. In particular, it is absent for $`n=4`$.
The case $`\alpha ^{}=\alpha `$ (and thus $`n=4`$) is a bit more complicated, but this should not be too surprising as now $`V_\mu `$ is only a conformal Killing vector,
$$L_Vg_{\mu \nu }=4\alpha V_1g_{\mu \nu },$$
(58)
and additional scale transformations of the scalars and fermions are required to produce an invariance of the Lagrangian density in that case. Recall that precisely when $`n=4`$ the scalar field action is conformally invariant so that this is feasible in principle.
The transformation on the gauge field is, as we have noted above, the standard one, which is fine since the Yang-Mills Lagrangian is conformally invariant precisely when $`n=4`$. The scalars now transform as
$`\alpha ^{}=+\alpha `$ $`\frac{1}{4}[\delta _1,\delta _2]\varphi _1=(L_V+2\alpha V_1)\varphi _1+\delta _V\varphi _1+\mathrm{\Delta }_V\varphi _1`$ (59)
$`\frac{1}{4}[\delta _1,\delta _2]\varphi _j=(L_V+2\alpha V_1)\varphi _j+\delta _V\varphi _j+\mathrm{\Delta }_V\varphi _j.`$
Here the modified Lie derivative $`L_V+2\alpha V_1`$ reflects the non-trivial conformal weight of the scalar fields, and
$`\mathrm{\Delta }_V\varphi _1`$ $`=`$ $`2\alpha V^i\varphi _i`$
$`\mathrm{\Delta }_V\varphi _j`$ $`=`$ $`2\alpha V_j\varphi _1`$ (60)
is a particular global (the $`V_i`$ are constant in this case) infinitesimal $`SO(d4)`$ rotation of the $`(d4)`$ scalar fields. This is only non-trivial for $`d=6`$ and for $`d=10`$. In the former case we find an $`SO(2)`$ rotation parametrized by $`2\alpha V_2`$, namely
$`\mathrm{\Delta }_V\varphi _1`$ $`=`$ $`2\alpha V_2\varphi _2`$
$`\mathrm{\Delta }_V\varphi _2`$ $`=`$ $`2\alpha V_2\varphi _1.`$ (61)
Note that in this case ($`\alpha `$ real) the bosonic generators of the algebra are conformal Killing vector fields that are also gradient vector fields (this is an extremely restrictive condition but solutions exist e.g. in de Sitter space). As a consequence, since the Lie bracket of two gradient vector fields is always zero, and also commutators of the modified operators $`L_V+2\alpha V_1`$ can be seen to vanish, the bosonic part of the algebra engendered in this way is Abelian, a situation apparently not covered by Nahm’s classification .
### 4.2 The Complete Superalgebra for $`n=4`$
Of course, to complete this discussion we should also calculate the commutator of two supersymmetry transformations on the fermions. At this point, because now Fierz identities are required, the discussion becomes somewhat dimension-dependent and we will only do this for $`n=4`$ which in many respects is the most interesting case to consider anyway.
For the $`(6,4)`$-theory, the supersymmetry variation of the spinor $`\mathrm{\Psi }`$ is
$`\delta \mathrm{\Psi }`$ $`=`$ $`\mathrm{\Gamma }^{MN}\epsilon F_{MN}4\alpha \varphi _m\mathrm{\Gamma }^m\mathrm{\Gamma }^1\epsilon `$ (62)
$`=`$ $`\mathrm{\Gamma }^{MN}\epsilon F_{MN}4\alpha (\varphi _1+\varphi _2\mathrm{\Gamma }^{21})\epsilon .`$
It follows that
$`\delta _1\delta _2\mathrm{\Psi }`$ $`=`$ $`2D_M(\overline{\epsilon }_1\mathrm{\Gamma }_N\mathrm{\Psi }\overline{\mathrm{\Psi }}\mathrm{\Gamma }_N\epsilon _1)\mathrm{\Gamma }^{MN}\epsilon _2`$ (63)
$`4\alpha (\overline{\epsilon }_1\mathrm{\Gamma }_m\mathrm{\Psi }\overline{\mathrm{\Psi }}\mathrm{\Gamma }_m\epsilon _1)\mathrm{\Gamma }^m\mathrm{\Gamma }^1\epsilon _2.`$
By the usual Fierz identity for SYM theories, the terms involving $`\epsilon _1`$ and $`\epsilon _2`$ will drop out after taking commutators and we drop them henceforth. From the other terms we find, using the Killing spinor equation
$$_\mu \overline{\epsilon }=\alpha ^{}\overline{\epsilon }\mathrm{\Gamma }_1\mathrm{\Gamma }_\mu ,$$
(64)
that
$$\delta _1\delta _2\mathrm{\Psi }=2\alpha ^{}\overline{\epsilon }\mathrm{\Gamma }_1\mathrm{\Gamma }_\mu \mathrm{\Gamma }_N\mathrm{\Psi }\mathrm{\Gamma }^{\mu N}\epsilon _2+2\overline{\epsilon }_1\mathrm{\Gamma }_ND_M\mathrm{\Psi }\mathrm{\Gamma }^{MN}\epsilon _24\alpha \overline{\epsilon }_1\mathrm{\Gamma }_m\mathrm{\Psi }\mathrm{\Gamma }^m\mathrm{\Gamma }^1\epsilon _2.$$
(65)
Taking commutators and using the Fierz rearrangement formula for Weyl spinors $`\mathrm{\Psi }_k`$ of the same chirality in $`d`$ dimensions,
$$\overline{\mathrm{\Psi }}_1M\mathrm{\Psi }_2\overline{\mathrm{\Psi }}_3N\mathrm{\Psi }_4=2^{d/2}\underset{p=0}{\overset{n/2}{}}c_p\overline{\mathrm{\Psi }}_1\mathrm{\Gamma }^{[p]}\mathrm{\Psi }_4\overline{\mathrm{\Psi }}_3N\mathrm{\Gamma }_{[p]}M\mathrm{\Psi }_2,$$
(66)
(here a sum over the antisymmetrized products of $`p`$ gamma matrices is understood) with
$`c_p`$ $`=`$ $`(1)^{\left(\genfrac{}{}{0pt}{}{p}{2}\right)}{\displaystyle \frac{2}{p!}}p<n/2`$
$`c_{n/2}`$ $`=`$ $`(1)^{\left(\genfrac{}{}{0pt}{}{n/2}{2}\right)}{\displaystyle \frac{1}{(n/2)!}},`$ (67)
one obtains
$`[\delta _1,\delta _2]\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \frac{1}{8}}2{\displaystyle \underset{p}{}}c_p(\overline{\epsilon }_1\mathrm{\Gamma }^{[p]}\epsilon _2\overline{\epsilon }_2\mathrm{\Gamma }^{[p]}\epsilon _1)\times `$
$`\times `$ $`[\mathrm{\Gamma }^{MN}\mathrm{\Gamma }_{[p]}\mathrm{\Gamma }_ND_M\mathrm{\Psi }+\alpha ^{}\mathrm{\Gamma }^{\mu N}\mathrm{\Gamma }_{[p]}\mathrm{\Gamma }_1\mathrm{\Gamma }_\mu \mathrm{\Gamma }_N\mathrm{\Psi }2\alpha \mathrm{\Gamma }^m\mathrm{\Gamma }^1\mathrm{\Gamma }_{[p]}\mathrm{\Gamma }_m\mathrm{\Psi }]`$
Now evidently only $`p=1`$ and $`p=3`$ contribute to the sum (this follows e.g. from the discussion leading to (130)), giving rise to terms involving the vectors $`V^M`$ and antisymmetric tensors $`V^{MNP}`$ encountered before. Upon using the equation of motion $`\mathrm{\Gamma }^MD_M\mathrm{\Psi }=0`$, the first term will just give the standard contribution proportional to
$$V^MD_M\mathrm{\Psi }=V^\mu _\mu \mathrm{\Psi }+[V^NA_N,\mathrm{\Psi }].$$
(69)
This has almost the right structure to be of the form diffeomorphism plus gauge transformation we encountered for the bosonic fields. However, the (covariant) derivative on the spinor alone is not part of the invariance of the action, i.e. the fermioic kinetic term is not invariant under
$$\delta \mathrm{\Psi }=V^\mu _\mu \mathrm{\Psi }$$
(70)
even if $`V`$ is Killing. Rather, for (conformal) Killing vectors the Lie derivatives on the bosonic fields have to be supplemented by the Lie derivative of the spinor field defined by
$$L_V\mathrm{\Psi }=V^\mu _\mu \mathrm{\Psi }+\frac{1}{4}_\mu V_\nu \mathrm{\Gamma }^{\mu \nu }\mathrm{\Psi }.$$
(71)
Let us note here that in the present case the second term only contributes when $`V`$ is a Killing vector ($`\alpha `$ imaginary), because $`V`$ is not only a conformal Killing vector but also a gradient vector when $`\alpha `$ is real. This additional contribution to the covariant derivative arises from the $`p=3`$ contributions to the second and third terms in (LABEL:fierz) in the form
$$_\mu V_\nu \mathrm{\Gamma }^{\mu \nu }=(\alpha ^{}\alpha )V_{\mu \nu 1}\mathrm{\Gamma }^{\mu \nu }.$$
(72)
The other $`p=1`$ contributions give rise to new terms in the supersymmetry algebra. After an altogether not particularly fascinating calculation one finds
$`\frac{1}{4}[\delta _1,\delta _2]\mathrm{\Psi }`$ $`=`$ $`L_V\mathrm{\Psi }+\delta _V\mathrm{\Psi }`$ (73)
$`+{\displaystyle \frac{1}{2}}(\alpha +5\alpha ^{})V_1\mathrm{\Psi }+{\displaystyle \frac{1}{2}}(\alpha +\alpha ^{})V_i\mathrm{\Gamma }^{i1}\mathrm{\Psi },`$
where $`i1`$. Now let us take a look at this for $`\alpha `$ real and imaginary respectively. For $`\alpha `$ imaginary, the complete commutator algebra reads
$`\frac{1}{4}[\delta _1,\delta _2]A_\mu `$ $`=`$ $`L_VA_\mu +\delta _VA_\mu `$
$`\frac{1}{4}[\delta _1,\delta _2]\varphi _m`$ $`=`$ $`L_V\varphi _m+\delta _V\varphi _m`$
$`\frac{1}{4}[\delta _1,\delta _2]\mathrm{\Psi }`$ $`=`$ $`L_V\mathrm{\Psi }+\delta _V\mathrm{\Psi }2\alpha V_1\mathrm{\Psi }.`$ (74)
Thus the only term we find in addition to the Lie derivative along a Killing vector and a gauge transformation is a constant ($`V_1`$ is constant) phase rotation ($`\alpha `$ is imaginary) of the spinor. The latter is of course an invariance of the Dirac action - in fact it is the diagonal $`U(1)`$ subgroup of the $`SU(2)`$ R-symmetry of the six-dimensional Weyl action. It is nevertheless interesting that this additional phase transformation appears in the commutator algebra for non-zero curvature. Its appearance in the $`(6,5)`$ theory has been noted in .
For $`\alpha `$ real, as we had seen before, already the algebra on the bosonic fields is more complicated. In this case we have
$`\frac{1}{4}[\delta _1,\delta _2]A_\mu `$ $`=`$ $`L_VA_\mu +\delta _VA_\mu `$
$`\frac{1}{4}[\delta _1,\delta _2]\varphi _1`$ $`=`$ $`(L_V+2\alpha V_1)\varphi _1+\delta _V\varphi _1+\mathrm{\Delta }_V\varphi _1`$
$`\frac{1}{4}[\delta _1,\delta _2]\varphi _j`$ $`=`$ $`(L_V+2\alpha V_1)\varphi _j+\delta _V\varphi _j+\mathrm{\Delta }_V\varphi _j`$
$`\frac{1}{4}[\delta _1,\delta _2]\mathrm{\Psi }`$ $`=`$ $`(L_V+3\alpha V_1)\mathrm{\Psi }+\delta _V\mathrm{\Psi }+\alpha V_i\mathrm{\Gamma }^{i1}\mathrm{\Psi }.`$ (75)
Once again we see the modified Lie derivative on the spinor field (the factor of 3 reflecting the familiar conformal weight 3/2 of a spinor field). We also see the constant R-symmetry transformation
$$\mathrm{\Delta }_V\mathrm{\Psi }=\alpha V_i\mathrm{\Gamma }^{i1}\mathrm{\Psi }$$
(76)
accompanying the rotation $`\mathrm{\Delta }_V\varphi _m`$ of the scalar fields. It is now straightforward to check that this indeed constitutes an invariance of the action, as it should.
### 4.3 The Superalgebra for Family B
We will now calculate the action of the commutator of two supersymmetry transformations on the bosonic fields for the family of Lagrangians (41) with supersymmetry transformation (42). Instead of $`V_M`$ and $`V_{MNP}`$, this algebra will now contain in addition to the vector $`V_M`$ the rank five anti-symmetric tensor
$$V_{MNPQR}=\overline{\epsilon }_1\mathrm{\Gamma }_{MNPQR}\epsilon _2\overline{\epsilon }_2\mathrm{\Gamma }_{MNPQR}\epsilon _1.$$
(77)
A straightforward calculation gives
$`\frac{1}{4}[\delta _1,\delta _2]A_\mu `$ $`=`$ $`V^NF_{N\mu }2\alpha V_{123i\mu }\varphi ^i`$
$`\frac{1}{4}[\delta _1,\delta _2]\varphi _a`$ $`=`$ $`V^NF_{Na}+2(n3)\alpha \epsilon _{abc}\varphi ^bV^c`$
$`\frac{1}{4}[\delta _1,\delta _2]\varphi _i`$ $`=`$ $`V^NF_{Ni}+2\alpha V_{123ij}\varphi ^j.`$ (78)
To interpret this, we proceed as in the analysis of (LABEL:dd1). First of all we note the following properties:
$`_\mu V_a`$ $`=`$ $`0`$
$`_\mu V_i`$ $`=`$ $`2\alpha V_{123i\mu }`$
$`_\mu V_\nu `$ $`=`$ $`2\alpha V_{123\mu \nu }`$
$`_\mu V_{123ij}`$ $`=`$ $`0`$
$`_\mu V_{123}`$ $`=`$ $`0.`$ (79)
This shows that $`V_\mu `$ is a Killing vector and that the coefficients of the scalar field rotations are constants. There is an $`SO(3)`$ rotation acting on the three scalar fields $`\varphi _a`$ and an $`SO(dn3)`$ rotation on the remaining scalars $`\varphi _i`$. The last relation, which we will only need later, shows that $`V_{123}`$ is a constant, an imaginary constant to be precise.
Moreover, the second relation allows us to write, as before,
$`\frac{1}{4}[\delta _1,\delta _2]A_\mu `$ $`=`$ $`V^\nu F_{\nu \mu }V^mD_\mu \varphi _m(_\mu V_i)\varphi ^i`$ (80)
$`=`$ $`L_VA_\mu +\delta _VA_\mu ,`$
so that all in all we have
$`\frac{1}{4}[\delta _1,\delta _2]A_\mu `$ $`=`$ $`L_VA_\mu +\delta _VA_\mu `$
$`\frac{1}{4}[\delta _1,\delta _2]\varphi _a`$ $`=`$ $`L_V\varphi _a+\delta _V\varphi _a+\mathrm{\Delta }_V\varphi _a`$
$`\frac{1}{4}[\delta _1,\delta _2]\varphi _i`$ $`=`$ $`L_V\varphi _i+\delta _V\varphi _i+\mathrm{\Delta }_V\varphi _i,`$ (81)
where
$`\mathrm{\Delta }_V\varphi _a`$ $`=`$ $`2(n3)\alpha \epsilon _{abc}\varphi ^bV^c`$
$`\mathrm{\Delta }_V\varphi _i`$ $`=`$ $`2\alpha V_{123ij}\varphi ^j.`$ (82)
Let us consider two special cases of this. The first is the $`(6,3)`$ theory. In this case evidently the commutator algebra is just the standard algebra, in agreement with the fact that the supersymmetry transformations themselves are just the standard ones - see (44). However, we will see below that in spite of this the commutator algebra acting on $`\mathrm{\Psi }`$ is different.
The second is the $`(10,4)`$ theory, i.e. the curved space counterpart of $`N=4`$ SYM theory in four dimensions. In this case we see that the supersymmetry algebra exhibits an $`SO(3)\times SO(3)`$ R-symmetry. I.e. from the point of view of the Poincaré supersymmetric theory the presence of curvature has broken the R-symmetry down from $`SO(6)`$ to $`SO(4)SO(3)\times SO(3)`$. This is in perfect agreement with what a look at the AdS superalgebras would lead one to conclude. The relevant superalgebra is now not the superconformal $`USp(N=4|4)`$ with its $`SU(4)`$ R-symmetry but the AdS superalgebra
$$OSp(N=4|4)O(3,2)\times SO(4).$$
(83)
It is rather pleasing to note that in the present context this reduction of the R-symmetry group can be traced back directly to the fact that the relevant Killing spinor equation involves the object $`\mathrm{\Gamma }_{123}`$. This itself came from the requirement of having a hermitan fermionic mass term for spinors that started off as ten-dimensional Majorana-Weyl spinors.
### 4.4 The Complete Superalgebra for $`n=3`$
We have seen above that in the $`(6,3)`$ theory the supersymmetry transformations (44) and the supersymmetry commutator algebra on the bosonic fields (82) are just the usual ones, and one might suspect that this essentially forces the commutator algebra on the fermionic fields to be the standard one as well. However, this is not necessarily the case.
First of all we know that the standard derivative term in the algebra has to be promoted to the spinorial Lie derivative (71) along a Killing vector field.
Secondly, in calculating $`[\delta _1,\delta _2]\mathrm{\Psi }`$ one encounters derivatives of the spinor parameters and in this way the fact that the $`\epsilon _i`$ are Killing spinors rather than parallel spinors feeds itself into the algebra.
Thirdly, in calculating this algebra one makes use of the $`\mathrm{\Psi }`$-equations of motion. A look at the action (44) reveals that these are
$$\mathrm{\Gamma }^MD_M\mathrm{\Psi }+\alpha \mathrm{\Gamma }_{123}\mathrm{\Psi }=0,$$
(84)
and therefore no longer describe a massless spinor.
And indeed one finds a new term in the commutator algebra even in this case, where such a term is not required by similar terms in the bosonic algebra. Starting from
$$\frac{1}{4}[\delta _1,\delta _2]\mathrm{\Psi }=\frac{1}{16}\underset{p}{}c_pV^{[p]}[\mathrm{\Gamma }^{MN}\mathrm{\Gamma }_{[p]}\mathrm{\Gamma }_ND_M\mathrm{\Psi }\alpha \mathrm{\Gamma }^{\mu N}\mathrm{\Gamma }_{[p]}\mathrm{\Gamma }_{123}\mathrm{\Gamma }_\mu \mathrm{\Gamma }_N\mathrm{\Psi }],$$
(85)
one finds that the first term contributes
$$\text{1st term}=\alpha V^LD_L\mathrm{\Psi }+\frac{3}{8}\alpha V^N\mathrm{\Gamma }_N\mathrm{\Gamma }_{123}\mathrm{\Psi }+\frac{1}{96}\alpha V^{[3]}\mathrm{\Gamma }_{[3]}\mathrm{\Gamma }_{123}\mathrm{\Psi },$$
(86)
which does not look particularly encouraging. However, the second term gives rise to
$`\text{2nd term}=`$ $`{\displaystyle \frac{5}{8}}\alpha V^\mu \mathrm{\Gamma }_\mu \mathrm{\Gamma }_{123}\mathrm{\Psi }{\displaystyle \frac{3}{8}}\alpha V^a\mathrm{\Gamma }_a\mathrm{\Gamma }_{123}\mathrm{\Psi }`$ (87)
$`{\displaystyle \frac{3}{96}}\alpha V^{[3]}\mathrm{\Gamma }_{[3]}\mathrm{\Gamma }_{123}\mathrm{\Psi }{\displaystyle \frac{1}{96}}\alpha V^{[3]}ϵ^{abc}\mathrm{\Gamma }_a\mathrm{\Gamma }^\mu \mathrm{\Gamma }_{[3]}\mathrm{\Gamma }_\mu \mathrm{\Gamma }_{bc}\mathrm{\Psi }.`$
The ‘mixed’ three-index terms, i.e. those involving $`V^{\mu \nu a}`$ and $`V^{\mu ab}`$ cancel, while the other two, those involving $`V^{\mu \nu \lambda }`$ and $`V^{abc}`$, add up and (using the chirality of $`\mathrm{\Psi }`$) give rise to a single term proportional to $`V_{123}\mathrm{\Psi }`$. The net result is then
$$\frac{1}{4}[\delta _1,\delta _2]\mathrm{\Psi }=V^LD_L\mathrm{\Psi }+\alpha V^\mu \mathrm{\Gamma }_\mu \mathrm{\Gamma }_{123}\mathrm{\Psi }2\alpha V_{123}\mathrm{\Psi }.$$
(88)
The second term is the missing contribution for the spinorial Lie derivative (71) as can be seen by using (79) and calculating
$`{\displaystyle \frac{1}{4}}_\mu V_\nu \mathrm{\Gamma }^{\mu \nu }\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \frac{1}{2}}\alpha V_{123\mu \nu }\mathrm{\Gamma }^{\mu \nu }\mathrm{\Psi }`$ (89)
$`=`$ $`\alpha V^\mu \mathrm{\Gamma }_\mu \mathrm{\Gamma }_{123}\mathrm{\Psi },`$
where the second equality follows from the chirality of $`\mathrm{\Psi }`$. Thus finally we have
$$\frac{1}{4}[\delta _1,\delta _2]\mathrm{\Psi }=L_V\mathrm{\Psi }+\delta _V\mathrm{\Psi }2\alpha V_{123}\mathrm{\Psi },$$
(90)
and only the last term requires some conmment. As we have seen in (79), $`V_{123}`$ is constant and, in fact, $`(V_{123})^{}=V_{123}`$, so that $`V_{123}`$ is an imaginary constant. But then the Lagrangian (44) is obviously invariant under this phase rotation of the fermions. Once again, as in (74), we find that the Killing spinor supersymmetry algebra includes this phase rotation for $`\alpha 0`$, i.e. for curved spaces.
## 5 The Coulomb Branch: Some Sample Calculations
### 5.1 Family A: Absence of a Maximally Supersymmetric Coulomb Branch
Recall that the standard Poincaré supersymmetric SYM theory has the Lagrangian (11)
$`L_{SYM}`$ $`=`$ $`\frac{1}{2}F_{\mu \nu }F^{\mu \nu }D_\mu \varphi _mD^\mu \varphi ^m\frac{1}{2}[\varphi _m,\varphi _n][\varphi ^m,\varphi ^n]`$ (91)
$`+`$ $`\overline{\mathrm{\Psi }}\mathrm{\Gamma }^\mu D_\mu \mathrm{\Psi }+\overline{\mathrm{\Psi }}\mathrm{\Gamma }^m[\varphi _m,\mathrm{\Psi }]`$
and the fermionic supersymmetry transformation (12)
$$\delta \mathrm{\Psi }=\mathrm{\Gamma }^{\mu \nu }\epsilon F_{\mu \nu }+2\mathrm{\Gamma }^{\mu m}\epsilon D_\mu \varphi _m+\mathrm{\Gamma }^{mn}\epsilon [\varphi _m,\varphi _n].$$
(92)
The quartic potential has flat directions for mutually commuting scalar fields. Thus there is a family of vacua parametrized by the constant expectation values of the scalar fields taking values in the Cartan subalgebra of the gauge group (modulo the action of the Weyl group). The supersymmetry transformations of the fermions are identically zero in such a background without any condition on $`\epsilon `$, and thus these configurations parametrize a family of maximally supersymmetric vacua of the SYM theory, the Coulomb branch.
We will now look for analogues of these solutions in the Killing SYM theories we have discussed above, and we will see that typically (because of the modified supersymmetry transformations and scalar potentials) there are no configurations which have all of the above properties, but that there are half-supersymmetric configurations which reduce to the above in the limit of vanishing curvature.
We begin by exploring the presence of a maximally supersymmetric purely scalar field configuration in the theories of section 3. We will first consider the Family A theories for $`d=6`$ and $`d=10`$. The supersymmetry variation of the fermions in a purely scalar background becomes
$$\delta \mathrm{\Psi }=2\mathrm{\Gamma }^{\mu n}\epsilon _\mu \varphi _n+\mathrm{\Gamma }^{mn}\epsilon [\varphi _m,\varphi _n]4\alpha [\underset{m=1}{\overset{dn}{}}\mathrm{\Gamma }^m\mathrm{\Gamma }^1\epsilon \varphi _m+(n4)\epsilon \varphi _1].$$
(93)
It is clear almost by inspection that, unless $`n=3`$ and without any further constraints on $`\epsilon `$ beyond the chirality constraint dictated by $`d`$-dimensional supersymmetry, vanishing of $`\delta \mathrm{\Psi }`$ implies vanishing of all the $`\varphi _m`$ because of the terms in $`\delta \mathrm{\Psi }`$ linear in the $`\varphi _m`$.
Indeed, first of all vanishing of the terms proportional to $`\mathrm{\Gamma }^{\mu m}`$ requires $`_\mu \varphi _m=0`$. The term linear in $`\varphi _1`$, proportional to the identity matrix acting on $`\epsilon `$ has to vanish sepreately, so one has $`\varphi _1=0`$. The coefficient of $`\mathrm{\Gamma }^{k1}`$, $`k1`$, is proportional to $`[\varphi _k,\varphi _1]2\alpha \varphi _k=2\alpha \varphi _k`$, and therefore also all the other scalar fields have to vanish, $`\varphi _k=0`$.
An exception occurs for $`n=3`$, as $`\varphi _1`$ does then not appear in the term in brackets proportional to $`\alpha `$ and can therefore be chosen to be constant but otherwise unconstrained. By gauge invariance, this constant can be chosen to lie in the Cartan subalgebra of the gauge group.
Thus for $`n3`$ there are no non-trivial maximally supersymmetric purely scalar configurations (switching on any scalar vev breaks at least some fraction of the supersymmetry), while for $`n=3`$ there is (for $`G=SU(2)`$) a one-dimensional Coulomb ‘twig’.
If these gauge theories can be shown to arise as worldvolume theories of branes, this should have implications for the possibility (or lack thereof) to move them apart, and thus also for the question of existence of marginal bound states among these branes.
### 5.2 Family B: Existence of a Discrete Family of Maximally Supersymmetric Scalar Field Configurations
For the Family B theories, with their supersymmetry variation
$$\delta \mathrm{\Psi }=\mathrm{\Gamma }^{MN}\epsilon F_{MN}4\alpha [\underset{m=1}{\overset{dn}{}}\varphi _m\mathrm{\Gamma }^m+(n4)\underset{a=1}{\overset{3}{}}\varphi _a\mathrm{\Gamma }^a]\mathrm{\Gamma }^{123}\epsilon ,$$
(94)
the situation is somewhat different.
In particular, as we had seen in (44), $`\alpha `$ disappears altogether from the supersymmetry transformation rules for $`(d=6,n=3)`$. In that particular case, we therefore find the ‘normal’ Coulomb branch parametrized by the three constant commuting scalars. These solutions are also the only maximally supersymmetric critical points of the scalar cubic plus quartic potential.
For the reductions of the $`d=10`$ theories to $`n7`$ dimensions the situation is the following. We once again set the gauge fields to zero. Then imposing $`\delta \mathrm{\Psi }=0`$ forces the scalars to be constants. The terms linear in the $`\varphi _k`$, $`k1,2,3`$ are proportional to $`\mathrm{\Gamma }^{k123}\epsilon `$ and have to vanish seperately. Thus $`\varphi _k=0`$. For the remaining scalar fields $`\varphi _a`$, by looking at the coefficients of $`\mathrm{\Gamma }^{ab}\epsilon `$ we find the condition
$$[\varphi _a,\varphi _b]=2\alpha (n3)ϵ_{abc}\varphi _c.$$
(95)
Up to an irrelevant scaling, this amounts to a homomorphism of the Lie algebra of $`SU(2)`$ into that of the gauge group $`G`$ and hence there are maximally supersymmetric vacua for each conjugacy class of such homomorphisms.
It can also be checked directly that this gives a critical point of the potential (with $`\varphi _m=0`$ for $`m1,2,3`$)
$$V(\varphi )=\frac{1}{2}\text{Tr}[\varphi _a,\varphi _b]^2+8\alpha ^2(n3)\text{Tr}\varphi _a^2+\frac{4}{3}\alpha (n4)ϵ_{abc}\varphi _a[\varphi _b,\varphi _c].$$
(96)
This is reminiscent of the analysis by Vafa and Witten of the vacua of the mass-perturbed $`N=4`$ SYM theory: in that case the cubic superpotential of the $`N=4`$ theory (in $`N=1`$ language) is perturbed by quadratic mass terms, and the equation for the critical points is equivalent to (95).<sup>4</sup><sup>4</sup>4For a recent discussion of these theories in the context of the AdS/CFT correspondence see . Here we find this solution even in the presence of an additional quartic term in the potential.
We see that for these theories there are indeed maximally supersymmetric vacua, but that their structure is rather different from that of the standard Coulomb branch. Instead of a continous we have a discrete family of vacua with unbroken supersymmetry, and this is reflected in the absence of flat directions in the scalar potential for $`n3`$.
### 5.3 Existence of a half-BPS Coulomb Branch for AdS Space-Times
In order to study configurations preserving some fraction of the supersymmetry, we need to know what kind of additional conditions can be imposed on an $`n`$-dimensional Killing spinor. Clearly a chirality condition (which is the natural condition for constant or parallel spinors) is incompatible with the Killing spinor equation
$$_\mu \eta =\alpha \gamma _\mu \eta .$$
(97)
Fortunately, very compact and explicit expressions are known for Killing spinors on AdS space-times, and these results will enable us to find half-supersymmetric scalar field configurations.
We begin by quickly reviewing the results obtained in . The $`AdS_n`$ metric takes a particularly simple form in horospheric (or the closely related Poincaré) coordinates, in which one has
$$ds^2=dr^2+\text{e}^{{\scriptscriptstyle \frac{2r}{\mathrm{}}}}\eta _{ij}dx^idx^j.$$
(98)
The scalar curvature of this metric is
$$R=\frac{1}{\mathrm{}^2}n(n1),$$
(99)
which identifies $`\mathrm{}`$ as the curvature radius of the space-time, related to our constant $`\alpha `$ by $`|\alpha |=1/2\mathrm{}`$. The spinorial covariant derivative in these coordinates is
$`_r\eta `$ $`=`$ $`_r\eta `$
$`_k\eta `$ $`=`$ $`_k\eta +{\displaystyle \frac{1}{2\mathrm{}}}\gamma _k\gamma _r\eta .`$ (100)
Hence the Killing spinor equation
$$_\mu \eta =\frac{1}{2\mathrm{}}\gamma _\mu \eta $$
(101)
can be written as the pair of equations
$`_r\eta `$ $`=`$ $`{\displaystyle \frac{1}{2\mathrm{}}}\gamma _r\eta `$
$`_k\eta `$ $`=`$ $`{\displaystyle \frac{1}{2\mathrm{}}}\gamma _k(1\gamma _r)\eta .`$ (102)
Clearly, if $`\gamma _r\eta =\eta `$, the solutions are
$$\eta ^+=\text{e}^{{\scriptscriptstyle \frac{r}{2\mathrm{}}}}\eta _0^+,$$
(103)
where $`\eta _0^+`$ is an arbitrary constant spinor satisfying
$$\gamma _r\eta _0^+=\eta _0^+.$$
(104)
These are the Killing spinors we will consider in the following. The general solution is
$$\eta =\text{e}^{{\scriptscriptstyle \frac{r}{2\mathrm{}}}\gamma _r}(1+\frac{1}{2\mathrm{}}x^k\gamma _{\underset{¯}{k}}(1\gamma _r))\eta _0,$$
(105)
where $`\eta _0`$ is now an arbitrary constant spinor and $`\gamma _{\underset{¯}{k}}`$ refers to an orthonormal basis. This shows that AdS has the maximal number of linearly independent Killing spinors, i.e. is maximally supersymmetric in the supergravity sense.
Armed with these solutions to the Killing spinor equations, we can now reconsider the issue of supersymmetric purely scalar field configurations. For concreteness we consider the Family A $`(6,n)`$ theories for $`n=4`$ and $`n=5`$.
For $`AdS_5`$ we choose gamma-matrices $`\gamma _k,k=0,1,2,3`$ satisfying
$$\{\gamma _k,\gamma _l\}=\text{e}^{{\scriptscriptstyle \frac{2r}{\mathrm{}}}}\eta _{kl}$$
(106)
and $`\gamma _r=\gamma ^{(5)}`$. A convenient basis for the $`d=6`$ Clifford algebra is then
$`\mathrm{\Gamma }_k`$ $`=`$ $`\sigma _1\gamma _kk=0,\mathrm{},n2=3`$
$`\mathrm{\Gamma }_r`$ $`=`$ $`\sigma _1\gamma ^{(5)}`$
$`\mathrm{\Gamma }_5`$ $`=`$ $`\sigma _2I`$ (107)
where we have now, for sanity’s sake, called the internal gamma matrix appearing in the Killing spinor equation
$$_\mu \epsilon =\alpha \mathrm{\Gamma }_\mu \mathrm{\Gamma }_5\epsilon ,$$
(108)
$`\mathrm{\Gamma }_5`$ instead of $`\mathrm{\Gamma }_1`$. For $`n=4`$ we will choose a dimensional reduction along the $`x^3`$-direction so that now $`\{\gamma _\mu \}=\{\gamma _k,\gamma ^{(5)}\}`$ with $`k=0,1,2`$.
For $`\epsilon `$ a six-dimensional Weyl spinor, $`\epsilon ^T=(\eta ^T,0)`$ the Killing spinor equation reduces to
$$_\mu \eta =i\alpha \gamma _\mu \eta $$
(109)
so we have the identification
$$i\alpha =\frac{1}{2\mathrm{}}.$$
(110)
Therefore the AdS Killing spinor equation becomes
$`_k\eta `$ $`=`$ $`{\displaystyle \frac{1}{2\mathrm{}}}\gamma _k\eta `$
$`_r\eta `$ $`=`$ $`{\displaystyle \frac{1}{2\mathrm{}}}\gamma ^{(5)}\eta ,`$ (111)
so that indeed $`\gamma _r=\gamma ^{(5)}`$ and the condition $`\gamma _r\eta =\eta `$ translates into a standard chirality condition in the four-dimensional sense.
We begin with the $`n=5`$ theory, denote the single scalar field simply by $`\varphi `$, and consider the fermionic variation (once again, we set the gauge fields to zero)
$$\delta \mathrm{\Psi }=2\mathrm{\Gamma }^{k5}\epsilon _k\varphi +2\mathrm{\Gamma }^{r5}\epsilon _r\varphi 8\alpha \varphi \epsilon .$$
(112)
Translating this into five-dimensional gamma matrices acting on $`\eta `$, one finds
$$\delta \mathrm{\Psi }=02i\gamma ^k\eta _k\varphi +2i\gamma _r\eta _r\varphi +\frac{4i}{\mathrm{}}\eta \varphi =0.$$
(113)
Now we find that for Killing spinors satisfying $`\gamma _r\eta =\eta `$, the supersymmetry condition becomes $`_k\varphi =0`$ and
$$_r\varphi =\frac{2}{\mathrm{}}\varphi ,$$
(114)
or
$$\varphi =\text{e}^{{\scriptscriptstyle \frac{2r}{\mathrm{}}}}\varphi ^0,$$
(115)
where $`\varphi ^0`$ is an arbitrary constant anti-hermitian matrix in the Lie algebra of the gauge group.
Let us note the following properties of this configuration:
1. By construction, this configuration leaves half of the supersymmetries (namely those associated with Killing spinors satisfying $`\gamma _r\eta =\eta `$) unbroken.
2. It is also a solution to the equations of motion. The equation of motion is (with the mass term expressed in terms of $`\mathrm{}`$)
$$\mathrm{}\varphi =\frac{4}{\mathrm{}^2}\varphi .$$
(116)
On functions depending only on $`r`$, this reduces to
$$(_r^2+\frac{4}{\mathrm{}}_r)\varphi =\frac{4}{\mathrm{}^2}\varphi ,$$
(117)
which is satisfied by $`\varphi \mathrm{exp}(2r/\mathrm{})`$.
3. In the flat space limit $`\mathrm{}\mathrm{}`$, $`\varphi `$ just reduces to a constant. In that limit there is a supersymmetry enhancement and $`\varphi ^0`$ parametrizes the maximally supersymmetric Coulomb branch of the five-dimensional $`N=2`$ theory.
For $`n=4`$ the situation is quite similar. We now have two scalar fields which, with the above conventions, would most naturally be called $`\varphi _3`$ (say) and $`\varphi _5`$. But I will just call them $`\varphi _{1,2}`$. Vanishing of the supersymmetry transformation in this case (for the $`\gamma _r=+1`$ Killing spinors) forces these fields to be $`x^k`$-independent and to commute, and the $`r`$-dependence is determined by
$$_r\varphi _{1,2}=\frac{1}{\mathrm{}}\varphi _{1,2},$$
(118)
leading to
$$\varphi _{1,2}=\text{e}^{\scriptscriptstyle \frac{r}{\mathrm{}}}\varphi _{1,2}^0.$$
(119)
These are once again half-supersymmetric solutions to the equations of motion, which in this case read
$$(_r^2+\frac{3}{\mathrm{}}_r)\varphi _{1,2}=\frac{2}{\mathrm{}^2}\varphi _{1,2},$$
(120)
and tend to the standard Coulomb branch of $`N=2`$ $`n=4`$ SYM as $`\mathrm{}\mathrm{}`$. Once again in that limit one finds a supersymmetry enhancement.
## 6 Open Issues: Interpretation and Applications
Above we have constructed two families of curved space counterparts of the standard Poincaré supersymmetric SYM theories which are globally supersymmetric on manifolds admitting Killing spinors, and we also began a preliminary investigation of their properties. But clearly a large number of issues still remain to be understood.
1. Foremost among them is perhaps the relevance of these theories to the dynamics of D-branes. For this one might also want to consider spacetimes of the form $`M=\mathrm{\Sigma }\times R`$ where $`\mathrm{\Sigma }`$ admits Killing spinors. The analysis closely resembles the one for Euclidean theories on $`\mathrm{\Sigma }`$ described in section 3.4.
If these theories play a role in that context, what are the consequences of the unusual properties of the Coulomb branch we have found in section 5? Where would one expect the mass or cubic potential terms to show up in applications? What about BPS configurations with non-trivial gauge fields (monopoles) in these theories? What is the relation to the BPS configurations in AdS space studied e.g. in ? What is the relation to the AdS calibrations of ? Are there interesting cohomological versions of these theories?
2. One might also want a better understanding of the superalgebras underlying these theories, depending on the number of available Killing spinors. What about the $`(d=10,n=8,9)`$ theories? How is the problem to construct such theories related to the absence of conventional AdS superalgebras beyond $`n=7`$? What about central charges and the addition of matter fields?
3. It would also be desirable to have a more conceptual understanding of the existence of these two classes of theories. For the Family A theories a possible approach may be the following. There is a one-to-one correspondence between (Riemannian, positive) Killing spinors on $`M`$ and parallel spinors on the so-called cone $`CM`$ over $`M`$ (see e.g. for a survey of these matters in the AdS/CFT context), with similar results for other signatures and signs. Thus the parallel spinors on $`CM`$ appear to play a dual role. On the one hand, they assure the supersymmetry of SYM theory on $`CM`$. On the other hand, they are invoked to establish the existence of Killing spinors on $`M`$ and hence supersymmetry of SYM theory on $`M`$. It is therefore natural to wonder if these two appearances of parallel spinors are related and if, indeed, a straightforward dimensional reduction of the supersymmetric theory on $`CM`$ might not have been a less roundabout way of arriving at the theory on $`M`$.
The problem with a naive dimensional reduction of a theory on $`CM`$ to one on $`M`$ is that there is no isometry in the cone direction but only a homothety. This suggests that perhaps one way to reduce a theory on $`CM`$ to a theory on $`M`$ is to perform a Scherk-Schwarz like reduction or gauging along the radial direction. The structure of the Family A theories is certainly suggestive: one ‘internal’ gamma-matrix $`\mathrm{\Gamma }_1`$ is singled out, which should be identified with $`\mathrm{\Gamma }_r`$, and the mass terms could arise from a Scherk-Schwarz like reduction. However, so far I have been unable to derive these theories in this way.
4. For the theories in Family B, an altogether different idea appears to be required to account for the Chern-Simons-like terms. The appearance of such a term in the $`n`$-dimensional action suggests an $`(n+3)`$-dimensional origin with a true CS term living in those extra three dimensions. Thus one should have a coupling
$$F^{(n)}(AdA+\mathrm{})$$
where $`F^{(n)}`$ is proportional to the volume form on $`M`$. Thinking of this as a RR field strength, one recognizes the Wess-Zumino coupling of a $`D(n+2)`$-brane world volume to a $`D(n2)`$-brane via the instanton action $`\text{Tr}FF`$. E.g. for $`n=5`$ and $`AdS_5`$ one has a $`D3D7`$ brane system. And indeed in the near-horizon limit of such a system one obtains $`AdS_5\times X_5`$, where $`X_5=S^5/Z_2`$ has a fixed $`S^3`$ over which the $`D7`$-branes are wrapped and the $`F^{(5)}`$ is proportional to the volume element on $`AdS_5`$ (plus its Hodge dual). Thus the $`D7`$-$`O7`$ couplings of the form
$$C^{(4)}\text{Tr}FF$$
could be responsible for the Chern-Simons like terms in the five-dimensional gauge theory obtained by reduction of the worldvolume theory of the $`D7`$-branes to $`AdS_5`$.
Of course, even if one can trace the Chern-Simons terms back to these configurations (and hence the corresponding supergravity theory), one still needs to understand why they are required by supersymmetry for a gauge theory on $`AdS_n`$ (or some other space-time admitting Killing spinors). However, perhaps the above considerations may at least provide a first step to such an understanding.
Alternatively, the existence of such terms in the action could be deduced from considerations as in , where D-brane actions in non-trivial antisymmetric tensor field backgrounds (and hence also non-trivial curvature by the Einstein equations) are studied.
## Appendix A Some Useful Identities for Fermion Bilinears
To understand the hermiticity properties of fermionic mass terms, which play an important role in the discussion of section 3, and in order to facilitate other manipulations, it is useful to know some identities for spinor bilinears involving gamma-matrices. First of all, let us introduce the unitary matrices $`A_\pm ,B_\pm ,C_\pm `$ by
$`\mathrm{\Gamma }_M^{}`$ $`=`$ $`\pm A_\pm \mathrm{\Gamma }_MA_\pm ^1`$
$`\mathrm{\Gamma }_M^{}`$ $`=`$ $`\pm B_\pm \mathrm{\Gamma }_MB_\pm ^1`$
$`\mathrm{\Gamma }_M^T`$ $`=`$ $`\pm C_\pm \mathrm{\Gamma }_MC_\pm ^1.`$ (121)
We can always choose $`A_{}=\mathrm{\Gamma }_0=A_{}^{}`$, and for $`d`$ even for $`A,B`$ and $`C`$ the $`\pm `$ matrices are related by multiplication by $`\mathrm{\Gamma }^{(d+1)}`$. For a general analysis see e.g. .
Majorana spinors are characterized by the condition
$$\mathrm{\Psi }^{}=B_\pm \mathrm{\Psi },$$
(122)
which is consistent provided that
$$B_\pm ^{}B_\pm =I.$$
(123)
Then for a Majorana spinor one has
$$\overline{\mathrm{\Psi }}=\mathrm{\Psi }^{}A_{}=\mathrm{\Psi }^TB_\pm ^TA_{}.$$
(124)
But one can easily check that, given the properties of $`A`$ and $`B`$, one has
$$B_\pm ^TA_{}\mathrm{\Gamma }_M(B_\pm ^TA_{})^1=\mathrm{\Gamma }_M^T,$$
(125)
and thus one can identify
$$C_{}=B_\pm ^TA_{}.$$
(126)
Hence the Majorana condition can also be written as
$$\overline{\mathrm{\Psi }}=\mathrm{\Psi }^TC_{},$$
(127)
which is perhaps more familiar. For the Majorana(-Weyl) theories in $`d=3+1`$ and $`d=9+1`$, we will usually choose $`B=B_+`$ to obtain
$$B=B_+\overline{\mathrm{\Psi }}=\mathrm{\Psi }^TC_{}.$$
(128)
In a Majorana basis of real gamma-matrices, one can always choose $`B_+=I`$ and $`A_{}=C_{}`$, since $`\mathrm{\Gamma }_M^{}=\mathrm{\Gamma }_M^T`$ and hence Majorana spinors are real in such a basis.
Now let us look quite generally at a spinor bilinear
$$\overline{\mathrm{\Psi }}\mathrm{\Gamma }^{[p]}\mathrm{\Phi }.$$
(129)
If $`\mathrm{\Psi }`$ and $`\mathrm{\Phi }`$ are chiral spinors, then it is easy to see that this bilinear is zero if $`p`$ is even and $`\mathrm{\Psi }`$ and $`\mathrm{\Phi }`$ have the same chirality (and likewise is zero if $`p`$ is odd and $`\mathrm{\Psi }`$ and $`\mathrm{\Phi }`$ have opposite chiralities). To see this one can calculate, using $`\mathrm{\Gamma }^{(d+1)}=\mathrm{\Gamma }^{(d+1)}{}_{}{}^{1}=\mathrm{\Gamma }^{(d+1)}`$,
$$\overline{\mathrm{\Gamma }^{(d+1)}\mathrm{\Psi }}\mathrm{\Gamma }^{[p]}\mathrm{\Phi }=(1)^{p+1}\overline{\mathrm{\Psi }}\mathrm{\Gamma }^{[p]}\mathrm{\Gamma }^{(d+1)}\mathrm{\Phi };,$$
(130)
from which the claim follows. Now let us check under which conditions the corresponding mass term is hermitian. To that end we calculate, noting an extra minus sign due to working with anticommuting spinors,
$$(\overline{\mathrm{\Psi }}\mathrm{\Gamma }^{[p]}\mathrm{\Phi })^{}=\eta _p\overline{\mathrm{\Phi }}\mathrm{\Gamma }^{[p]}\mathrm{\Psi }$$
(131)
($`\eta _p`$ was defined in (24)) so that $`\overline{\mathrm{\Psi }}\mathrm{\Gamma }^{[p]}\mathrm{\Psi }`$ is hermitian for $`\eta _p=+1`$, i.e. $`p=0,3,4,7,8\mathrm{}\mathrm{}`$ while for $`\eta _p=1`$, one has to multiply this term by $`i`$ to obtain a hermitian mass term.
For $`\mathrm{\Psi }`$ and $`\mathrm{\Phi }`$ Majorana, one has, using also $`C^T=C`$ (in a Majorana basis)
$$\overline{\mathrm{\Psi }}\mathrm{\Gamma }^{[p]}\mathrm{\Phi }=(\overline{\mathrm{\Psi }}\mathrm{\Gamma }^{[p]}\mathrm{\Phi })^T=\eta _p\overline{\mathrm{\Phi }}\mathrm{\Gamma }^{[p]}\mathrm{\Psi }$$
(132)
consistent with the fact that in a Majorana basis transposition and hermitian conjugation are the same operation. Thus the potential mass term $`\overline{\mathrm{\Psi }}\mathrm{\Gamma }^{[p]}\mathrm{\Psi }`$ is zero unless $`\eta _p=+1`$ (and in this case we are not permitted to render the mass term hermitian for $`\eta _p=1`$ by multiplying it by $`i`$).
Summarizing the above discussion, we see that for the $`d=2+1`$ Majorana theory, the only posibility is $`p=3`$, equivalent to $`p=0`$ because $`\mathrm{\Gamma }_{012}`$ is a multiple of the identity in that case. Likewise, for the $`d=3+1`$ Majorana theory, the only possibilities are $`p=0,3,4`$. For the chiral version of that theory, we have $`p=1`$ or $`p=3`$ (with imaginary and real coefficients respectively, related to the fact that $`\mathrm{\Gamma }^{(5)}`$ has a factor of $`i`$). For the chiral theory in $`d=5+1`$, one necessarily has $`p`$ odd, and therefore either $`p=1`$ (equivalent to $`p=5`$) with a factor of $`i`$, or $`p=3`$ with a real coefficient. We will find supersymmetric gauge theories for either choice of mass term. Finally, the only possibility for the Majorana-Weyl theory in $`d=9+1`$ is $`p=3`$.
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# Fermi Surface Properties of Low Concentration CexLa1-xB6: dHvA.
## I Introduction
The rare earth (RE) and divalent hexaborides have a variety of electrical, magnetic and thermodynamic properties and all have the same cubic structure. Among these materials are metallic LaB<sub>6</sub> , Kondo insulating SmB<sub>6</sub> , semi-metallic CaB<sub>6</sub>\[3 - 5\], heavy Fermion (HF) CeB<sub>6</sub>, and ferromagnetic EuB<sub>6</sub> . Extensive experimental and theoretical investigations have been done in order to understand their varying physical properties. One of the most decisive techniques to study the electronic properties of these materials is the de Haas - van Alphen (dHvA) effect with which the extremal cross-sectional areas of the Fermi surface (FS) and effective masses can be measured accurately. Pure LaB<sub>6</sub> and pure CeB<sub>6</sub> have been studied using this technique, having nearly identical prolate ellipsoidal FS’s, with the FS of CeB<sub>6</sub> being larger than that of LaB<sub>6</sub> by about 10% . For example, the values of the dHvA frequencies for LaB<sub>6</sub> and CeB<sub>6</sub> for the same minimum FS ellipsoid cross-section are 7.89kT and 8.66kT respectively. Yet, the effective masses are quite different being 0.65m<sub>e</sub> and 30m<sub>e</sub>(at 5-7 T) for LaB<sub>6</sub> and CeB<sub>6</sub> respectively, where m<sub>e</sub> is the free electron mass.
There have been several electrical, magnetic and thermal studies carried out to explore how this transition from light metallic LaB<sub>6</sub> to the HF CeB<sub>6</sub> takes place when La ions are gradually replaced by Ce ions that introduce 4f electrons into the metal. In addition to this, experimental work has been carried out, using the dHvA effect at high magnetic fields ($`>`$20 T), to explore the development of the HF behavior in Ce<sub>x</sub>La<sub>1-x</sub>B<sub>6</sub> . Here, it was reported that both the FS topology and effective masses transform continuously from that of pure LaB<sub>6</sub> to that of pure CeB<sub>6</sub> as the Ce concentration $`x`$ is increased from $`0`$ to $`1`$. Furthermore, beginning at very low values of $`x`$ (about 0.05), the contribution to the dHvA signal was observed to originate from only a single spin FS sheet.
Here we report detailed dHvA measurements, using both the field modulation technique at intermediate fields (6-15 T) and cantilever torque measurements to 30 T to investigate how the spin polarization manifests itself in the topological changes of the FS, and changes in effective masses of Ce<sub>x</sub>La<sub>1-x</sub>B<sub>6</sub> alloys for 0 $``$ $`x`$ $``$ 0.05. The results of these measurements are then compared with the previous pulsed field measurements and found to be in excellent agreement. In this paper, the spin polarization of the FS is investigated both qualitatively and quantitatively, and it is found that the spin up component dominates the dHvA signal as the Ce concentration increases.
## II Experiment and Sample Preparation
Single crystals of Ce<sub>x</sub>La<sub>1-x</sub>B<sub>6</sub> with $`x`$ = 0, 0.01, 0.02, 0.03, 0.04, and 0.05 were grown in Al flux in the shape of rectangular parallelepipeds (1$`\times `$0.5$`\times `$2 mm) with each face along a axis of the cubic structure . Most of the dHvA magnetization measurements on these samples were made in a 0$``$18 T superconducting magnet using the field modulation technique where a small time oscillating field $`h_0`$ is superimposed on a steady state field. The sample was placed inside an astatic pair of pick up coils that was balanced to 1 part in 10<sup>4</sup> at zero field. The remaining imbalance signal voltage in the pick up pair due to the magnetization of the sample oscillating with changing magnetic field was detected using a lock-in amplifier.
Using field modulation, constant angle dHvA measurements were made in the field range 6-15 T with the field sweeping at a rate of 0.05 T/min and the field direction rotated within (100) crystal plane. From these dHvA oscillations one can determine the extremal cross-sectional areas and effective masses from the temperature dependence of signal amplitudes over the entire FS. The sample was further rotated continuously in a fixed field of 10 T to observe the detailed dependences of the FS cross sectional areas on angle. For these measurements, the sample rotator angle was calibrated so that the orientation of the sample was known to a precision of rotation of 0.1. The above measurements were made at five or six different temperatures depending on the sample in the temperature range of 1.4 K to 4.2 K with the sample immersed in a pumped He<sup>4</sup> bath. The sample temperature was measured using a calibrated Cernox thermometer and the vapor pressure of the bath. The field was calibrated with NMR.
To verify the reproducibility of these results, torque measurements on an $`x`$= 0.01 Ce sample were made to 30 T at the National High Magnetic Field Laboratory, Tallahassee, FL.
## III Experimental Results
Figure 1 shows typical dHvA oscillations for Ce<sub>x</sub>La<sub>1-x</sub>B<sub>6</sub> ($`x`$= 0.01) for $`\theta `$ = 0 (i.e., the direction) in the field range of 10$``$11 T and at a temperature of 1.4 K. A discrete Fourier transform (DFT) of the signal is shown on the same graph. For measurements on pure LaB<sub>6</sub>, the frequency of the minimum area or $`\alpha _3`$ orbit is found to be 7.894 $`\pm `$ 0.004 kT, which is in good agreement with the original measurements of Arko et al.. With this as a point of reference, complete angular dependent studies of the frequencies in all of the Ce concentrations were made. This study was made in order to check the assumption that the FS is represented by an ellipsoid of revolution, because this assumption had previously been used to calculate FS volumes . Both constant angle field sweeps and constant field angle sweeps were obtained. Example data from the constant field rotation measurements for $`x`$ = 0.01 at T = 1.4 K and H = 10 T is shown in the inset of Figure 2. The oscillations with angle are caused by the fact that the dHvA phase, $`2\pi F(\theta )/H,`$ changes by $`2\pi `$ for each complete oscillation as $`\theta `$ is varied. The angular variation of $`F`$ can be determined from these rotation measurements using the counting method first implemented by Halse . Using this technique, the angular variations of the minimum area or $`\alpha _3`$ and maximum area or $`\alpha _{12}`$ orbits were obtained. As a further check, field dependent measurements and DFTs were also made at several angles. An example of the complete data for $`x`$ = 0.01 is shown in Figure 2.
The effective masses of the different samples (0 $``$ $`x`$ $``$ 0.05) are extracted from the temperature dependences of the oscillation amplitudes. In comparison, the value of the effective mass of the $`\alpha _3`$ orbit for x = 0 or pure LaB<sub>6</sub> was found to be 0.66$`\pm `$0.03m<sub>e</sub> compared to the results of Arko et al. which give 0.65m<sub>e</sub>. Thus, the two values are the same to within experimental uncertainty.
## IV Discussion
It is well known that the cross-sectional area $`A`$ of the FS in atomic units (a.u.) is related to the dHvA frequency $`F`$ by the Onsager relation ,
$$A(a.u)=(\frac{2\pi e}{\mathrm{}c})F=2.673\times 10^5F$$
(1)
where $`F`$ is the dHvA frequency in tesla. From the measured values of the dHvA frequencies, the extremal cross - sectional areas of the $`\alpha _3`$ and $`\alpha _{1,2}`$ orbits for the field applied along the crystal axis of Ce concentrations between $`x`$ = 0 and $`x`$ = 0.05 can be calculated. As shown in Figure 3, both of the frequencies corresponding to the $`\alpha _3`$ and $`\alpha _{1,2}`$ orbits are observed to increase with $`x`$. From these two frequencies, the volume of the FS, and hence the number of charge carriers per unit volume can be evaluated assuming that the FS is an ellipsoid of revolution. The $`x`$ dependence of the carrier density, n, calculated in this manner is shown in the inset of Figure 3 .
It has been reported that both LaB<sub>6</sub> and CeB<sub>6</sub> have similar prolate electron ellipsoidal FS’s situated at the six X points of the cubic Brillouin zone(BZ) that overlap along the $`\mathrm{\Gamma }`$R symmetry axes . This situation is shown schematically in Figure 4. The minimum cross-sectional area of the ellipse corresponding to the $`\alpha _3`$-orbit can be measured directly by applying the magnetic field H along the axis, while the maximum area for the $`\alpha _{1,2}`$ orbit is only observed through magnetic breakdown (MB) for this same field orientation . MB through the necks also leads to a multitude of frequencies $`\alpha _{1,2}+n\rho `$ within the $`\mathrm{\Gamma }XM`$ plane, $`\rho `$ being the orbit associated with a small FS orbit inside the necks and n is an integer . The value of $`\rho `$ is $``$ 4 - 8 T , which is smaller than our experimental uncertainty, and therefore cannot be resolved at the fields used for our measurements.
The expected angular dependence of the dHvA frequency from an ellipsoid of revolution is :
$$\frac{1}{F(\theta )}=A\mathrm{cos}^2\theta +B\mathrm{sin}^2\theta +C\mathrm{sin}\theta \mathrm{cos}\theta $$
(2)
Here $`F`$($`\theta )`$ is the dHvA frequency and $`\theta `$ is the angle of rotation from the principal axis of the ellipsoid normal to the field direction. A fit of the data to this equation is a useful criterion for deciding whether or not the FS is an ellipsoid of revolution . We are interested here in the angular variations of the $`\alpha _3`$ and $`\alpha _{1,2}`$ orbits. If the axis is rotated through an angle relative to the field direction in the (100) plane, the cross-sectional area of the FS normal to the field direction corresponding to the semi-minor axis of the ellipse, or the $`\alpha _3`$ orbit, increases while the one corresponding to the semi-major axis or the $`\alpha _{1,2}`$ orbit decreases, and these two areas or frequencies become degenerate at 45 (see Figure 2). These two frequencies and their fit to Equation 2 are plotted on the same graph in Figure 2, showing excellent agreement between the expected angular dependence from ellipsoidal FS and the data. This same agreement is obtained for all of the samples with 0 $``$ $`x`$ $``$ 0.05 measured here. Thus, all of the measurements support the assumption used previously that the FS is an ellipsoid of revolution.
Figure 5 shows the concentration ($`x`$) dependence of m\* for the $`\alpha _3`$ orbit and the curve fit to the data has the quadratic form m$`{}_{e}{}^{}(c+bx+ax^2)`$. According to Gor’kov and Kim, a linear dependence of the specific heat coefficient $`\gamma `$ (proportional to m\*) and the magnetic susceptibility $`\chi `$ of Ce and U based alloys would be a signature of contributions from independent impurity centers . However, at larger values of concentration in a system of localized spins, the linear dependence on $`x`$ does not hold and Gor’kov and Kim calculated an additional $`x^2`$ correction term to both the specific heat and magnetic susceptibility, using the Fermi liquid formulation. Therefore, the very good fit of our data to a quadratic equation relating m\* and $`x`$ would indicate that impurity centers are coupled even at low Ce concentrations. This observation is consistent with the model of coupled Ce atoms giving rise to the antiquadrupolar state in CeB<sub>6</sub> . We would like to point out that the effective mass m\* could be spin-dependent. At this stage in the analysis, the effective mass m\* is the one determined from the temperature dependence of the dHvA signal which has contributions from spin up and spin down electrons. The spin dependence of m\* will be discussed later in Section V.
In the usual case when both spin states have the same mass, the magnetic field dependence of the dHvA amplitude can be used to determine the average Dingle temperature $`\overline{T}_D`$ for the two spin states (see Section V). The effect of finite relaxation time due to impurity or point defects is to broaden the Landau levels leading to a reduction in amplitude roughly equivalent to that which would be caused by a rise of temperature to $`\overline{T}_D.`$ The field dependence of the dHvA amplitude can be expressed as
$$A_p=\frac{C_pTH^nR_D}{\mathrm{sinh}(\alpha pT/H)}$$
(3)
where $`A_p`$ is the amplitude of the $`p^{th}`$ harmonic, $`R_D`$ $`=exp(\alpha p\overline{T}_D/H)`$ is the Dingle reduction factor, $`p`$ is the harmonic number, $`\alpha `$ =$`14.69(m/m_e)T`$ $`/K`$, and $`C_p`$ and $`n`$ depend on the particular method of measurement. For the torque method the value of $`n=1/2`$ and a plot of $`ln(A_pH^{1/2}\mathrm{sinh}(\alpha pT/H))`$ versus $`1/H`$ yields a straight line with a slope of $`\alpha p\overline{T}_D`$ and a linear fit to the data gives $`m^{}\overline{T}_D`$. A Dingle plot for the high field cantilever data for 1%Ce in LaB<sub>6</sub> at T=1.73 K and in the field range $`1025`$ T is given in Figure 6. From the slope of the straight line the value of the average Dingle temperature $`\overline{T}_D`$ was found to be 3.5 K at high fields on the assumption that m = 0.73m<sub>e</sub> for both spin states. We have analyzed the field dependence of the amplitude for all six samples (0 $``$ $`x`$ $``$ 0.05) using field modulation in the range 7 $``$ B $``$ 15 T and we found that $`m^{}\overline{T}_D`$ is the same within the measurement uncertainty. This means that the arbitrary substitution of La by Ce (or Ce by La) contributes little or nothing to the mean free path l of the dominant spin channel which is given by
$$\frac{l}{l_c}=\omega _c\tau $$
(4)
where l<sub>c</sub> is the cyclotron length for a particular orbit given by
$$l_c=\left(\frac{2\mathrm{}F}{eB^2}\right)^{1/2}.$$
(5)
The average scattering rate is related to the average Dingle temperature by the relation
$$\overline{T}_D=\frac{\mathrm{}}{2\pi k\tau }.$$
(6)
Thus, since $`m^{}\overline{T}_D`$ is independent of $`x`$, the mean free path l is also independent of $`x`$ for a given field range. This observation is in agreement with the results of Goodrich et. al. that the dominant source of scattering should originate from other forms of crystallographic imperfections not from the Ce or La impurities.
## V Spin Dependent Scattering (SDS)
One of the effects of an applied magnetic field is to lift the spin degeneracy of the energy levels and the contributions to the dHvA signal from the spin-up and spin-down electrons. In conventional metals, the effect of the Zeeman splitting is to reduce the amplitude by a spin reduction factor $`R_S`$ given by
$$R_S=\mathrm{cos}(\frac{1}{2}p\pi g\frac{m^{}}{m_e})\mathrm{cos}(p\pi S)$$
(7)
where $`g`$ is the spin-splitting factor and m<sub>e</sub> is the free electron mass. We had reported earlier that above 5% Ce in LaB<sub>6</sub>, the contribution to the dHvA amplitude originates from a single spin FS. One explanation of this observation is that scattering from spin fluctuations does not occur with equal strength for the two spin directions. There is a large negative magnetoresistance in Ce<sub>x</sub>La<sub>1-x</sub>B<sub>6</sub> alloys that can be explained by the suppression of spin fluctuation scattering. However, from magnetoresistance measurements one cannot determine if one or two spin states are contributing to the scattering. Part of the purpose of the present work was to study in detail how the single spin dHvA signal develops in CeB<sub>6</sub> from the two spin signal in pure LaB<sub>6</sub>.
In the presence of high magnetic fields, if there is only one spin contribution to the signal, then a plot of $`\mathrm{ln}(A_p/`$ $`p^{1/2})`$ against the harmonic number, $`p,`$ yields a straight line because the spin splitting reduction factor as is given by Equation 7 is no longer present. However, if there are contributions to the dHvA amplitude from spin-up and spin-down states, there is spin reduction of the amplitude and we observe non linear dependence of $`\mathrm{ln}(A_p/`$ $`p^{1/2})`$ on $`p`$ as shown in Figure 7 for 0 $``$ x $``$ 0.05 except for $`x`$ $``$ 0.05, the concentration at which only one spin component is observed.
For pure LaB<sub>6</sub>, the dHvA amplitudes associated with the spin-up and spin-down electrons are equal. However, if magnetic impurities are involved, we could expect spin dependent scattering and the amplitudes for spin-up and spin-down oscillations could be unequal corresponding to unequal Dingle temperatures. So, the signal amplitude measured with the field modulation technique, which has contributions from both the spin-up and spin-down components, has to be modified in order to account for differences in Dingle temperatures and masses between the two spin channels. In the first harmonic detection, this signal voltage is related to the oscillatory magnetization by
$$\stackrel{~}{V}(\zeta )=G\underset{p}{}\stackrel{~}{M}_pJ_1(p\mathrm{\Lambda })\mathrm{sin}(p\zeta +\theta _p),$$
(8)
where $`G`$ represents the system gain, $`\stackrel{~}{M}_p`$ is the magnetization due to the $`p^{th}`$ harmonic of the dHvA signal, $`J_1`$ is a Bessel function of order one, $`\mathrm{\Lambda }`$ = $`2\pi h/H^2,`$ $`h`$ is the modulation amplitude, $`\zeta `$ $`=2\pi F/H,`$ $`\sigma `$ $`=`$ $`\pm 1,`$ and $`\theta _p`$ is the phase. The magnetization $`\stackrel{~}{M}`$ can be written :
$$\stackrel{~}{M}=\underset{p=1}{\overset{\mathrm{}}{}}\underset{\sigma }{}C_pD^pE^{\sigma p}\mathrm{sin}(p\zeta +p\frac{\pi }{4}\sigma p\pi S)$$
(9)
where
$`D`$ $`=`$ $`\mathrm{exp}(Km^{}\overline{T}_D/H),`$ (10)
$`\overline{T}_D`$ $`=`$ $`(T_D^{}+T_D^{})/2`$ (11)
$`E`$ $`=`$ $`\mathrm{exp}(Km^{}(\delta T_D)/H)`$ (12)
$`\delta T_D`$ $`=`$ $`(T_D^{}T_D^{})/2`$ (13)
and
$$C_p=\frac{\nu TF}{(A^\mathrm{"}p\mathrm{})^{1/2}}\frac{1}{\mathrm{sinh}(pKm^{}T/H)}$$
(14)
where $`\nu `$ $`=1.304\times 10^5(Oe^{1/2}/K),`$ $`K`$ $`=`$ $`14.69(T/K),`$ $`F`$ is the dHvA frequency and $`A`$ is the extremal cross sectional area of the FS.
If the phase difference $`\varphi _p(=p\pi S)`$ between spin-up and spin-down oscillations is field dependent, the first two harmonics of the magnetization $`\stackrel{~}{M}`$ may also be written as
$`\stackrel{~}{M}_1`$ $`=`$ $`C_1\left[z\mathrm{sin}(\psi +{\displaystyle \frac{\varphi }{2}})+z^{}\mathrm{sin}(\psi {\displaystyle \frac{\varphi }{2}})\right]`$ (15)
$`=`$ $`C_1(z^2+z^2+2zz^{}\mathrm{cos}\varphi )^{1/2}\mathrm{sin}(\psi +\theta _1)`$ (16)
and
$$\stackrel{~}{M}_2=C_2(z^4+z^4+2z^2z^2\mathrm{cos}(2\varphi ))^{1/2}\mathrm{sin}(2\psi \frac{\pi }{4}+\theta _2)$$
(17)
where $`z`$ is the Dingle reduction factor for the spin up electrons, $`z^{}`$ is the Dingle reduction factor for the spin down electrons, and $`\psi `$ is defined to be $`(2\pi F/H)\pm \pi /4`$ where the upper sign is for a minimum FS area and the lower is for a maximum FS area. Other higher harmonics can be written in a similar way. The relative phase between the spin up and down components of the signal is given by
$$\mathrm{tan}\theta p=\frac{z^pz^p}{z^p+z^p}\mathrm{tan}(\frac{p\varphi }{2})$$
(18)
and the spin up and the spin down Dingle reduction factors for the $`p^{th}`$ harmonic are given by
$$z^p=\mathrm{exp}(p\alpha T_D^{}/H)$$
(19)
and
$$z^p=\mathrm{exp}(p\alpha T_D^{}/H)$$
(20)
The relative phases are obtained by fitting the data to the first three harmonics of the magnetization. From the measured signal harmonic amplitude ratios $`M_2/M_1,`$ $`M_3/M_1`$ and the relative phases between the harmonics ($`\theta _2`$ $`2\theta _1)`$ and ($`\theta _33\theta _1)`$ , we calculate, at a given H and T, the values of the amplitudes $`D`$ and $`E`$ (i.e., $`m^{}`$ $`\overline{T_D}`$ and $`m^{}\delta T_D)`$ , and the value of $`S`$. Once the value of $`S`$ is known, we can determine the amplitude ratio $`z^{}/z`$ (or $`(m^{}T_D^{}m^{}T_D^{})`$ ). The average Dingle temperature, $`\overline{T_D},`$ is spin independent while the difference $`\delta T_D`$ = $`(T_D^{}T_D^{})`$ $`/2`$ shows the spin-dependence. As mentioned earlier, the effective mass m could also depend on spin. Therefore, since in all of the above expressions the product $`m^{}T_D`$ where both m\* and T<sub>D</sub> are spin-dependent always occurs, it is not trivial to single out the spin dependence of either one. Therefore, we will first assume that $`m^{}=m^{}`$ so that $`C_p`$ $`(`$ Eqn. 14) is the same for both spin states. The concentration dependence of $`m(T_D^{}`$ $``$ $`T_D^{}`$ $`)`$ is then calculated from the amplitude ratios $`M_2/M_1,`$ $`M_3/M_1`$ and relative phases. Next, we assume $`m^{}m^{}such`$ that $`C_p^{}`$ $`C_p^{}`$ and again calculate the x dependence, now of $`(m^{}T_D^{}m^{}T_D^{})`$. Figure 8 shows the concentration dependence of both $`m(T_D^{}`$ $``$ $`T_D^{}`$ $`)`$ and $`(m^{}T_D^{}m^{}T_D^{})`$ and there is clear evidence of SDS. In other words, if there is no SDS, then $`\delta T_D`$ $`=`$ $`T_D^{}`$ $`T_D^{}`$ $`=`$ $`0`$, that is the slope of the $`x`$ dependence of $`(m^{}T_D^{}m^{}T_D^{})`$ is non-zero. The circles represent the case that $`m^{}=m^{},`$ while the squares represent the case $`m^{}m^{}.`$ The slope of the line corresponding to $`m^{}m^{}`$ is approximately eight times that corresponding to $`m^{}=m^{}`$ indicating that $`m^{}T_D^{}`$ becomes greater than $`m^{}T_D^{}`$ as the Ce concentration increases, that additional increase in slope arising from the difference in mass.
For LaB<sub>6</sub>, the two spin components have equal amplitudes and the amplitude ratio $`z^{}/z`$ is equal to one or $`\delta T_D`$ $`=`$ $`T_D^{}`$ $`T_D^{}`$ $`=`$ $`0`$. In other words, the scattering rates for spin up and spin down are equal. As the Ce concentration is increased to 5%, the ratio of the spin-up to the spin-down amplitude increases confirming the complete observed spin polarization of the FS at 5%Ce in LaB<sub>6</sub> to one spin channel which is the spin up.
There are two contrasting theories concerning whether the FS polarizes to the spin up or spin down state. The first is the theory developed by Wasserman et al. for quantum oscillations in heavy fermion materials. This model, along with its zero field predecessor , is successful in accounting for the heavy effective masses as well as small topological changes in the FS caused by the presence of additional f electrons. However, this model also predicts that the dHvA signal is dominated by the spin down channel, and that its associated effective mass should decrease in a magnetic field. While apparent evidence for this was reported in very heavy compounds such as CeCu<sub>6</sub> , it is CeB<sub>6</sub> that shows perhaps the most dramatic mass changes with increasing magnetic field , but the polarity of the spin was not identified.
A large number of measurements have been performed on CeB<sub>6</sub> \[26-34\], which is regarded as a typical dense Kondo lattice with a very low Kondo temperature of $`12`$ K. Previous experimental data have appeared to be entirely consistent with the theoretical model of Wasserman et al. , that is, the effective mass is dramatically suppressed in a magnetic field . Recent measurements have shown that in addition to the suppression of the effective mass, there is notable deformation in the topology of the FS in a magnetic field, and this result is not entirely consistent with the mean field theory of Ref. . One aspect of the result that does appear to be consistent with this theoretical model, though, is that the dHvA signal originates from only a single spin FS sheet , even though the theory must be fundamentally incorrect because it predicts the wrong spin state to be observed.
All of the dHvA measurements on CeB<sub>6</sub> are made in the high magnetic field regime well above the metamagnetic transition where the dipole moments of the f electrons are essentially aligned. According to Edwards and Green , in this regime the theory developed by Wasserman et al. is no longer applicable. This is due to the fact that this is a mean field approach in which the interactions are assumed not to change in a magnetic field. Making such a description of the dHvA effect in HF systems is really only valid at low magnetic fields, that is, magnetic fields less than the Kondo temperature scale. Edwards and Green instead make the analogy of a HF compound in a magnetic field to a itinerant ferromagnet, in which spin fluctuations play a decisive role. Edwards and Green also anticipate that the dHvA effect should be dominated by only a single spin, but the up spin instead of the down spin. Therefore, one can see that our measurements are in agreement with the predictions of Edwards and Green that the down spin mass enhancement is larger than that of the up spin and does not contribute to the dHvA signal amplitude.
As further verification of this mass difference, we use Equations (19) and (20) to write
$$\frac{H}{K}\mathrm{ln}(z/z^{})=(m^{}m^{})T+(m^{}T_D^{}m^{}T_D^{})$$
(21)
The quantity on the left hand side of Equation 21 is calculated for each value of $`x`$ and plotted versus T in Figure 9. It can be seen that the data is linear in $`T`$ with a slope of $`(m^{}m^{})`$ and intercept $`(m^{}T_D^{}m^{}T_D^{})`$. The value of $`(m^{}m^{})`$ ranges from 0.003 for $`x`$ = 0 or pure LaB<sub>6</sub> to 0.09 for $`x`$ = 0.05 respectively. In addition, the value of $`(m^{}T_D^{}m^{}T_D^{})`$ ranges from 0.03 for $`x`$ = 0 to 0.3 for $`x`$ = 0.05. Thus, for pure LaB<sub>6</sub> , which is not spin polarized, $`m^{}=`$ $`m^{}`$ as expected. As the Ce concentration increases to 5%, $`(m^{}`$ $`m^{})`$ increases with $`m^{}`$ being greater than $`m^{}`$ by about 10%. If this mass difference continues to increase with $`x`$ the observed discrepancy between specific heat and dHvA mass measurements in CeB<sub>6</sub> is explained. Moreover, the difference in the scattering parameter,$`(m^{}T_D^{}m^{}T_D^{})`$ , increases with the Ce concentration confirming that the observed FS is due only to the spin up state. These observations lead us to the conclusion that it is the combination of $`\mathrm{\Delta }m^{}`$ and SDS that takes the down spin out of the dHvA signal. Therefore, both $`\mathrm{\Delta }m^{}`$ and SDS are equally important in understanding many of the properties of CeB<sub>6</sub> .
## VI Conclusion
We have performed a detailed microscopic dHvA study of low concentration (up to x = 0.05) in Ce<sub>x</sub>La<sub>1-x</sub>B<sub>6</sub> alloys and determined the development of the size and geometry of the FS from that of spin unpolarized LaB<sub>6</sub> to the spin polarized Ce<sub>x</sub>La<sub>1-x</sub>B<sub>6</sub> $`(x`$ $`0.05)`$ alloys. We have shown:
* The spin up signal amplitude dominates the dHvA signal when the Ce concentration is $``$ 5%.
* We determined that the down spin mass is greater than that of the up spin and the spin down contribution to the dHvA signal amplitude is small.
* The angular dependence of the dHvA extremal areas of the FS show that the assumption that the FS is an ellipsoid of revolution is valid for all concentrations measured.
* The dependence of the effective mass on concentration is in agreement with the existing theories of magnetic impurity interactions .
Overall, this work is the most detailed study using dHvA of alloy systems involving magnetic ions with concentrations greater than 1% that has been reported.
Acknowledgment
A portion of this work was performed at the National High Magnetic Laboratory, which is supported by NSF Cooperative agreement No. DMR - 9527035 and by the State of Florida. Additional support from the NSF (DMR9971348) is acknowledged by Z. Fisk.
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FIGURE CAPTIONS
FIG. 1 An example of dHvA oscillations from the field modulation measurements for $`x`$ = 0.01 sample for H along the axis. The inset shows the DFT of the oscillations for a field range of 10 $``$ 11 T.
FIG. 2 The dHvA frequencies F<sub>3</sub> and F<sub>1,2</sub> as a function of orientation for $`x`$ = 0.01 from field sweep and rotation (high density points in the Fig.) measurements. The solid lines are the fits to the to the ellipsoid Eqn. 2. The inset shows the raw data from angular sweep measurements at 10 T.
FIG. 3 Ce concentration dependence of the dHvA frequencies F<sub>3</sub> and F<sub>1,2</sub> at fields $`10`$ $``$ $`11`$ T. The inset shows the number of electrons per unit volume, n, as a function of $`x.`$
FIG. 4 FS of LaB<sub>6</sub> .
FIG. 5 The Ce concentration dependence of cyclotron mass for the $`\alpha _3`$ orbit at 10 T. The solid line is a quadratic fit to the data.
FIG. 6 Dingle plot for $`x`$ = 0.01 at 1.73 K in the field range of 10 $``$ 25 T using the cantilever technique.
FIG. 7 A plot of $`\mathrm{ln}(A_p`$ $`/`$ $`p^{1/2})`$ versus the harmonic number $`p`$ for $`x`$ = 0 to 0.05. Note that it becomes linear at $`x`$ = 0.05.
FIG. 8 Ce concentration dependence of $`m^{}T_D^{}`$ $``$ $`m^{}T_D^{}`$ .
FIG. 9 Temperature dependence of ($`H/K)\mathrm{ln}(z/z^{})`$ for all the alloys including pure LaB$`_6.`$ From the linear fits to the data, mass differences between the two spin states were determined.
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# The Rapidly Fading Afterglow from the Gamma-Ray Burst of 1999 May 6The observations reported herein were obtained with the Very Large Array (VLA) operated by the National Radio Astronomy Observatory which is a facility of the National Science Foundation operated under a cooperative agreement by Associated Universities, Inc.; and with the W. M. Keck Observatory, which is operated by the California Association for Research in Astronomy, a scientific partnership among California Institute of Technology, the University of California and the National Aeronautics and Space Administration.
## 1 Introduction
GRB 990506 was detected by the BATSE (trigger No. 7549) on board the Compton Gamma-Ray Observatory (CGRO) on 1999 May 6.47 UT (Kippen (1999)). It was a bright burst, lasting approximately 150 s, with a fluence in the 50-300 keV band of 2.23$`\times 10^4`$ erg cm<sup>-2</sup> – placing it in the top 2% of BATSE bursts as ranked by fluence. The PCA instrument on board the Rossi X-Ray Timing Explorer (RXTE) began scanning a 8$`{}_{}{}^{}\times 6^{}`$ region centered on the BATSE position some three hours after the burst (Marshall & Takeshima (1999)) and discovered a previously uncataloged X-ray source, which faded significantly over the observing interval (Takeshima & Marshall (1999)). The initial 12′ error circle for the RXTE-PCA position of GRB 990506 was further refined to a 30 arcmin<sup>2</sup> region through the addition of timing annuli from the Interplanetary Network (IPN), involving the NEAR, Ulysses and CGRO satellites (Hurley et al. (1999)).
Despite the fact that at least nine different optical telescopes imaged the IPN/RXTE-PCA region containing the X-ray transient (some as early as 90 minutes after the burst), no corresponding optical transient was detected. In this paper we report the discovery of a radio transient (VLA J115450.1$``$2640.6) in this field, which we propose is the radio afterglow of GRB 990506. We also present upper limits on the brightness of the corresponding optical afterglow. We discuss the radio light curve of this transient, and its implications for the lack of detection of the optical afterglow. Finally, we present the discovery of the probable host galaxy of the radio transient, and by implication the GRB itself. GRB 990506 joins a small but growing class of bursts which evidently lack a bright optical transient but are seen at both X-ray and radio wavelengths.
## 2 Discovery of the Radio Afterglow
Very Large Array (VLA) observations were initiated 1.5 days after the $`\gamma `$-ray burst. Details of this and all subsequent VLA observations are given in Table 1. In order to image the entire 11.5′$`\times `$2.5′ error region with the VLA at 8.46 GHz two pointings were required, each with a field-of-view at half power response of 5.3′ in diameter. In the southern field we detect two sources whose flux density ($`150`$ $`\mu `$Jy) did not vary significantly from 1999 May 8 to 1999 May 21. There were two additional sources in the northern pointing, only one of which was within the IPN/RXTE-PCA error box. This source, hereafter VLA J115450.1$``$2640.6, was clearly detected during the first two epochs but had faded by at least a factor of three thirteen days later. VLA J115450.1$``$2640.6 was located at (epoch J2000) $`\alpha `$ = $`11^h54^m50.13^s`$ ($`\pm 0.02^s`$) $`\delta `$ = $`26^{}40^{}35.0^{\prime \prime }`$ ($`\pm 0.4^{\prime \prime }`$), where the errors are the quoted 1$`\sigma `$ uncertainties in the Gaussian fit obtained from the combined images on May 8 and May 9. The source has not been detected since at any frequency (see Fig. 1). A search was made for polarized emission by imaging the combined datasets on May 8 and 9 in all four Stokes parameters. The 3$`\sigma `$ upper limits were 30% for both linear and circular polarization.
Of immediate concern is the possibility that VLA J115450.1$``$2640.6 is a background variable radio source lying by chance within the moderately large IPN/RXTE-PCA error box. On the basis of the radio data alone we cannot eliminate this possibility altogether. However, the order of magnitude change in source strength over the two weeks immediately following the GRB, and the failure to detect any subsequent radio emission after monitoring for 375 days, strongly suggests that VLA J115450.1$``$2640.6 was the afterglow of GRB 990506. We now proceed with this hypothesis.
Reaching a flux density of 580 $`\mu `$Jy, GRB 990506 is in the top half of the dozen radio radio afterglows detected to date (ranked by flux density), fainter only than GRB 970508, GRB 980703, GRB 991208, and GRB 991216. There is some indication that the flux density increased from $`\mathrm{\Delta }t`$=1.66 to 2.70 days (see Table 1) but the change in flux $`\mathrm{\Delta }`$S=134$`\pm `$68 $`\mu `$Jy is only marginally significant. Some time between 2.5 days after the burst and 15.5 days, the radio afterglow declined below a detectable level. If we describe this decline in flux density, $`F_R`$, as a power-law in time (i.e., $`F_Rt^{\alpha _d}`$) then the minimum slope of the decline depends on the poorly determined time, $`t_m`$, for the onset of the decline such that $`\alpha _d^1\mathrm{log}t_m\mathrm{log}15.66`$.
## 3 Limits on the Optical Afterglow, and Discovery of the Host Galaxy
We obtained our first optical image of the field of GRB 990506 using the 60-inch CCD Camera (160 arcmin<sup>2</sup> field of view) approximately 18 hr after the GRB and continued for the next four nights. More details of the R-band optical observations to date are presented in Table 2. No confirmed optically variable sources in the IPN error box were found by our group and others (e.g., Henden et al. 1999, Pedersen et al. 1999a ).
At the position of the radio transient, an upper limit as strict as $`R>23.5`$ (Pedersen et al. 1999b) is placed on the brightness of any point-source optical afterglow starting 12 hours after the burst. Table 2 summarizes the $`R`$-band optical limits; our $`V`$-band observations on the first fours days also revealed no counterpart to less stringent flux densities.
On 11 June 1999 UT, 35.78 days after the GRB, we re-observed the field of GRB 990506 at the position of the radio transient using the Low Resolution Imaging Spectrometer (LRIS; Oke et al. 1995) on the Keck II 10-meter Telescope on Mauna Kea, Hawaii. Six images of the field totaling 26 minutes of integration were obtained in the Cousins $`R`$-band. Observations of the standard field PG 1323-086 (Landolt (1992)) were used for magnitude zero-point calibration, which we also checked against the photometry by Vrba et al. (1999). An astrometric plate solution was obtained relative to the USNO A2.0 catalog (Monet et al. 1998) with a statistical error of 0.24, 0.27 arcsec ($`\alpha `$, $`\delta `$).
Figure 2 depicts the position of the radio transient and persistent 1.4 GHz emission overlaid atop the optical Keck image. Coincident with the position of the fading radio source is a faint extended (NE-SW) galaxy with an irregular morphology. The putative host appears as two knots of roughly equal brightness, with the radio transient position lying within the southwest knot. We find $`R=24.0\pm 0.3`$ and $`R=24.4\pm 0.3`$ mag for the northeast and southwest knots, respectively, based on photometry from the USNO (Vrba et al. 1999). The error in these magnitudes reflects the uncertainty in the color of the galaxy and aperture correction. The morphology of the host is suggestive of a merging or interacting galaxy pair.
## 4 Field radio Galaxies
In a deep radio image taken at 1.4 GHz on 1999 June 12, two weak background radio sources are seen within 30<sup>′′</sup> of VLA J115450.1$``$2640.6 (Fig. 2). The stronger of the two sources (R1) is 23<sup>′′</sup> to the SE and has a flux density of 540 $`\mu `$Jy at 1.4 GHz. This source is undetected at 8.46 GHz, requiring either that it is resolved out at higher resolution or that it has a moderately steep spectral index, $`\alpha <1.1`$, where $`S\nu ^\alpha `$. The weaker of the two nearby sources (R2) is detected some 30<sup>′′</sup> to the NW of VLA J115450.1$``$2640.6. This source has a flux density of 170 $`\mu `$Jy at 1.4 GHz and a spectral index steeper than $`0.5`$. From recent spectroscopic observations (see table 2) we identify the two radio sources (R1 and R2) with an early-type spiral and a QSO at redshifts $`z=0.326`$ and $`z=0.273`$, respectively.
## 5 Discussion
The simplest explanation for the behavior of the radio afterglow from GRB 990506 is that the emission originates in the forward shock driven into the surrounding medium by the relativistically expanding blast wave (Sari, Piran & Narayan (1998)). In this case the apparent rise in the 8.46 GHz light curve to a maximum $`F_m`$, followed by a decay is the result of an evolving synchrotron spectrum, produced from electrons accelerated in the shock, whose peak frequency $`\nu _m`$ passed through the band during our observing interval. As noted in §2, the exact time, $`t_m`$, for the onset of the decay and the value of the power-law slope, $`\alpha _d`$, are uncertain owing to the undersampling of the light curve. Given the data, a range of $`t_m`$ values from 1 to 5 days are possible. The upper limit on $`t_m`$ is determined by requiring that $`\alpha _d>2`$. Significantly steeper values have yet to be seen for other afterglows. Furthermore, this value is entirely consistent with the X-ray decay $`\alpha _d=1.9\pm 0.6`$ for this burst as measured by the RXTE-ASM (Takeshima & Marshall (1999)). For $`1t_m5`$ days the corresponding values of the decay slope are $`0.8\alpha _d2`$.
In Figure 1 we plot a 8.46 GHz light curve expected for a adiabatic forward shock, propagating in a homogeneous medium to demonstrate that such models are consistent with the data. In this representative model, the peak flux is 580 $`\mu `$Jy, $`t_m`$=2.5 days, and $`\alpha _d=1.25`$. The tightest constraint on $`\alpha _d`$ and/or $`t_m`$ comes from the 1.4 GHz flux density upper limit on 1999 June 12.92 UT. Reasonable model inputs predict peak flux densities at this frequency of order 100-200 $`\mu `$Jy. Therefore, at least this one epoch favors the steeper decay indices and the smaller $`t_m`$ values from the ranges given above.
The estimated $`t_m`$ value suggests that GRB 990506 entered the decay phase at 8.46 GHz considerably earlier than previous radio afterglows. At this same frequency the decay timescales for GRB 970508 (Frail, Waxman & Kulkarni (2000)), GRB 980329 (Frail et al. 2000a ), GRB 980519 (Frail et al. 2000b ), GRB 990510 (Harrison et al. (1999)), GRB 981216 (Frail et al. (1999)) range from 8.5 days to 90 days. Only GRB 990123 exhibited a faster decay (Kulkarni et al. (1999); see below).
One consequence of the early radio decay is that the apparent absence of the optical afterglow can now be understood rather simply without the need to invoke excess extinction along the line of sight. If the synchrotron peak moved through the radio band $`\nu _R`$ on the order of $`t_m`$=1-5 days, then the corresponding timescale for the peak to move through the optical bands $`\nu _o`$ is $`t_o=t_m(\nu _R/\nu _o)^{2/3}`$, or $`60t_m300`$ s (Mészáros & Rees (1997)). This small timescale can be contrasted with GRB 971214 for which Ramaprakash et al. (1998) derived $`t_o`$=0.6 days at $`\nu _m=3\times 10^{14}`$ Hz. In general, the flux density decays from its peak $`F_m`$ as a power law with index $`\alpha _d`$ for $`t>t_o`$. Provided that the synchrotron peak frequency $`F_m`$ remains constant (as it should if dealing with a spherical, adiabatic expansion), then the peak 8.4 GHz flux density of 580 $`\mu `$Jy predicts a peak optical afterglow emission of $`R=17`$ mag (after correcting for foreground extinction). Even a relatively shallow decay, $`\alpha _d=1.25`$, predicts that the OT was already undetectable at $`R=22`$ mag 90 minutes after the burst when the field of GRB 990506 was observed by Zhu & Zhang (1999). By the time of the deep $`R`$-band observations (e.g., Masetti et al. (1999), Pedersen et al. 1999b ) the OT had $`R>25`$ mag and was likewise undetectable.
Perhaps more importantly, the early radio decay is indicative of unusual physical conditions. Short-lived but moderately bright afterglows can result if the GRB exploded with a high energy ($`E_{}`$) into a low density medium ($`n_{}`$), or, as noted by Galama et al. (1999), the magnetic field in the forward shock was weak. Unfortunately, while we have constrained $`F_m`$ and $`\nu _m(t_m)`$, the paucity of broadband data for this burst does not allow us to determine other important observables such as the synchrotron self-absorption frequency $`\nu _a`$ and the cooling frequency $`\nu _c`$. Without this information it is not possible to fully constrain the physical properties of the afterglow and the surrounding medium (Wijers & Galama (1999)).
Alternatively, prompt radio emission, analogous to the well-known optical emission (Akerlof et al. (1999)), was also seen in the afterglow from GRB 990123 (Kulkarni et al. (1999)) and perhaps GRB 970828 as well. This radio emission arises in the reverse shock in the days following a $`\gamma `$-ray burst (Sari & Piran (1999)). It is suppressed at early times by synchrotron self-absorption, while at later times the light curve decays rapidly ($`t^2`$).
We note that it is possible that the southwest knot may itself be the optical transient since the GRB field has not been reimaged since 35 days after the burst. However, even a shallow decay in the optical of $`t^1`$ then predicts an $`R`$-band magnitude of 20.1 one day after the burst which was not observed. We thus find it unlikely that the southwest knot could be an optical transient given the early non-detections. Instead, we interpret the two knots as part of the host galaxy of GRB 990506. The apparent $`R`$-band magnitude of the putative host of GRB 990506 is similar to that of other GRB host galaxies. Spectroscopic observations are necessary in order to establish the redshift and other physical properties of this host galaxy. The bimodal morphology may be indicative of merger activity. In this respect the host appears similar to a growing number of GRB hosts with irregular morphology: GRB 980613 (Djorgovski et al. 1999), GRB 970828 (Djorgovski et al. 2000), and GRB 990123 (Bloom et al. 1999).
## 6 Conclusions
We have identified the radio afterglow from GRB 990506. This GRB was unusual in that it produced a radio afterglow that began to fade at very early times (between 1 and 5 days after the burst) and in that no optical afterglow was detected in spite of the numerous deep images obtained. Both these observations may be explained by a high energy spherical fireball expanding into a low density environment. In this simple picture there is no need to invoke dust extinction to account for the lack of detection of an optical afterglow. We cannot rule out, however, that the radio emission originated in a reverse shock. If the reverse shock produced the radio afterglow, then the emission from the forward shock (in both optical and radio) was presumably too faint to be seen. To find additional GRBs of this type will require more rapid followup to precise burst localizations than has typically been achievable. Upcoming satellite missions such as HETE II and SWIFT should improve upon this situation.
We are grateful to D. Frayer, A. Eichelberger, G. Oelmer for observations at Palomar. W. W. Sargent, T. Small, A. Diercks, and T. J. Galama are thanked for their contribution at Keck. JSB gratefully acknowledges support from the Fannie and John Hertz Foundation. SRK’s research is supported by grants from NSF and NASA. SGD acknowledges partial support from the Bressler Foundation.
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# Untitled Document
YITP-00-21 math-ph/0005008
Six-Vertex Model with Domain Wall Boundary Conditions and One-Matrix Model
P. Zinn-Justin<sup>*</sup><sup>*</sup>e-mail: pzinn@insti.physics.sunysb.edu
C.N. Yang Institute for Theoretical Physics
State University of New York at Stony Brook
Stony Brook, NY 11794–3840, USA
The partition function of the six-vertex model on a square lattice with domain wall boundary conditions (DWBC) is rewritten as a hermitean one-matrix model or a discretized version of it (similar to sums over Young diagrams), depending on the phase. The expression is exact for finite lattice size, which is equal to the size of the corresponding matrix. In the thermodynamic limit, the matrix integral is computed using traditional matrix model techniques, thus providing a complete treatment of the bulk free energy of the six-vertex model with DWBC in the different phases. In particular, in the anti-ferroelectric phase, the bulk free energy and a subdominant correction are given exactly in terms of elliptic theta functions.
05/2000
1. Introduction
In , V. Korepin and the author brought up the issue of the sensitivity of the six-vertex model to its boundary conditions (even in the thermodynamic limit). The motivation came mostly from some recent work on domino tilings \[2,,3,,4\], in which boundary conditions seemed to affect greatly the typical arrangement of dominos. The problem of counting domino tilings is equivalent to the six-vertex model with particular Boltzmann weights; this is schematically described on Fig. 1. Therefore it seems natural to investigate the corresponding problem for the general six-vertex model with arbitrary weights.
Fig. 1: Correspondence between vertices of the six-vertex model and small patches of a domino tiling.
The usual studies of the six-vertex model (see and references therein) are made by assuming periodic boundary conditions (PBC). In , different boundary conditions, the so-called domain wall boundary conditions (DWBC), were used (Fig. 2a), and the thermodynamic limit of the model was investigated using determinant formulae for the partition function \[6,,7\]. The main result found was an expression for the bulk free energy in the disordered phase of the model, which is different from the usual expression for the case of periodic boundary conditions. It should be noted that the DWBC correspond to the Aztec shape in the domino tiling language (see Fig. 2), which is precisely the type of tiling which was considered in \[2,,3\].
Fig. 2: a) A configuration of the six-vertex model with DWBC, and b) one possible corresponding tiling of the Aztec diamond.
Here, we use a new method to compute the bulk free energy with DWBC in all phases of the model; in particular, we obtain an independent confirmation of the results of . In section 2, starting from the determinant formula for the partition function, we shall rewrite the latter as a matrix integral, but with a measure on the space of hermitean matrices which is not necessarily smooth. In the disordered phase (section 4), the measure will turn out to be smooth, whereas in the ferroelectric and anti-ferroelectric phase (section 3 and 5) it will be discrete (when expressed in terms of the eigenvalues). The size $`N`$ of the matrices is the size of the original square lattice, and therefore the thermodynamic limit can be investigated using tools from large $`N`$ matrix models. Since the results of section 5 (concerning the anti-ferroelectric phase) are new, they are analyzed in more detail by considering various limits of the parameters, and the subleading correction of the free energy is calculated.
2. Properties of the determinant formula
We use the same notations as in . We consider the homogeneous six-vertex model, with the following parameterization of the Boltzmann weights attached to the vertices:
$$a=\mathrm{sinh}(t\gamma )b=\mathrm{sinh}(t+\gamma )c=\mathrm{sinh}(2\gamma )$$
The domain wall boundary conditions (DWBC) mean that external horizontal arrows are outgoing, whereas external vertical arrows are incoming (Fig. 2a). These boundary conditions only exist for square lattices. In \[6,,7\], it was shown that the partition function of the six-vertex model with DWBC on a $`N\times N`$ lattice could be written as:
$$Z_N=\frac{(\mathrm{sinh}(t+\gamma )\mathrm{sinh}(t\gamma ))^{N^2}}{\left(_{n=0}^{N1}n!\right)^2}\tau _N$$
where $`\tau _N`$ is a Hänkel determinant:
$$\tau _N=\underset{1i,kN}{det}\left[\frac{\mathrm{d}^{i+k2}}{\mathrm{d}t^{i+k2}}\varphi (t)\right]$$
Here,
$$\varphi (t)\frac{\mathrm{sinh}(2\gamma )}{\mathrm{sinh}(t+\gamma )\mathrm{sinh}(t\gamma )}$$
It is known that such determinants are tau-functions of the Toda semi-infinite chain hierarchy in terms of appropriate parameters. Here, as a function of $`t`$, the $`\tau _N`$ satisfy the ususal Toda equations under the bilinear form \[8,,1\]:
$$\tau _N\tau _N^{\prime \prime }\tau _N^{}{}_{}{}^{2}=\tau _{N+1}\tau _{N1}N1$$
This equation was used in to derive the bulk free energy of the model in the ferroelectric and disordered phase by making an appropriate Ansatz on the large $`N`$ form of $`\tau _N`$. Unfortunately, the Ansatz in the anti-ferroelectric phase is not that simple, as we shall see, and would be hard to justify at this point.
We shall therefore use another approach here, based on the equivalence of Hänkel determinants with one-matrix models \[9,,10\]. Let us write formally $`\varphi (t)`$ as a Laplace transform:
$$\varphi (t)=dm(\lambda )\mathrm{e}^{t\lambda }$$
where $`\mathrm{d}m(\lambda )`$ is a measure. We then notice that the derivatives of $`\varphi (t)`$ are the moments:
$$\frac{\mathrm{d}^i}{\mathrm{d}t^i}\varphi (t)=dm(\lambda )\lambda ^i\mathrm{e}^{t\lambda }$$
Inserting this into (2.1) leads to:
$$\tau _N=dm(\lambda _1)\mathrm{}dm(\lambda _N)\underset{\sigma 𝒮_N}{}(1)^\sigma \underset{i=1}{\overset{N}{}}\left[\mathrm{e}^{t\lambda _i}\lambda _i^{i+\sigma (i)2}\right]$$
We see that appears naturally the Van der Monde determinant $`\mathrm{\Delta }(\lambda _i)=det(\lambda _i^{j1})=_{i<j}(\lambda _i\lambda _j)`$. After a few elementary manipulations we find:
$$\tau _N=\frac{1}{N!}dm(\lambda _1)\mathrm{}dm(\lambda _N)\mathrm{\Delta }(\lambda _i)^2\mathrm{e}^{t_i\lambda _i}$$
If $`\mathrm{d}m(\lambda )`$ is a smooth positive measure of the form $`\mathrm{d}m(\lambda )=\mathrm{d}\lambda \mathrm{e}^{V(\lambda )}`$, then we recognize in (2.1) the expression in terms of its eigenvalues of the matrix integral:
$$\tau _NdM\mathrm{e}^{\mathrm{tr}\left[tM+V\left(M\right)\right]}$$
where $`M`$ is a hermitean $`N\times N`$ matrix, and $`\mathrm{d}M`$ is the flat measure.
As we shall see, if the measure is not smooth, we shall end up with expressions which can still be treated using appropriately adapted matrix model techniques. This is typically the case of discrete measures that appear in sums over Young diagrams \[11,,12,,13,,14,,15,,16,,17,,18\].
Expressions of the type (2.1) have been widely studied in the literature (on random matrices in particular). One important goal is to find their large $`N`$ asymptotic behavior. Here we shall mention the simplest method to find their leading large $`N`$ behavior: the saddle point method. The basic idea is that $`\mathrm{log}\mathrm{\Delta }(\lambda _i)^2`$, being a sum of $`N^2`$ terms, scales as $`N^2`$ in the large limit, whereas there are only $`N`$ variables of integration. Therefore the integral is dominated by a saddle point. An important remark is that, in order to find the saddle point, we must write our action (i.e. log of the function integrated) in such a way that all terms are of the same order $`N^2`$. Here, the term $`t_i\lambda _i`$ is naïvely of order $`N`$, and we reach the important conclusion that the $`\lambda _i`$ will scale as
$$\lambda _iN\mu _i$$
After the change of variables $`\lambda _i\mu _i`$, one can use the saddle point approximation, which gives us access to the function $`f`$ defined by
$$f=\underset{N\mathrm{}}{lim}\frac{\mathrm{log}(\tau _N/c_N)}{N^2}$$
where $`c_N(_{n=0}^{N^2}n!)^2`$. $`f`$ is essentially the bulk free energy, cf Eq. (2.1). Note that the saddle point is a very crude approximation in the sense that it does not naturally allow for a systematic computation of subleading corrections; however it will be sufficient for our purposes. We now proceed with a separate discussion of the different phases of the model.
3. Ferroelectric phase
This is the phase in which the weights are given by (2.1) with $`t`$ and $`\gamma `$ real, $`|\gamma |<t`$. We use the following decomponsition:
$$\varphi (t)=\frac{\mathrm{sinh}(2\gamma )}{\mathrm{sinh}(t+\gamma )\mathrm{sinh}(t\gamma )}=4\underset{l=0}{\overset{\mathrm{}}{}}\mathrm{e}^{2tl}\mathrm{sinh}(2\gamma l)$$
We are in the situation where the measure $`\mathrm{d}m`$ is discrete. The determinant takes the form
$$\tau _N=2^{N^2}\underset{l_1,\mathrm{},l_N=0}{\overset{\mathrm{}}{}}\mathrm{\Delta }(l_i)^2\mathrm{e}^{2t_il_i}\underset{i}{}\mathrm{sinh}(2\gamma l_i)$$
(we have neglected here, as in all subsequent calculations, constant factors which manifestly do not contribute to the bulk free energy). This expression is very close to what one encounters when studying the Plancherel measure (or other similar measures) on Young diagrams . In the context of Young diagrams, the $`l_i`$ represent the shifted highest weights $`l_i=m_i+Ni`$, where the $`m_i`$ are the usual highest weights (sizes of the rows of the diagram), and one is usually interested in the limiting shape of the Young diagram when its size is sent to infinity. There has been a lot of work on this type of expressions, both in the mathematical literature \[11,,15,,16,,17\] (the recent work being concerned with fluctuations around the limiting shape, which we shall not discuss here) and the physical literature \[12,,13,,14,,18\]. One relevant observation from is the following: after the rescaling $`\mu =l/N`$, all sums look like Rieman sums and one is tempted to replace them with integrals, and then apply the saddle point method. This is correct on condition that one imposes an additional constraint coming from the discreteness of the $`l_i`$. In Eq. (3.1), all $`l_i`$ must be distinct integers (due to the Van der Monde determinant), and therefore
$$|l_il_j|1ij$$
If we introduce the density $`\rho (\mu )\mathrm{d}\mu `$ of the $`\mu _i=l_i/N`$, normalized so that $`\rho (\mu )d\mu =1`$, then (3.1) implies that it must satisfy the inequality
$$\rho (\mu )1$$
In general, when the $`l_i`$ are trapped in a well of the potential (as is the case here), there will be a saturated region at the bottom of the well where $`\rho (\mu )=1`$, and an unsaturated region where $`\rho (\mu )<1`$.
Let us now proceed with the solution. Once the rescaling $`\mu _i=l_i/N`$ is performed, one notices that up to corrections exponentionally small in $`N`$, $`\mathrm{sinh}(2\gamma N\mu _i)\frac{1}{2}\mathrm{e}^{2|\gamma |N\mu _i}`$. Therefore
$$\tau _Nc_N^{}\mathrm{\hspace{0.17em}2}^{N^2}\underset{\mu _1,\mathrm{},\mu _N\frac{1}{N}_+}{}\mathrm{\Delta }(\mu _i)^2\mathrm{e}^{2N(t|\gamma |)_i\mu _i}$$
where $`c_N^{}N^{N^2}`$. Of course, once this simplification is made, we regognize a well-known expression; in fact, going back now to the original variables $`l_i`$ one can compute $`\tau _N`$ directly using the Cauchy identity for Schur functions. However, to emphasize the similarity with the other phases (which do not possess such a simple group-theoretic interpretation), we shall use the saddle point method, following the solution of . Since $`\tau _N`$ only depends on $`t|\gamma |`$, we temporarily set $`\gamma =0`$.
The support of the saddle point density $`\rho (\mu )`$ is expected to be of the form $`[0,\beta ]`$; the saturated region is $`[0,\alpha ]`$, whereas the unsaturated region is $`[\alpha ,\beta ]`$. We define the resolvent
$$\omega (z)=_0^\beta \frac{\mathrm{d}\mu \rho (\mu )}{z\mu }$$
for all complex $`z[0,\beta ]`$. The saddle point equations can be written in terms of $`\omega `$:
$$\omega (\mu +i0)+\omega (\mu i0)=2t\mu [\alpha ,\beta ]$$
In order to solve the equation, we first remove the logarithmic cut of $`\omega `$ with the redefinition: $`\stackrel{~}{\omega }(z)=\omega (z)\mathrm{log}\frac{\mu }{\mu \alpha }`$. $`\stackrel{~}{\omega }(z)`$ is analytic everywhere except on $`[\alpha ,\beta ]`$ and satisfies
$$\stackrel{~}{\omega }(\mu +i0)+\stackrel{~}{\omega }(\mu i0)=2t2\mathrm{log}\frac{\mu }{\mu \alpha }$$
This completely determines it to be:
$$\stackrel{~}{\omega }(z)=t\sqrt{(z\alpha )(z\beta )}_{\alpha i0}^{\beta i0}\frac{\mathrm{d}z^{}}{2i\pi (zz^{})\sqrt{(z^{}\alpha )(z^{}\beta )}}\mathrm{log}\frac{z^{}}{z^{}\alpha }$$
After some calculations, we find that
$$\omega (z)=t2\mathrm{log}\left[\frac{\sqrt{\beta (z\alpha )}+\sqrt{\alpha (z\beta )}}{\sqrt{z(\beta \alpha )}}\right]$$
The endpoints $`\alpha `$ and $`\beta `$ are determined by imposing $`\omega (z)\frac{1}{z}`$ as $`z\mathrm{}`$. This gives rise to two equations:
$$\{\begin{array}{cc}t=\mathrm{log}\frac{\sqrt{\beta }+\sqrt{\alpha }}{\sqrt{\beta }\sqrt{\alpha }}\hfill & \\ \sqrt{\alpha \beta }=1\hfill & \end{array}$$
whose solution is:
$$\alpha =\mathrm{coth}\frac{t}{2}\beta =\mathrm{tanh}\frac{t}{2}$$
In order to conclude, one expands further the function $`\omega (z)`$:
$$\omega (z)=\frac{1}{z}+\frac{\alpha +\beta }{4}\frac{1}{z^2}+\mathrm{}$$
and uses the fact that
$$\frac{f}{t}=2\mu =\frac{\alpha +\beta }{2}=\mathrm{coth}t$$
Integrating once and restoring $`\gamma `$, we have the final result
$$\mathrm{e}^f=\frac{1}{\mathrm{sinh}(t|\gamma |)}$$
which coincides with what was found in .
4. Disordered phase
In this phase, one usually rewrites the weights
$$a=\mathrm{sin}(\gamma t)b=\mathrm{sin}(\gamma +t)c=\mathrm{sin}(2\gamma )$$
with redefined parameters $`t`$ and $`\gamma `$, $`|t|<\gamma `$, and the function $`\varphi (t)=\mathrm{sin}(2\gamma )/(\mathrm{sin}(t\gamma )\mathrm{sin}(t+\gamma ))`$; the partition function is then given by
$$Z_N=\frac{(\mathrm{sin}(\gamma +t)\mathrm{sin}(\gamma t))^{N^2}}{\left(_{n=0}^{N1}n!\right)^2}\tau _N$$
with $`\tau _N`$ still given by (2.1). The Laplace transform is:
$$\varphi (t)=\frac{\mathrm{sin}(2\gamma )}{\mathrm{sin}(\gamma +t)\mathrm{sin}(\gamma t)}=_{\mathrm{}}^+\mathrm{}d\lambda \mathrm{e}^{t\lambda }\frac{\mathrm{sinh}\frac{\lambda }{2}(\pi 2\gamma )}{\mathrm{sinh}\frac{\lambda }{2}\pi }$$
This time the measure is smooth and $`\tau _N`$ is a matrix integral in the usual sense.
We must now rescale the variables $`\lambda _i`$. We choose to define $`\mu _i=\gamma \lambda _i/N`$. Then:
$$\tau _N=c_N^{}\gamma ^{N^2}_{\mathrm{}}^+\mathrm{}d\mu _1\mathrm{}d\mu _N\mathrm{\Delta }(\mu _i)^2\underset{i=1}{\overset{N}{}}\left[\frac{\mathrm{sinh}N\mu _i(\frac{\pi }{2\gamma }1)}{\mathrm{sinh}N\mu _i\frac{\pi }{2\gamma }}\mathrm{e}^{N\frac{t}{\gamma }\mu _i}\right]$$
One then simplifies the potential by using: $`\frac{\mathrm{sinh}N\mu (\frac{\pi }{2\gamma }1)}{\mathrm{sinh}N\mu \frac{\pi }{2\gamma }}\mathrm{e}^{N|\mu |}`$. Therefore,
$$\tau _Nc_N^{}\gamma ^{N^2}_{\mathrm{}}^+\mathrm{}d\mu _1\mathrm{}d\mu _N\mathrm{\Delta }(\mu _i)^2\mathrm{e}^{N_i(\frac{t}{\gamma }\mu _i|\mu _i|)}$$
Note that the matrix integral only depends on the ratio $`\zeta t/\gamma `$.
The matrix model (4.1) is fairly simple and can be solved easily in the large $`N`$ limit via the saddle point method. One introduces again the saddle point density of eigenvalues $`\rho (\mu )\mathrm{d}\mu `$, normalized so that $`\rho (\mu )d\mu =1`$. The support of $`\rho (\mu )`$ is assumed to be a single interval $`[\alpha ,\beta ]`$ ($`\alpha <0<\beta `$), due to the shape of the potential (single well centered around $`0`$). The resolvent is defined as before. The saddle point equations read:
$$\omega (\mu +i0)+\omega (\mu i0)=\zeta +\mathrm{sign}(\mu )\mu [\alpha ,\beta ]$$
where the right hand side is simply the derivative of the potential. The solution of this equation:
$$\omega (z)=\frac{1\zeta }{2}+\frac{2}{i\pi }\mathrm{log}\left[\frac{\sqrt{\beta (z\alpha )}i\sqrt{\alpha (z\beta )}}{\sqrt{z(\beta \alpha )}}\right]$$
is very similar to the ferroelectric phase; and the rest of the calculation goes along the same lines.
Requiring that $`\omega (z)\frac{1}{z}`$ as $`z\mathrm{}`$, we obtain the $`2`$ equations:
$$\{\begin{array}{cc}1\zeta =\frac{2}{i\pi }\mathrm{log}\frac{\sqrt{\beta }+i\sqrt{\alpha }}{\sqrt{\beta }i\sqrt{\alpha }}\hfill & \\ \sqrt{\alpha \beta }=\pi \hfill & \end{array}$$
which we solve for $`\alpha `$ and $`\beta `$:
$$\alpha =\pi \mathrm{tan}\frac{\pi }{4}(1\zeta )\beta =\pi \mathrm{tan}\frac{\pi }{4}(1+\zeta )$$
Noting that
$$\frac{f}{\zeta }=\frac{1}{N}\mathrm{tr}M=\frac{\alpha +\beta }{4}$$
we find
$$f=\mathrm{log}\mathrm{cos}\frac{\pi }{2}\zeta +\mathrm{cst}$$
We shall not discuss how to fix the constant of integration, since this will be addressed in the next section in a more general setting. Reintroducing the $`\gamma `$ dependence coming from Eq. (4.1), we have the final expression:
$$\mathrm{e}^f=\frac{\pi }{2\gamma }\frac{1}{\mathrm{cos}\frac{\pi t}{2\gamma }}$$
which reproduces the result of .
5. Anti-ferroelectric phase
We finally study the most interesting phase, in which the weights are given by
$$a=\mathrm{sinh}(\gamma t)b=\mathrm{sinh}(\gamma +t)c=\mathrm{sinh}(2\gamma )$$
with $`|t|<\gamma `$, and the partition function by
$$Z_N=\frac{(\mathrm{sinh}(\gamma +t)\mathrm{sinh}(\gamma t))^{N^2}}{\left(_{n=0}^{N1}n!\right)^2}\tau _N$$
with $`\varphi (t)=\mathrm{sinh}(2\gamma )/(\mathrm{sinh}(\gamma +t)\mathrm{sinh}(\gamma t))`$.
5.1. Bulk free energy
We have the expansion
$$\varphi (t)=\frac{\mathrm{sinh}(2\gamma )}{\mathrm{sinh}(\gamma +t)\mathrm{sinh}(\gamma t)}=2\underset{l=\mathrm{}}{\overset{+\mathrm{}}{}}\mathrm{e}^{2tl}\mathrm{e}^{2\gamma |l|}$$
We perform the rescaling $`\mu _i=2\gamma l_i/N`$ and find that $`\tau _N`$ takes the form:
$$\tau _N=c_N^{}\gamma ^{N^2}\underset{\mu _1,\mathrm{},\mu _N\frac{2\gamma }{N}}{}\mathrm{\Delta }(\mu _i)^2\mathrm{e}^{N_i(\frac{t}{\gamma }\mu _i|\mu _i|)}$$
The remarkable feature is that Eq. (5.1) is identical to Eq. (4.1) up to the discrete nature of the variables! We shall comment on this later.
The situation is a bit more complicated than in the previous cases, since we now expect a saturated region $`[\alpha ^{},\beta ^{}]`$ at the bottom of the well ($`\alpha ^{}<0<\beta ^{}`$) and two unsaturated regions $`[\alpha ,\alpha ^{}]`$ and $`[\beta ^{},\beta ]`$ on each side. This is a two-cut situation, which is in fact the reason why the naïve approach of fails in the anti-ferroelectric phase (see section 5.3 for more on this). Let us define as before $`\zeta =t/\gamma `$, the density $`\rho (\mu )`$ and its resolvent $`\omega (\mu )`$. The constraint coming from the discreteness of the $`\mu _i`$ reads
$$\rho (\mu )\frac{1}{2\gamma }\mu $$
Therefore we have in the saturated region the equation
$$\rho (\mu )=\frac{1}{2i\pi }(\omega (\mu i0)\omega (\mu +i0))=\frac{1}{2\gamma }\mu [\alpha ^{},\beta ^{}]$$
whereas in the unsaturated regions, the saddle point equations are
$$\omega (\mu +i0)+\omega (\mu i0)=\zeta +\mathrm{sign}(\mu )\mu [\alpha ,\alpha ^{}][\beta ^{},\beta ]$$
with $`\zeta =t/\gamma `$.
We could proceed as in the previous sections; this would lead to a representation of $`\omega (z)`$ in terms of elliptic integrals. However, this would be fairly cumbersome and we proceed instead as follows. Introduce an elliptic parameterization
$$u(\mu )=\frac{1}{2}\sqrt{(\beta ^{}\alpha )(\beta \alpha ^{})}_\beta ^\mu \frac{\mathrm{d}z}{\sqrt{(z\alpha )(z\alpha ^{})(z\beta ^{})(z\beta )}}$$
which corresponds to setting: $`\frac{\beta ^{}\alpha }{\beta \alpha }\frac{\beta \mu }{\beta ^{}\mu }=sn^2(u,k)`$ with $`k=\sqrt{\frac{(\beta \alpha )(\beta ^{}\alpha ^{})}{(\beta ^{}\alpha )(\beta \alpha ^{})}}`$. With an appropriate choice of path of integration, this maps the $`\mu `$ complex plane (resp. upper half-plane, lower half-plane) onto the rectangle $`[0,K]\times [iK^{},iK^{}]`$ (resp. $`[0,K]\times [0,iK^{}]`$, $`[0,K]\times [iK^{},0]`$), where $`K`$ and $`K^{}`$ are the usual complete elliptic integrals of the first kind. Similarly, the second sheet of the double covering is mapped onto the other half of the torus, which can be chosen to be $`[K,0]\times [iK^{},iK^{}]`$. The point of this parameterization is that the resolvent $`\omega `$ is now a well-defined function of $`u`$. In fact we have the following properties:
(i) The function $`\omega (u)`$ can be extended to a holomorphic function in the whole $`u`$ plane.
(ii) The function $`\omega (u)`$ satisfies the following functional relations (for all complex $`u`$):
$$\begin{array}{ccc}\hfill \omega (u+2iK^{})& =\omega (u)\frac{i\pi }{\gamma }\hfill & (5.1a)\hfill \\ \hfill \omega (u+2K)& =\omega (u)2\hfill & (5.1b)\hfill \\ \hfill \omega (u)+\omega (u)& =1\zeta \hfill & (5.1c)\hfill \end{array}$$
Eq. $`(5.1a)`$ is the analytic continuation of Eq. (5.1). Similarly, by combining the analytic continuations of the two equations contained in (5.1), one obtains Eqs. $`(5.1b,c)`$.
(iii) The function $`\omega (u)`$ has the following expansion near $`u_{\mathrm{}}=u(z=\mathrm{})`$:
$$\omega =\frac{2}{\sqrt{(\beta ^{}\alpha )(\beta \alpha ^{})}}(uu_{\mathrm{}})+O(uu_{\mathrm{}})^2$$
This is a rewriting of the condition $`\omega (z)\frac{1}{z}`$ at infinity. Using properties (i) and (ii) (Eqs. $`(5.1a,b)`$), we conclude that $`\frac{\mathrm{d}}{\mathrm{d}u}\omega (u)`$ is a doubly periodic holomorphic function, and so is a constant. In order to restore the coefficients of $`\omega (u)`$ we can use properties (ii) or (iii). We find that
$$\omega (u)=\frac{1}{K}(uu_{\mathrm{}})$$
plus several conditions relating the different parameters of the problem:
$$\begin{array}{ccc}\hfill \frac{K^{}}{K}& =\frac{\pi }{2\gamma }\hfill & (5.2a)\hfill \\ \hfill \sqrt{(\beta ^{}\alpha )(\beta \alpha ^{})}& =2K\hfill & (5.2b)\hfill \\ \hfill \frac{u_{\mathrm{}}}{K}& =\frac{1\zeta }{2}\hfill & (5.2c)\hfill \end{array}$$
Relation $`(5.2a)`$ is particularly interesting since it shows that the elliptic nome $`q=\mathrm{e}^{\pi K^{}/K}=\mathrm{e}^{\pi ^2/2\gamma }`$ depends only on $`\gamma `$ (and not on $`\zeta `$). Also, the dual nome (under modular transformation) $`\stackrel{~}{q}=\mathrm{e}^{2\gamma }`$ is up to a sign the quantum group deformation parameter of the model.
We can rewrite the three conditions in terms of the endpoints; we find
$$\begin{array}{cc}\hfill \beta \alpha & =2K\frac{dnu_{\mathrm{}}}{snu_{\mathrm{}}cnu_{\mathrm{}}}\hfill \\ \hfill \beta \alpha ^{}& =2K\frac{cnu_{\mathrm{}}}{snu_{\mathrm{}}dnu_{\mathrm{}}}\hfill \\ \hfill \beta \beta ^{}& =2K\frac{cnu_{\mathrm{}}dnu_{\mathrm{}}}{snu_{\mathrm{}}}\hfill \end{array}$$
In order to completely fix the four endpoints $`\alpha `$, $`\alpha ^{}`$, $`\beta ^{}`$, $`\beta `$, we need one extra relation; this is the equality of chemical potentials in the two unsaturated regions. This relation takes the form
$$_\alpha ^{}^\beta ^{}(\omega (\mu +i0)+\omega (\mu i0))d\mu =(1\zeta )\beta ^{}+(1+\zeta )\alpha ^{}$$
Using the expression (5.2) of $`\omega (u)`$, we can rewrite it as
$$\beta ^{}(\beta \beta ^{})\frac{snu_{\mathrm{}}}{cnu_{\mathrm{}}dnu_{\mathrm{}}}Z(u_{\mathrm{}})=0$$
where $`Z`$ is Jacobi’s Zeta function; this fixes $`b^{}`$ to be
$$\beta ^{}=2KZ(u_{\mathrm{}})$$
The endpoints are now determined by (5.3) and (5.3), supplemented by the value $`(5.2c)`$ of $`u_{\mathrm{}}`$.
At this point, we are ready to calculate the free energy. We first rewrite explicitly the resolvent (Eq. (5.2)) under the form
$$\omega (z)=_z^{\mathrm{}}\frac{\mathrm{d}z^{}}{\sqrt{(z^{}\alpha )(z^{}\alpha ^{})(z^{}\beta ^{})(z^{}\beta )}}$$
Next we expand it to order $`1/z^2`$ to find
$$\frac{f}{\zeta }=\frac{\alpha +\alpha ^{}+\beta ^{}+\beta }{4}$$
which generalizes Eq. (4.1); using some known identities satisfied by Zeta and theta functions, we obtain
$$\frac{f}{\zeta }=\frac{\pi }{2}\frac{\theta _2^{}(\pi \zeta /2)}{\theta _2(\pi \zeta /2)}$$
where we recall that $`\theta _2(z)`$ is
$$\theta _2(z)=2\underset{n=0}{\overset{\mathrm{}}{}}q^{(n+1/2)^2}\mathrm{cos}(2n+1)z$$
There are a variety of ways to find the integration constant. One is to calculate explicitly $`f`$ (for a particular value of $`\zeta `$, e.g. $`\zeta =0`$) using this matrix model solution, and then restore the $`\gamma `$ dependence coming from (5.1); this is a straightforward but tedious exercise. Another possibility is to use the known limits $`\zeta \pm 1`$, that is $`t\pm \gamma `$, where we should have (see )
$$\mathrm{e}^f\frac{1}{\gamma t}$$
Either way, we finally find:
$$\mathrm{e}^f=\frac{\pi }{2\gamma }\frac{\theta _1^{}(0)}{\theta _2(\frac{\pi t}{2\gamma })}$$
where we recall that the elliptic nome is $`q=\mathrm{e}^{\frac{\pi ^2}{2\gamma }}`$.
As a simple check of our calculation, note that if one sends $`\gamma `$ to $`0`$ (keeping $`\zeta `$ fixed), since the constraint (5.1), which was the only difference with the disordered phase, disappears, one should recover the results of the previous section. This is indeed what happens when one replaces the theta functions with their $`q0`$ limit. Also, (5.3) has been numerically checked with high accuracy.
This concludes the calculation of the bulk free energy in the anti-ferroelectric phase. Restated more explicitly, this is the result we have obtained: the partition function $`Z_N`$ of the six-vertex model on a $`N\times N`$ lattice with DWBC and Boltzmann weights given by (5.1) has the following large $`N`$ behavior:
$$\underset{N\mathrm{}}{lim}Z_N^{1/N^2}=\mathrm{sinh}\gamma (1\zeta )\mathrm{sinh}\gamma (1+\zeta )\frac{\pi }{2\gamma }\frac{\theta _1^{}(0)}{\theta _2(\frac{\pi \zeta }{2})}$$
where $`\zeta =t/\gamma `$, and the elliptic nome of the theta functions is $`q=\mathrm{e}^{\frac{\pi ^2}{2\gamma }}`$. Note that this expression is different from the corresponding expression for PBC. Let us now consider the two limits $`\gamma 0`$ and $`\gamma \mathrm{}`$. In both cases we shall assume that $`\zeta `$ remains fixed.
5.2. Small $`\gamma `$ limit
As one sends $`\gamma `$ to $`0`$, one reaches the line of the disordered/anti-ferroelectric phase transition. As noted earlier, the bulk free energy of the disordered phase is essentially obtained from that of the anti-ferroelectric phase by setting $`q=0`$ in the theta functions (and performing the rotation $`\gamma i\gamma `$, $`tit`$ in the prefactors). Considering that $`q=\mathrm{e}^{\frac{\pi ^2}{2\gamma }}`$, we expect a very smooth phase transition. More explicitly, we have the following expansion of $`f`$:
$$f=\mathrm{log}\left[\frac{\pi }{2\gamma }\frac{1}{\mathrm{cos}(\frac{\pi t}{2\gamma })}\right]+2\underset{m=1}{\overset{\mathrm{}}{}}\frac{1}{m}\frac{q^{2m}}{1q^{2m}}(1(1)^m\mathrm{cos}(m\pi t/\gamma ))$$
After subtraction of the analytic continuation of the disordered phase free energy (note that this analytic continuation is trivial since $`f`$ only depends on $`t/\gamma `$), we obtain the singular part of the free energy, which has a leading singularity
$$f_{\mathrm{sing}}=4\mathrm{e}^{\pi ^2/\gamma }\mathrm{cos}^2\left(\frac{\pi t}{2\gamma }\right)+\mathrm{}$$
This is the same type of singularity that appears in the model with periodic boundary conditions . In more physical terms, if we introduce a temperature $`T`$ which is near the critical temperature $`T_c`$, we have
$$f_{\mathrm{sing}}\mathrm{e}^{C/\sqrt{T_cT}}$$
that is an infinite order phase transition.
5.3. Large $`\gamma `$ limit
Next, let us consider the $`\gamma \mathrm{}`$ limit, i.e. $`\mathrm{\Delta }=\mathrm{cosh}(2\gamma )\mathrm{}`$. This is a typical zero temperature limit, and we expect that the free energy will be dominated by the contribution of a ground state. After a modular transformation, the bulk free energy reads
$$\begin{array}{cc}\hfill F& =\mathrm{log}(\mathrm{sinh}(\gamma t)\mathrm{sinh}(\gamma +t))f\hfill \\ & =\frac{\gamma }{2}\frac{t^2}{2\gamma }\mathrm{log}\mathrm{sinh}(\gamma +t)+t+2\underset{m=1}{\overset{\mathrm{}}{}}\frac{1}{m}\frac{\mathrm{e}^{2m\gamma }}{\mathrm{sinh}(2m\gamma )}\mathrm{sinh}^2(m(\gamma t))\hfill \end{array}$$
We can interpret the first terms when $`\gamma \mathrm{}`$
$$F=\frac{3}{2}\gamma \frac{t^2}{2\gamma }+O(\mathrm{e}^{2\gamma })$$
as coming from the family of ground states described by Fig. 3. The pattern of a rectangle inscribed inside a square is reminiscent of the circle inscribed inside a square characteristic of the disordered phase .
Fig. 3: Ground states of the anti-ferroelectric phase. In regions $`a`$ and $`b`$ the arrows are aligned, whereas in region $`c`$ they alternate in direction.
5.4. Subdominant corrections
As a final note, it is interesting to understand why the approach of fails in the anti-ferroelectric phase. There, the idea was to find an appropriate Ansatz on the asymptotic behavior of the determinant $`\tau _N`$ and plug it in the Toda equation (2.1). The simplest assumption is that only the leading behavior (bulk free energy) must be taken into account, which leads to replacing $`\tau _N`$ with $`c_N\mathrm{e}^{N^2f}`$, where $`c_N=(_{n=0}^{N1}n!)^2`$. The Toda equation then reduces to the ordinary differential equation for $`f`$:
$$f^{\prime \prime }=\mathrm{e}^{2f}$$
We can now use some insight from matrix models to understand whether this assumption was justified or not. In the ferroelectric and disordered phases, we reduced the computation of $`\tau _N`$ to a matrix model with eigenvalues in one single interval $`[a,b]`$ (disregarding the saturated region which plays no role here); it is known that such models have a regular large $`N`$ limit. In fact, in the ferroelectric phase one can easily prove that
$$\tau _Nc_N\mathrm{e}^{N^2f}\mathrm{e}^{N(t|\gamma |)}$$
up to only exponentially small corrections; whereas in the disordered phase, one expects an asymptotic expansion which starts with
$$\tau _Nc_N\mathrm{e}^{N^2f}N^\kappa C$$
and continues with inverse powers of $`N`$ (note that this is not quite the usual topological expansion of $`2D`$ gravity since the potential is not polynomial). In either case, the assumption on the corrections is valid, and indeed, one can check that the expressions (3.1) and (4.1) do satisfy the ODE (5.3).
On the contrary, in the anti-ferroelectric regime, we have found that the support of the eigenvalues contains two intervals $`[a,a^{}]`$ and $`[b^{},b]`$ and therefore we expect to be in a situation similar to what was studied in \[19,,20\]. The analysis shows that $`\tau _N`$ should in this case display a pseudo-periodic behavior, which is indeed what is found in numerical computations. More precisely, after some calculations along the lines of , one finds that
$$\tau _Nc_N\left[\frac{\pi }{2\gamma }\frac{\theta _1^{}(0)}{\theta _2(\frac{\pi \zeta }{2})}\right]^{N^2}\theta _4\left(\frac{\pi }{2}(1+\zeta )N\right)C$$
where $`\zeta =t/\gamma `$ and the elliptic nome $`q`$ of the theta function is as before $`q=\mathrm{e}^{\frac{\pi ^2}{2\gamma }}`$. The constant $`C`$ depends only on $`\gamma `$. One can check that the right hand side of Eq. (5.3) does satisfy the Toda equation (2.1), even though the bulk free energy alone does not satisfy the ODE (5.3).
Acknowledgements
I thank V. Korepin for discussions, and the Centre de Recherches Mathématiques de l’Université de Montréal, where part of this work was performed, for its hospitality.
References
relax V. Korepin and P. Zinn-Justin, preprint cond-mat/0004250. relax W. Jockush, J. Propp and P. Shor, preprint math.CO/9801068. relax H. Cohn, N. Elkies and J. Propp, Duke Math. Journal 85 (1996), 117. relax R. Kenyon, The planar dimer model with boundary: a survey, preprint (http://topo.math.u-psud.fr/$``$kenyon/papers.html). relax R.J. Baxter, Exactly Solved Models in Statistical Mechanics (San Diego, CA: Academic). relax A.G. Izergin, Sov. Phys. Dokl. 32 (1987), 878. relax A.G. Izergin, D.A. Coker and V.E. Korepin, J. Phys. A 25 (1992), 4315. relax K. Sogo, Journal of the Physical Society of Japan 62, 6 (1993), 1887. relax A. Gerasimov, A. Marshakov, A. Mironov, A. Morozov and A. Orlov, Nucl. Phys. B 357 (1991), 565. relax M. Adler and P. van Moerbeke, Duke Math. Journal 80 (1995), 863 (preprint solv-int/9706010); preprint math.CO/9912143. relax A.M. Vershik and S.V. Kerov, Soviet. Math. Dokl. 18 (1977), 527. relax M.R. Douglas and V.A. Kazakov, Phys. Lett. B319 (1993), 219. relax V.A. Kazakov, M. Staudacher and T. Wynter, Commun. Math. Phys. 177 (1996), 451; 179 (1996), 235; Nucl. Phys. B471 (1996), 309; I. Kostov, M. Staudacher and T. Wynter, Commun. Math. Phys. 191 (1998), 283. relax V. Kazakov and P. Zinn-Justin, Nucl. Phys. B546 (1999), 647 (preprint hep-th/9808043). relax J. Baik, P. Deift and K. Johansson, J. Amer. Math. Soc. 12 (1999) no. 4, 1119 (preprint math.CO/9810105). relax K. Johansson, preprint math.CO/9906120. relax A. Borodin, A. Okounkov and G. Olshanski, preprint math.CO/9905032. relax E. Brézin and V. Kazakov, preprint math-ph/9909009. relax P. Deift, T. Kriecherbauer, K.T-R. McLaughlin, S. Venakides and X. Zhou, Commun. on Pure and Applied Math. 52 (1999), 1491. relax G. Bonnet, F. David and B. Eynard, preprint cond-mat/0003324.
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# 1 Introduction
## 1 Introduction
Fermion zero modes of the Abelian Dirac operator in three dimensional Euclidean space (i.e., of the Pauli operator) are a rather young subject of study, and still many features remain unknown. The first example of such a zero mode of the Pauli operator has been given in , where it was used to prove that one-electron atoms with sufficiently high nuclear charge in an external magnetic field are unstable, see . Further examples of zero modes were discussed in , where their relevance for QED was briefly mentioned (this point is discussed in more detail in ). In it was proven by explicit construction of a class of examples that the phenomenon of zero mode degeneracy (i.e., Pauli operators with more than one zero mode) occurs, and a relation between the number of zero modes of a Pauli operator and the Hopf index of the corresponding magnetic field was established. This point was further elaborated in .
In an example of a zero mode was given where the corresponding gauge potential (and magnetic field) are non-zero only within a finite region of space (within a ball with finite radius). In fact, this example belongs to the types of zero modes that were discussed in . Here we want to construct a whole class of zero modes of Pauli operators where the related gauge potentials and magnetic fields vanish outside a finite region of space. This demonstrates that the possibility of having zero modes in magnetic fields of finite range is not just a curiosity that is related to some very special examples, but a rather generical feature of the Pauli operator. We think that this observation is interesting from a physical perspective as well, because magnetic fields with finite range are precisely the types of magnetic fields that may be realised experimentally.
## 2 Construction of the zero modes
We want to study specific solutions of the equation
$$i\sigma _i_i\mathrm{\Psi }(x)=A_i(x)\sigma _i\mathrm{\Psi }(x)$$
(1)
(here $`x=(x_1,x_2,x_3)`$, $`r=|x|`$, and $`\sigma _i`$ are the Pauli matrices) where $`\mathrm{\Psi }`$ is square-integrable, $`\mathrm{\Psi }(x)L^2(𝐑^3)`$, $`A_i`$ and $`B_i=ϵ_{ijk}_jA_k`$ are non-singular everywhere in $`𝐑^3`$ and are different from zero only in a finite region of space. Further $`\mathrm{\Psi }`$, $`A_i`$ and $`B_i`$ have to be smooth everywhere. For this purpose, let us first observe that the spinor ($`x_\pm x_1\pm ix_2`$)
$$\mathrm{\Psi }^0(x)=\frac{i}{r^3}\left(\begin{array}{c}x_3\\ x_+\end{array}\right)$$
(2)
solves the free Dirac equation
$$i\stackrel{}{\sigma }\stackrel{}{}\mathrm{\Psi }^0=0$$
(3)
(the $`i`$ in (2) is chosen for later convenience). The spinor (2) is singular at $`r=0`$ but it is well behaved for large $`r`$. So we might ask whether there exist spinors that are equal to $`\mathrm{\Psi }^0`$ outside a ball of radius $`r=R`$ (where they solve the free Dirac equation, i.e., $`A_i=0`$ for $`r>R`$), whereas they differ from $`\mathrm{\Psi }^0`$ inside $`r=R`$. Inside the ball they are supposed to solve the Dirac equation for some nonzero $`A_i`$ such that they are nonsingular and smooth everywhere. We shall find a whole class of such zero modes among the zero modes that were discussed in , therefore we want to review the results of briefly. There the ansatz
$$\mathrm{\Psi }=g(r)\mathrm{exp}(if(r)\frac{\stackrel{}{x}}{r}\stackrel{}{\sigma })\left(\begin{array}{c}1\\ 0\end{array}\right)=g(r)[\mathrm{cos}f(r)\mathrm{𝟏}+i\mathrm{sin}f(r)\frac{\stackrel{}{x}}{r}\stackrel{}{\sigma }]\left(\begin{array}{c}1\\ 0\end{array}\right)$$
(4)
for the spinor lead to a zero mode for the gauge field
$$A_i=h(r)\frac{\mathrm{\Psi }^{}\sigma _i\mathrm{\Psi }}{\mathrm{\Psi }^{}\mathrm{\Psi }}$$
(5)
provided that $`g(r)`$ and $`h(r)`$ are given in terms of the independent function $`f(r)`$ as ($`{}_{}{}^{}d/dr`$)
$$g^{}=\frac{2}{r}\frac{t^2}{1+t^2}g.$$
(6)
$$h=(1+t^2)^1(t^{}+\frac{2}{r}t)$$
(7)
where
$$t(r):=\mathrm{tan}f(r).$$
(8)
A sufficient condition on $`t(r)`$ leading to smooth, non-singular and $`L^2`$ spinors and smooth, non-singular gauge potentials with finite energy ($`(\stackrel{}{B})^2`$) and finite Chern–Simons action ($`\stackrel{}{A}\stackrel{}{B}`$) is
$$t(0)=0,t(r)c_1r+o(r^2)\mathrm{for}r0$$
(9)
$$t(\mathrm{})=\mathrm{}$$
(10)
which we shall assume in the sequel. Observe that in the limit $`t\mathrm{}`$ $`\stackrel{}{A}`$ vanishes whereas $`\mathrm{\Psi }`$ becomes $`\mathrm{\Psi }^0`$. Therefore, if we find some $`t`$ that become infinite at some finite $`r=R`$ in a smooth way and stay infinite for $`r>R`$, we have found precisely what we want. To get a more manageable condition, let us re-express things in terms of
$$c(r):=\mathrm{cos}f(r)=\frac{1}{(1+t(r)^2)^{1/2}}$$
(11)
which leads to
$$g^{}=\frac{2}{r}(1c^2)g$$
(12)
$$h=c(\frac{c^{}}{(1c^2)^{1/2}}+\frac{2}{r}(1c^2)^{1/2}).$$
(13)
Further $`c`$ has to behave like
$$c(r)1c_2r^2+\mathrm{}\mathrm{for}r0.$$
(14)
Now let us assume that $`c`$ approaches zero in a smooth way for $`r=R`$ and stays zero for $`rR`$, and further $`c^{}(r=R)=0`$, $`|c^{\prime \prime }(r=R)|<\mathrm{}`$. This implies that $`\stackrel{}{A}=0`$ for $`r>R`$ and that
$$g(r)=kr^2,k=\mathrm{exp}(2_0^R𝑑r\frac{1c(r)^2}{r})\mathrm{for}rR$$
(15)
which precisely leads to $`\mathrm{\Psi }=\mathrm{const}\mathrm{\Psi }^0`$ for $`r>R`$, see (4).
Finally, let us give some examples, where we choose $`R=1`$ for convenience. A first example is
$$c(r)=(1r^2)^2\mathrm{for}r<1,c(r)=0\mathrm{for}r1$$
(16)
leading to
$$g(r)=\mathrm{exp}(4r^2+3r^4\frac{4}{3}r^6+\frac{1}{4}r^8)\mathrm{for}r<1$$
$$g(r)=\mathrm{exp}(\frac{25}{12})r^2\mathrm{for}r1$$
(17)
and
$$h(r)=\frac{2(1r^2)^2(24r^2+4r^4r^6)}{(46r^2+4r^4r^6)^{1/2}}\mathrm{for}r<1$$
$$h(r)=0\mathrm{for}r1.$$
(18)
Another example is
$$c(r)=\mathrm{exp}(\frac{r^2}{r^21})\mathrm{for}r<1,c(r)=0\mathrm{for}r1.$$
(19)
There is one difference between example (16) and example (19). Both lead to $`L^2`$ zero modes (i.e., bound states), and both lead to magnetic fields with are smooth, non-singular and have finite energy. However, higher derivatives of $`c`$ in (16) are discontinuous, whereas all derivatives of $`c`$ in (19) are smooth. In some situations (or for mathematical reasons) it may be preferable to have only such $`c`$ that have only smooth higher derivatives, then functions $`c`$ like in (16) may be treated as follows with the help of functions like (19). Define a function $`c_a`$
$$c_a(r)=(1r^2)^2\mathrm{exp}(\frac{ar^2}{r^21})\mathrm{for}r<1,c(r)=0\mathrm{for}r1$$
(20)
where $`a0`$. For $`a0`$ $`c_a`$ is a $`C^{\mathrm{}}`$ function. In the limit $`a0`$ $`c_a`$ is equal to the $`c`$ of (16). Further, for $`a`$ sufficiently small, $`c_a`$ approximates the $`c`$ of (16) with arbitrary precision. Therefore, the function $`c_a`$ with a sufficiently small $`a`$ may be used as a $`C^{\mathrm{}}`$ substitute for (16).
## 3 Summary
As should be clear from the above discussion, there is an infinite number of possible functions $`c(r)`$, therefore already for the special ansatz (4) there exists a whole class of zero modes in finite range magnetic fields. One obvious generalisation of the above result is the existence of zero modes in magnetic fields that are non-zero inside a ball $`r<R_1`$, zero between $`R_1`$ and $`R_2>R_1`$, non-zero again between $`R_2`$ and $`R_3>R_2`$, etc., forming an onion-like structure.
We started from the spherically symmetric ansatz (4), therefore the finite regions where $`\stackrel{}{B}0`$ are all spherically symmetric balls. It is plausible to assume that by relaxing the symmetry condition on the zero modes one could find zero modes for magnetic fields which vanish outside finite regions of different shapes.
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# Oscillating Color Transparency in 𝜋𝐴→𝜋𝑝(𝐴-1) and 𝛾𝐴→𝜋𝑁(𝐴-1)
## Abstract
Abstract: The energy dependence of $`90^o`$ $`cm`$ fixed angle scattering of $`\pi p\pi ^{}p^{}`$ and $`\gamma p\pi ^+n`$ at large momentum transfer are found to be well described in terms of interfering short and long distance amplitudes with dynamical phases induced by Sudakov effects. We calculate the color transparency ratio for the corresponding processes in nuclear environments $`\pi A\pi ^{}p(A1)`$ and $`\gamma A\pi N(A1)`$ taking nuclear filtering into account. A prediction that the transparency ratio for these reactions will oscillate with energy provides an important test of the Sudakov phase shift and nuclear filtering hypothesis which is testable in upcoming experiments.
PACS: 13.85.Dz, 14.20.Dh, 14.40.-n
The strong interactions remain a mystery and a phenomenology. Color transparency separates conventional strong interaction physics from perturbative $`QCD`$. The perturbative calculation predicts suppression of strong interactions in certain exclusive reactions containing a large momentum transfer $`Q^2>>GeV^2`$ subprocess . Suppression is supposed to occur in initial or final state interactions with nuclear targets. The perturbative $`QCD`$ ($`pQCD`$) prediction is dramatic because it apparently contradicts the older theory in a domain of its validity . Indeed it is not clear whether color transparency is capable of being described using hadronic coordinates . At the same time, the many shortcomings of the $`pQCD`$ description at the moderate $`Q^2`$ values of experiments are well known: Hence the phenomena of color transparency play a pivotal role from either point of view.
The BNL E850 experiment of Carroll et al compared proton-proton elastic collisions with corresponding quasi-elastic nuclear processes $`pAp^{}p^{\prime \prime }(A1)`$. A transparency ratio oscillating with energy was observed. The origin of the oscillations remains controversial, and the underlying mechanism has generally been assumed to be unique to proton-proton reactions. Here we show that oscillatory color transparency is also expected in several processes involving the pion. Only recently have the full perturbative kernels needed for a $`pQCD`$ description of color transparency been completely evaluated so that the calculations have a workable paradigm . Here we report calculations predicting new phenomena observable in experiments currently underway at CEBAF , which will provide fundamental information on how pQCD may be applied to exclusive processes both in free space and in a nuclear medium. Other experimental predictions can be checked at BNL or other hadron beam laboratories. The predictions are rather distinctive, and tests of the entire framework of color transparency become available.
Consider the reaction $`\pi pm^{}p^{}`$ compared to $`\pi Am^{}N^{}(A1)`$, where $`m^{}`$ represents a meson and $`N^{}`$ represents a nucleon. These processes contain Landshoff pinch singularities, and in $`pQCD`$ are expected to show oscillations about power-law energy dependence at fixed angle. The wavelength of the oscillations (imaginary anomalous dimensions) are calculable, but have not yet been calculated. These important calculations have only been performed so far for selected diagrams occuring in $`pp`$ scattering . Now considering the $`\pi p`$ reactions, the same physics of pinch singularities again demands a factorization scheme more general than the short-distance “quark-counting” method, which is sometimes misunderstood to define $`pQCD`$. We re-iterate that among the various competing renditions –the short distance factorization of Brodsky-LePage, the asymptotic impact-parameter factorization of Botts and Sterman, and the finite $`Q^2`$ factorization of Gousset et al incorpating spin-effects, all represent $`pQCD`$ and concepts such as “asymptotic” or “short distance” are not synonymous with “$`pQCD`$”.
To represent all the diagrams and integration regions, our calculations include the transverse spatial separation $`b`$ between quarks . Several remarkable things emerge: First, the naive association of $`b1/Q`$ breaks down, and hard reactions depend on the entire region $`1/Q<b<1/\mathrm{\Lambda }_{QCD}`$. There is generic violation of the short-distance selection rule known as “hadron-helicity conservation”, a model-independent test of the short-distance framework , and $`pQCD`$ predicts non-trivial transverse and helicity-violating spin effects for large $`Q^2>>GeV^2`$ . Next, Sudakov factors regulate the approach to the pinch configurations and must be included among the kernels. The Sudakov factors reinstate the geometrical strong-interaction $`Fm`$-scale by drastically cutting off amplitudes at distances larger than $`1/\mathrm{\Lambda }_{QCD}`$. The Sudakov-improved amplidues must obey analyticity, exhibited in $`pQCD`$ by color and flavor matrix phase factors of the form $`\mathrm{exp}[i\pi c\mathrm{ln}(\mathrm{ln}(s/\mathrm{\Lambda }_{QCD}^2))]`$, handled by extending the notion of anomalous dimensions to purely imaginary numbers. In a nuclear medium large-$`b`$ regions interact inelastically with exponential attenuation, while those regions of small $`b`$ interact proportional to $`b^20`$, resulting in transparency . By depleting the long distance amplitudes, “nuclear filtering” quantum mechanically favors short distance processes in large nuclei .
All of these elements are seen in free-space $`pp`$ reactions and the $`BNL`$ color transparency experiment. In particular the free space cross section $`s^{10}d\sigma /dt`$ at fixed $`cm`$ angle $`90^0`$ oscillates with $`\mathrm{ln}(\mathrm{ln}(s/\mathrm{\Lambda }_{QCD}^2))`$. (Here $`s,t`$ are the Mandlestam variables for $`cm`$ energy-squared and momentum transfer-squared). The color transparency ratio was found to show oscillations $`180^o`$ out of phase with the free-space oscillations. In a two-component model this rather unambiquously indicates strong attenuation or “filtering” in the nuclear medium of one long-distance amplitude, and little attenuation of another short-distance component. In addition, attenuation cross sections extracted from the data are substantially smaller than the traditional 40 $`mb`$ of conventional strong interaction physics at these energies, and scaling in the variable $`Q^2/A^{1/3}`$ was observed. Consistently, the cross section in the nuclear target shows negligible oscillations with energy and apparently conforms to predictions of short-distance physics . In contrast, a model based on the hadronic basis (Farrar et al Ref. ) fails to describe the data by many standard deviations.
Correlating these observations of $`pp`$ reactions with the dynamical similarity of $`\pi p`$ reactions suggests similar phenomena should be observed. Parts of the calculations are stymied by a major difficulty: no systematic method exists to find the relative phases of exclusive amplitudes. It is not enough to calculate the phase of the asymptotically largest amplitude (the procedure of Ref. ) but it is necessary to find any sizable phase coefficient of any sizable amplitude.
As a practical resolution we have fit the $`90^ocm`$ fixed angle $`s^8d\sigma /dt`$ data for $`\pi p\pi ^{}p^{}`$ with a two-component model. The existence of this data for fixed angle scattering compiled by Blazey (Fig. 1a) appears not to be widely appreciated. Oscillations in this data show much the same features as the free-space $`pp`$ data. With $`M`$ denoting the $`22`$ amplitude for the reaction, our fit is given by:
$`s^8{\displaystyle \frac{d\sigma }{d|t|}}=|A_0`$ $`+`$ $`{\displaystyle \frac{A_1\sqrt{s}e^{ic_1\mathrm{log}\mathrm{log}q^2/\mathrm{\Lambda }_{QCD}^2}}{(\mathrm{log}s)^{d_1}}}`$ (1)
$`+`$ $`{\displaystyle \frac{A_2e^{ic_2\mathrm{log}\mathrm{log}q^2/\mathrm{\Lambda }_{QCD}^2}}{(\mathrm{log}s)^{d_2}}}|^2`$ (2)
where $`A_0,A_1,A_2,c_1,c_2,d_1,d_2`$ are real parameters. The functional form of our exponents have been updated compared to Ref. and come from expanding the imaginary parts of Sudakov exponents. In accord with the discussion, the two components, $`A_1`$ and $`A_2`$ represent regions of large $`b`$, associated Sudakov effects, and logarithmically-varying phases, while small-$`b1/Q`$ regions are described by short-distance theory. The best fit gives $`A_0=0.638,A_1=5.1,c_1=25.6,d_1=5.13,A_2=0.065,c_2=26.3,d_2=1.16`$ with $`\chi ^2/dof=1.97`$. If we include only one long distance amplitude setting $`A_2=0`$, then the best fit gives $`A_0=0.661,A_1=7.67,c_1=23.2,d_1=6.03`$ with $`\chi ^2/dof=5.01`$. In comparison the short-distance model $`s^8`$ fit gives $`\chi ^2/dof=99`$.
Now turn to the corresponding pion-initiated reaction with a nuclear target. In the two-amplitude model each component interacts with the nuclear target by a different rule. For the long distance pieces, the target measures the integration region (“transverse size”) via attenuation by the rule $`I_j=exp(k\sigma _jn𝑑z)`$ where $`z`$ is the straight-line propagation distance across the target, and $`n`$ is the nucleon density; $`\sigma _{ab}`$ is the absorptive cross section for particles $`a,b`$. We used $`\sigma _{pp}=40mb`$, and $`\sigma _{\pi p}=26mb`$. The short distance amplitude is attenuated with a model inspired by $`\sigma _S=k/(x_1x_2Q^2)`$, where $`x_1,x_2`$ are the momentum fractions of the quarks inside the proton. Since this amplitude is short-distance we set $`x_1=x_2=0.5`$ and so $`\sigma _S=k(1.6mb)(GeV^2/Q^2)`$. Short-range nuclear correlations are included in both cases. We then calculate the cross section in the nuclear case by
$`s^8{\displaystyle \frac{d\sigma _A}{d|t|}}`$ $`=`$ $`{\displaystyle \frac{1}{A}}{\displaystyle d^3xn(x)}|A_0I_S^iI_S^{fa}I_S^{fb}`$ (3)
$`+`$ $`{\displaystyle \frac{A_1\sqrt{s}e^{ic_1\mathrm{log}\mathrm{log}q^2/\mathrm{\Lambda }_{QCD}^2+i\varphi _A}}{(\mathrm{log}s)^{d_1}}}I_{\pi p}^iI_{\pi p}^fI_{pp}^f`$ (4)
$`+`$ $`{\displaystyle \frac{A_2e^{ic_2\mathrm{log}\mathrm{log}q^2/\mathrm{\Lambda }_{QCD}^2+i\varphi _A}}{(\mathrm{log}s)^{d_2}}}I_{\pi p}^iI_{\pi p}^fI_{pp}^f|^2`$ (5)
where $`A`$ is the nuclear number; superscripts $`i`$ and $`f`$ refer to initial and final state attenuation factors, respectively. The formula indicates we took into account a potential relative phase $`\varphi _A`$ between the two amplitudes due to interaction with the nucleus.
We assume Fermi-motion is taken out experimentally by overdetermined kinematic reconstruction (such as possible at $`BNL`$) and so this has not been included in the calculations. We treat $`\varphi _A`$ and $`k`$ as parameters subject to considerable uncertainty. However for the entire range of $`0<\varphi _A<2\pi `$ and varying $`5<k<10`$ the calculations are sufficiently robust to predict rather dramatic effects. In Fig. 1b we show the results for the transparency ratio, $`T=d\sigma (\pi Am^{}N^{}(A1);90^o)/dt/\left[Zd\sigma (\pi p\pi p^{};90^o)/dt\right]`$ for the two different models. The plots (Fig 1, b-c) show a striking $`180^o`$ phase shift between the oscillations of the transparency ratio and those seen in the free-space reaction. $`T`$ is less sensitive to variations of $`\varphi _A`$ compared to $`k`$: for all values of the $`\varphi _A`$ we find that $`T`$ shows significant oscillations with energy. Only for very large values of $`k>>10`$ do these oscillations disappear, a limit in which no short distance contribution effectively exists. The plots are given for large nuclei where the calculation indicates filtering will be effective: for $`A>>1`$, short distance physics predicts scaling in the variable $`Q^2/A^{1/3}`$. The theory may be extended to smaller $`A12`$, where our calculations also show a substantial effect, with less confidence regarding the importance of the short-distance component. As in the $`pp`$ case, the measurement of the transparency ratio as a function of $`s`$, or the $`A`$ dependence at fixed large $`s`$, would be capable of ruling out the hadronic-basis predictions for the same reaction, which are either monotonic (Glauber theory) or linear in the energy (exploding point-like classical expansion theory (Farrar et al , ).
Most exciting are experimental data and rapidily upcoming prospects for the processes $`\gamma p\pi ^+n`$ and $`\gamma n\pi ^{}p`$. Data exists for $`s<16GeV^2`$ and $`s<4GeV^2`$, respectively . The Jefferson Lab and CEBAF is soon expected to extend the energy range of $`\gamma n`$ reactions to about $`s=16`$ GeV<sup>2</sup> with high precision, as well as measure the color transparency ratio for this process . The short distance theory predicts $`d\sigma /dt_{90^o}s^7`$, within which framework it has been shown for asymptotically large momentum transfer that Landshoff pinches are absent. Unlike $`pp`$ and $`\pi p`$ reactions, then, where the pinch regions actually constitute the asymptotic prediction, here the Landshoff and associated Sudakov phase physics is subleading. But the asymptotic limit (infinte energy) has little weight for laboratory $`Q^2`$, and ubiquitous observations of spin-effects forbidden by short-distance theory but part of regular $`pQCD`$ argue for the more complete treatment including the pinches. Most interestingly, the existing data show considerable oscillations around power dependence (Fig. 2a). Like the $`\pi p`$ case, the existence of this data also appears not to be widely appreciated.
We fit the experimental data for $`\gamma p\pi ^+n`$ with center of mass scattering angle $`90^o`$ and $`\sqrt{s}>2`$ GeV. The best fit to the 17 data points available is shown in Fig. (2). We use the same amplitude ansatz as Eq (1) for $`s^7d\sigma /dt`$, obtaining $`A_0=0.90,A_1=2.65/s,c_1=64.5,A_2=8.01/s,c_2=126.4`$ with $`\chi ^2/dof=0.69`$. Here we have set $`d_1=d_2=4`$, as the quality of fit does not depend substantially on these parameters. The values of $`d_1`$ and $`d_2`$ were chosen to obtain a relatively flat free space behavior beyond $`\sqrt{s}=3.0`$ GeV, where the presence or absence of oscillations remains experimentaly unstudied. We arbitrarily imposed a model of short-distance physics for this region. If we set $`A_2=0`$ then the best fit gives, $`A_0=0.89`$, $`A_1=4.15`$ and $`c_1=79.8`$ with $`\chi ^2`$ per degree of freedom of 1.09. For comparison the short-distance $`s^7`$ model gives $`\chi ^2/dof=2.9`$. While our fit is favored statistically, including effects of extra parameters, the short-distance model is not ruled out in comparison. Cutting the experimental uncertainties in half would be pivotal. We mention this because the uncertainties are expected to decrease with the experiments imminent.
For the nuclear process $`\gamma A\pi ^+n+(A1)`$, we calculate $`s^7\frac{d\sigma }{d|t|}`$ with the same format as Eq (2). Results for the transparency ratio for $`A=12,56,197`$ are shown in Fig. 2. In calculating filtering factors we conservatively assume that the incident photon does not attenuate significantly. While there are many models to attenuate the photon somewhat, this allows a conservative presentation, because the effects of filtering which generate the oscillating transparency ratio are minimized. Let us note that experimentally the final state $`N`$ can be a proton or a neutron, but to predict the neutron case definitively we would need free space neutron scattering data that we do not currently have.
Observing Fig. (2), the predicted transparency ratio $`T`$ again oscillates 180<sup>o</sup> out of phase with the free space cross-section. This simple fact has so far not been appreciated as generic, and previous hadronic-basis estimates for the transparency ratio have not taken the oscillations in free space data into account, yielding monotonically increasing energy dependence. The upcoming photon-initiated experiments, then, may be on the verge of confirming a third case of oscillating fixed angle data, and oscillating color transparency.
Color transparency with a photon beam remains significantly different from hadron initiated processes. The distinction becomes clear when the $`Q^2`$ dependence of a virtual photon is used as an experimental tool. In the limit of large $`Q^2>>GeV^2`$, experimental evidence from deeply-inelastic scattering overwhelming supports the concept of a point-like photon interaction, with negligible attenuation and pertubatively understood hadronic components in scattering. The lack of pinch singularities of the large-$`Q^2`$ framework predicts power-law fading of oscillations in both free space cross section and transparency ratio in the limit of large-$`Q^2`$. The regime of large-$`Q^2`$ for photons should coincide with the regime of Bjorken scaling, so that the moderate $`Q^2`$ of existing electron beams should suffice. This would be extremely interesting and productive area to explore experimentally.
Direct tests of the hadron-helicity non-conserving character of the pinch-singularity regions are very interesting. A pion beam suggests studying reactions involving a final-state $`\rho `$ meson: again the process has pinch singularities. The $`pQCD`$ analysis indicates that oscillations of fixed-angle scattering with energy will occur, and indeed one of the points of this paper is that such oscillations are generic. The failure of short-distaince models, and dynamical importance of the pinch regions for $`\pi p\rho p`$ is supported by observations of final-state $`\rho `$-polarization density matrix elements $`\rho _{1,1}`$ of order unity. If this is due to the pinch regions, as expected , then filtering in a large nucleus should remove them. Oscillating polarization effects would be very dramatic: $`\rho _{1,1}`$ oscillating with energy at fixed angle is expected if the dynamical phases are correlated with exchange of orbital angular momentum. Counting powers of the internal coordinate $`b`$ and the units of orbital angular momentum, we can predict that at fixed large $`Q^2`$, each power of $`b^2`$ in amplitude calculations will scale like $`A^{1/3}`$ due to nuclear filtering.
To conclude, oscillating color transparency is a generic prediction of $`pQCD`$, testable with imminent experiments. We believe that the observation of oscillations in experimental data for the transparency ratio, consistently $`180^o`$ out of phase with the free space counterparts, and in three independent reactions, will be strong confirmation of nuclear filtering and the basic $`pQCD`$ understanding of color transparency.
Acknowledgements We thank Haiyan Gao for very useful discussions and for providing data for $`\gamma p`$. We also thank Jerry Blazey for help with data and Bernard Pire for helpful comments. This work was supported in part under the Department of Energy, and the Kansas Institute for Theoretical and Computational Science/ K\*STAR program.
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# Abstract
### Abstract
In this paper we complete in several aspects the picture of locally compact quantum groups. First of all we give a definition of a locally compact quantum group in the von Neumann algebraic setting and show how to deduce from it a C-algebraic quantum group. Further we prove several results about locally compact quantum groups which are important for applications, but were not yet settled in our paper . We prove a serious strengthening of the left invariance of the Haar weight, and we give several formulas connecting the locally compact quantum group with its dual. Loosely speaking we show how the antipode of the locally compact quantum group determines the modular group and modular conjugation of the dual locally compact quantum group.
## Introduction
Building on the work of Kac & Vainerman , Enock & Schwartz , Baaj & Skandalis , Woronowicz and Van Daele a precise definition of a locally compact quantum group was recently introduced by the authors in , see and for an overview. For an overview of the historic development of the theory we refer to and the introduction to . Because commutative C-algebras are always of the form $`C_0(X)`$, where $`X`$ is a locally compact space and $`C_0(X)`$ denotes the C-algebra of continuous functions on $`X`$ vanishing at infinity, arbitrary C-algebras are sometimes thought of as the algebra of continuous functions vanishing at infinity on a (non-existing) locally compact quantum space. For this reason the C-algebra framework is the most natural one to study locally compact quantum groups. The most general commutative example of a locally compact quantum group is $`C_0(G)`$ with comultiplication $`\mathrm{\Delta }:C_0(G)C_b(G\times G)`$ given by $`(\mathrm{\Delta }f)(x,y)=f(xy)`$, where $`G`$ is a locally compact group and $`C_b`$ denotes the algebra of continuous bounded functions. This philosophy is followed in where we defined ‘reduced C-algebraic quantum groups’ as the proper notion of a locally compact quantum group in the C-algebra framework. As already explained, the most general commutative example is $`C_0(G)`$ where $`G`$ is a locally compact group. Further the theory unifies compact quantum groups and Kac algebras and it includes known examples as the quantum Heisenberg group, quantum $`E(2)`$-group, quantum Lorentz group and quantum $`az+b`$-group. Within this theory one can construct a dual reduced C-algebraic quantum group and prove a Pontryagin duality theorem.
On a technical level it is often more easy to work with von Neumann algebras rather than C-algebras, certainly when dealing with weights. So, already in , we associated with every reduced C-algebraic quantum group a von Neumann algebraic quantum group and we used it to prove several results on the C-algebra level. The first aim of this paper is to give an intrinsic definition of a von Neumann algebraic quantum group and to associate with it, in a canonical way, a reduced C-algebraic quantum group. This can be thought of as the quantum analogue of the classical result of Weil (see \[18, Appendice I\]), stating that every group with an invariant measure has a unique topology turning it into a locally compact group.
A second aim of this paper is to prove some new results on both C-algebraic and von Neumann algebraic quantum groups, which are indispensable for applications. In our definition of either C-algebraic or von Neumann algebraic quantum groups we assume the existence of left and right invariant weights. But the property of invariance we assume is quite weak, and in this paper we show how a much stronger notion of invariance can be proved. The same kind of result is stated for Kac algebras in , but not proved. The first proof was given by Zsidó in (see also remark 18.23 in ). This stronger invariance property is needed whenever an action of a von Neumann algebraic quantum group on a von Neumann algebra appears: see and , but also for Kac algebra actions, and it will certainly be useful in future investigations as well.
Further we will complete the picture of the quantum group and its dual with several formulas giving a link between the antipode of the quantum group and the modular theory of its dual. Roughly speaking we obtain that
$$\widehat{T}^{}\mathrm{\Lambda }(x)=\mathrm{\Lambda }\left(S(x^{})\right)$$
for nice $`xM`$, where $`M`$ is the von Neumann algebraic quantum group, $`\mathrm{\Lambda }`$ is the GNS-map of the left invariant weight $`\phi `$ on $`M`$, $`S`$ is the antipode and $`\widehat{T}`$ is the operator appearing in the modular theory of the left invariant weight $`\widehat{\phi }`$ on the dual von Neumann algebraic quantum group: it is the closure of $`\widehat{\mathrm{\Lambda }}(\omega )\widehat{\mathrm{\Lambda }}(\omega ^{})`$ where $`\widehat{\mathrm{\Lambda }}`$ is the GNS-map of $`\widehat{\phi }`$. To these results and formulas will be referred in further research, see e.g. and .
We end this introduction with some conventions and references concerning weights and operator valued weights.
We assume that the reader is familiar with the theory of normal semi-finite faithful weights (in short, n.s.f. weights) on von Neumann algebras. Nevertheless, let us fix some notations. So let $`\phi `$ be a n.s.f. weight on a von Neumann algebra $`M`$. Then we define the following sets:
1. $`_\phi ^+=\{xM^+\phi (x)<\mathrm{}\}`$, so $`_\phi ^+`$ is a hereditary cone in $`M^+`$,
2. $`𝒩_\phi =\{xMx^{}x_\phi ^+\}`$, so $`𝒩_\phi `$ is a left ideal in $`M`$,
3. $`_\phi =`$ the linear span of $`_\phi ^+`$ in $`M`$, so $`_\phi `$ is a -subalgebra of $`M`$.
There exists a unique linear map $`F:_\phi `$ such that $`F(x)=\phi (x)`$ for all $`x_\phi ^+`$. For all $`x_\phi `$, we set $`\phi (x)=F(x)`$.
A GNS-construction for $`\phi `$ is a triple $`(H_\phi ,\pi _\phi ,\mathrm{\Lambda }_\phi )`$, where $`H_\phi `$ is a Hilbert space, $`\pi _\phi :MB(H_\phi )`$ is a normal -homomorphism and $`\mathrm{\Lambda }_\phi :𝒩_\phi H_\phi `$ is a $`\sigma `$-strong closed linear map with dense range such that (1) $`\mathrm{\Lambda }_\phi (x),\mathrm{\Lambda }_\phi (y)=\phi (y^{}x)`$ for all $`x,y𝒩_\phi `$ and (2) $`\mathrm{\Lambda }_\phi (xy)=\pi _\phi (x)\mathrm{\Lambda }_\phi (y)`$ for all $`xM`$ and $`y𝒩_\phi `$. As usual we introduce the closed densely defined linear operator $`T`$ in $`H_\phi `$ as the closure of the map $`\mathrm{\Lambda }_\phi (x)\mathrm{\Lambda }_\phi (x^{})`$ for $`x𝒩_\phi 𝒩_\phi ^{}`$. Making the polar decomposition $`T=J^{\frac{1}{2}}`$ of $`T`$ we obtain the modular operator $``$ and modular conjugation $`J`$ of $`\phi `$ with respect to the GNS-construction $`(H_\phi ,\pi _\phi ,\mathrm{\Lambda }_\phi )`$.
Consider two von Neumann algebras $`M`$, $`N`$. Let $`\phi `$ be a n.s.f. weight on $`M`$ with GNS-construction $`(H_\phi ,\pi _\phi ,\mathrm{\Lambda }_\phi )`$ and let $`\psi `$ be a n.s.f. weight on $`N`$ with GNS-construction $`(H_\psi ,\pi _\psi ,\mathrm{\Lambda }_\psi )`$. The tensor product weight $`\phi \psi `$ is a n.s.f. weight on $`MN`$ (see e.g. definition 8.2 of for a definition). This tensor product weight has a GNS-construction $`(H_\phi H_\psi ,\pi _\phi \pi _\psi ,\mathrm{\Lambda }_\phi \mathrm{\Lambda }_\psi )`$ where $`\mathrm{\Lambda }_\phi \mathrm{\Lambda }_\psi :𝒩_{\phi \psi }H_\phi H_\psi `$ is the $`\sigma `$-strong closure of the algebraic tensor product $`\mathrm{\Lambda }_\phi \mathrm{\Lambda }_\psi :𝒩_\phi 𝒩_\psi H_\phi H_\psi `$.
Let $`M`$ be any von Neumann algebra. For the definition of the extended positive part $`M^+\text{ext}`$ we refer to definition 1.1 of . For $`TM^+\text{ext}`$ and $`\omega M_{}^+`$, we set $`T,\omega =T(\omega )[0,\mathrm{}]`$. Recall that there exists an embedding $`M^+M^+\text{ext}:xx^{\mathrm{}}`$ such that $`x^{\mathrm{}},\omega =\omega (x)`$ for all $`xM^+`$ and $`\omega M_{}^+`$. We will use this embedding to identify $`M^+`$ as a subset of $`M^+\text{ext}`$.
Consider a von Neumann algebra $`M`$ and a von Neumann subalgebra $`N`$ of $`M`$. The definition of an operator valued weight from $`M`$ to $`N`$ is given in definition 2.1 of .
Now consider two von Neumann algebras $`M`$ and $`N`$ and a n.s.f. weight $`\phi `$ on $`M`$. We identify $`N`$ with $`N`$ as a von Neumann subalgebra of $`MN`$ to get into the framework of operator valued weights. The operator valued weight $`\phi \iota :(MN)^+N^+\text{ext}`$ is defined in such a way that for $`x(MN)^+`$, we have that
$$\omega \left((\phi \iota )(x)\right)=\phi \left((\iota \omega )(x)\right).$$
As for weights we define the following sets:
1. $`_{\phi \iota }^+=\{x(MN)^+(\phi \iota )(x)N^+\}`$, so $`_{\phi \iota }^+`$ is a hereditary cone of $`(MN)^+`$,
2. $`𝒩_{\phi \iota }=\{xMNx^{}x_{\phi \iota }^+\}`$, so $`𝒩_{\phi \iota }`$ is a left ideal in $`MN`$,
3. $`_{\phi \iota }=`$ the linear span of $`_{\phi \iota }^+`$ in $`MN`$, so $`_{\phi \iota }`$ is a -subalgebra of $`MN`$.
There exists a unique linear map $`G:_{\phi \iota }N`$ such that $`G(x)=(\phi \iota )(x)`$ for all $`x_{\phi \iota }^+`$. For all $`x_{\phi \iota }`$, we set $`(\phi \iota )(x)=G(x)`$. Let $`a_\phi `$ and $`bN`$. Then it is easy to see that $`ab`$ belongs to $`_{\phi \iota }`$ and $`(\phi \iota )(ab)=\phi (a)b`$.
Thanks to the remark after lemma 1.4 of , we also have the following characterization of $`_{\phi \iota }^+`$: Let $`x(MN)^+`$, then $`x`$ belongs to $`_{\phi \iota }^+`$ $``$ $`\phi ((\iota \omega )(x))<\mathrm{}`$ for all $`\omega M_{}^+`$.
Let $`x𝒩_{\phi \iota }`$ and $`\omega N_{}`$. The inequality $`(\iota \omega )(x)^{}(\iota \omega )(x)\omega (\iota |\omega |)(x^{}x)`$ will imply that $`(\iota \omega )(x)𝒩_\phi `$ and
$$\mathrm{\Lambda }_\phi ((\iota \omega )(x))\omega (\phi \iota )(x^{}x)^{\frac{1}{2}}.$$
When $`L`$ is some set of elements of a space we denote by $`L`$ the linear span of $`L`$ and by $`[L]`$ the closed linear span. The symbol $``$ will denote either a von Neumann algebraic tensor product or a tensor product of Hilbert spaces and $`\iota `$ will denote the identity map. Finally we use the symbol $`\chi `$ to denote the flip map from $`MN`$ to $`NM`$, where $`N`$ and $`M`$ are von Neumann algebras. We use $`\mathrm{\Sigma }`$ to denote the flip map from $`HK`$ to $`KH`$ when $`H`$ and $`K`$ are Hilbert spaces.
## 1. Von Neumann algebraic quantum groups
We state the definition of a von Neumann algebraic quantum group and discuss how the C$`^{}`$-algebraic theory can be translated to the von Neumann algebraic setting. The major difference between both approaches is the absence of density conditions in the definition of von Neumann algebraic quantum groups: these will follow automaticly!
###### Definition 1.1.
Consider a von Neumann algebra $`M`$ together with a unital normal -homomorphism $`\mathrm{\Delta }:MMM`$ such that $`(\mathrm{\Delta }\iota )\mathrm{\Delta }=(\iota \mathrm{\Delta })\mathrm{\Delta }`$. Assume moreover the existence of
1. a n.s.f. weight $`\phi `$ on $`M`$ that is left invariant: $`\phi ((\omega \iota )\mathrm{\Delta }(x))=\phi (x)\omega (1)`$ for all $`\omega M_{}^+`$ and $`x_\phi ^+`$.
2. a n.s.f. weight $`\psi `$ on $`M`$ that is right invariant: $`\psi ((\iota \omega )\mathrm{\Delta }(x))=\psi (x)\omega (1)`$ for all $`\omega M_{}^+`$ and $`x_\psi ^+`$.
Then we call the pair $`(M,\mathrm{\Delta })`$ a von Neumann algebraic quantum group.
In the next part of this section we will list the essential properties of these quantum groups. Most of the time the proofs in can be easily translated to the von Neumann algebraic setting by replacing the norm and strict topology in the considerations by the $`\sigma `$-strong topology. However, some care has to be taken to prove the density conditions and we will discuss this in detail.
For the rest of this section we fix a von Neumann algebraic quantum group $`(M,\mathrm{\Delta })`$. Without loss of generality, we may and will assume that $`M`$ is in standard form with respect to a Hilbert space $`H`$.
At the same time we fix a n.s.f. left invariant weight $`\phi `$ on $`(M,\mathrm{\Delta })`$ together with a GNS-construction $`(H,\iota ,\mathrm{\Lambda })`$ (which is possible because $`M`$ is in standard form). We let $``$ denote the modular operator and $`J`$ the modular conjugation of $`\phi `$ with respect to this GNS-construction $`(H,\iota ,\mathrm{\Lambda })`$.
We also choose a n.s.f. right invariant weight $`\psi `$ on $`(M,\mathrm{\Delta })`$ together with a GNS-construction $`(H,\iota ,\mathrm{\Gamma })`$ (later on, we will introduce some canonical choice for $`\psi `$ and $`\mathrm{\Gamma }`$).
By left invariance of $`\phi `$, we get that $`(\omega \iota )\mathrm{\Delta }(x)𝒩_\phi `$ and $`\mathrm{\Lambda }((\omega \iota )\mathrm{\Delta }(x))\omega \mathrm{\Lambda }(x)`$ for all $`x𝒩_\phi `$ and $`\omega M_{}`$. Arguing as in result 2.6 of , we also get for $`a,b𝒩_\psi `$ and $`x𝒩_\phi `$ that $`(\psi \iota )(\mathrm{\Delta }(b^{}x)(a1))𝒩_\phi `$ and $`\mathrm{\Lambda }\left((\psi \iota )(\mathrm{\Delta }(b^{}x)(a1))\right)\mathrm{\Gamma }(a)\mathrm{\Gamma }(b)\mathrm{\Lambda }(c)`$.
Along the way to the proof of theorem 1.2, one also translates proposition 3.15 of , giving rise to the important equalities
$`H`$ $`=`$ $`[\mathrm{\Lambda }((\omega \iota )\mathrm{\Delta }(x))x𝒩_\phi ,\omega M_{}]`$ (1.1)
$`=`$ $`[\mathrm{\Lambda }\left((\psi \iota )(\mathrm{\Delta }(b^{}x)(a1))\right)x𝒩_\phi ,a,b𝒩_\psi ].`$ (1.2)
The left invariance of $`\phi `$ implies that $`\mathrm{\Delta }(y)(x1)𝒩_{\phi \phi }`$ for all $`x,y𝒩_\phi `$ and that
$$(\mathrm{\Lambda }\mathrm{\Lambda })(\mathrm{\Delta }(y_1)(x_11)),(\mathrm{\Lambda }\mathrm{\Lambda })(\mathrm{\Delta }(y_2)(x_21))=\mathrm{\Lambda }(x_1)\mathrm{\Lambda }(y_1),\mathrm{\Lambda }(x_2)\mathrm{\Lambda }(y_2)$$
for all $`x_1,x_2,y_1,y_2𝒩_\phi `$. The proof of the next result is an easy translation of the proof of theorem 3.16 of .
###### Theorem 1.2.
There exists a unique unitary element $`WB(HH)`$ such that $`W^{}(\mathrm{\Lambda }(x)\mathrm{\Lambda }(y))=(\mathrm{\Lambda }\mathrm{\Lambda })(\mathrm{\Delta }(y)(x1))`$ for all $`x,y𝒩_\phi `$.
It should be noted that $`(\omega \iota )(W^{})\mathrm{\Lambda }(x)=\mathrm{\Lambda }((\omega \iota )\mathrm{\Delta }(x))`$ for all $`x𝒩_\phi `$ and $`\omega B(H)_{}`$ (see e.g. result 2.10 of ). Using the commutant theorem for the tensor product of von Neumann algebras, this implies that $`W`$ is a unitary element in $`MB(H)`$.
Using the formula for $`W^{}`$ above, one sees that $`\mathrm{\Delta }(x)=W^{}(1x)W`$ for all $`xM`$. Applying the techniques of proposition 3.18 of , one checks that $`W`$ satisfies the pentagonal equation: $`W_{12}W_{13}W_{23}=W_{23}W_{12}`$. We call $`W`$ the multiplicative unitary of $`(M,\mathrm{\Delta })`$ with respect to the GNS-construction $`(H,\iota ,\mathrm{\Lambda })`$.
It goes without saying that all these results also have their right invariant counterparts. For later purposes we introduce the unitary element $`VB(H)M`$ such that $`V(\mathrm{\Gamma }(x)\mathrm{\Gamma }(y))=(\mathrm{\Gamma }\mathrm{\Gamma })(\mathrm{\Delta }(x)(1y))`$ for all $`x,y𝒩_\psi `$. As in result 2.10 of , one proves that
$$(\omega _{\mathrm{\Gamma }(a),\mathrm{\Gamma }(b)}\iota )(V^{})=(\psi \iota )(\mathrm{\Delta }(b^{})(a1))\text{for all}a,b𝒩_\psi .$$
(1.3)
The proof of proposition 3.22 of survives the translation to the von Neumann algebra setting. Combining this with equation (1.2) we arrive at the following conclusion.
###### Proposition 1.3.
There exists a unique densely defined closed antilinear operator $`G`$ in $`H`$ such that
$$\mathrm{\Lambda }\left((\psi \iota )(\mathrm{\Delta }(x^{})(y1))\right)x,y𝒩_\phi ^{}𝒩_\psi $$
is a core for $`G`$ and
$$G\mathrm{\Lambda }\left((\psi \iota )(\mathrm{\Delta }(x^{})(y1))\right)=\mathrm{\Lambda }\left((\psi \iota )(\mathrm{\Delta }(y^{})(x1))\right)$$
for $`x,y𝒩_\phi ^{}𝒩_\psi `$. We have moreover that $`G`$ is involutive.
By taking the polar decomposition of $`G`$, we get the following essential operators in $`H`$.
###### Notation 1.4.
We define $`N=G^{}G`$, so $`N`$ is a strictly positive operator in $`H`$. We also define the anti-unitary operator $`I`$ on $`H`$ such that $`G=IN^{\frac{1}{2}}`$.
Because $`G`$ is involutive, we have that $`I=I^{}`$, $`I^2=1`$ and $`INI=N^1`$.
A careful analysis of the proof of proposition 5.5 of reveals that this result remains true in the present setting. By equation (1.3) and the techniques used in the proof of proposition 5.8 of , this is equivalent to saying that
$$(\omega _{v,w}\iota )(V^{})GG(\omega _{w,v}\iota )(V^{})\text{and}(\omega _{v,w}\iota )(V)G^{}G^{}(\omega _{w,v}\iota )(V)$$
for all $`v,wH`$. Hence, appealing to the proof of result 5.10 of , we arrive at the vital commutation relation
$$V(_\psi N)=(_\psi N)V,$$
(1.4)
where $`_\psi `$ denotes the modular operator of $`\psi `$ with respect to the GNS-construction $`(H,\iota ,\mathrm{\Gamma })`$.
Up till now, we did not need the density conditions that are present in the definition of reduced C$`^{}`$-algebraic quantum groups (see definition 4.1 of ). This is the case because we were only working on the Hilbert space level for which the relevant density conditions are already established in equations (1.1) and (1.2). In order to further develop the theory along the lines of , we will now prove similar density conditions on the level of the von Neumann algebra $`M`$. The idea of the proof is taken from \[4, 2.7.6\].
###### Proposition 1.5.
Denoting by <sup>-</sup> the $`\sigma `$-strong closure we have
$`M`$ $`=(\omega \iota )\mathrm{\Delta }(x)xM,\omega M_{}^{}`$
$`=(\iota \omega )\mathrm{\Delta }(x)xM,\omega M_{}^{}`$
$`=\{(\omega \iota )(V)\omega B(H)_{}\}^{}.`$
###### Proof.
Define $`𝒯_\psi `$ to be the Tomita -algebra of $`\psi `$. From formula (1.3) it follows that
$`\{(\omega \iota )(V)\omega B(H)_{}\}^{}`$ $`=(\psi \iota )\left((ca^{}1)\mathrm{\Delta }(b)\right)a,b𝒩_\psi ,c𝒯_\psi ^{}`$
$`=(\psi \iota )\left((a^{}1)\mathrm{\Delta }(b)(\sigma _i^\psi (c)1)\right)a,b𝒩_\psi ,c𝒯_\psi ^{}`$
$`=(\omega \iota )\mathrm{\Delta }(x)xM,\omega M_{}^{}.`$
Now we define
$$M_r=(\omega \iota )\mathrm{\Delta }(x)xM,\omega M_{}^{}.$$
Because $`V`$ is a multiplicative unitary the linear space $`\{(\omega \iota )(V)\omega B(H)_{}\}`$ is an algebra that acts non-degenerately on $`H`$. Because $`M_r`$ is clearly self-adjoint, we get that $`M_r`$ is a von Neumann subalgebra of $`M`$. Working with the von Neumann algebraic quantum group $`(M,\chi \mathrm{\Delta })`$ instead of $`(M,\mathrm{\Delta })`$ we obtain that also
$$M_l=(\iota \omega )\mathrm{\Delta }(x)xM,\omega M_{}^{}$$
is a von Neumann subalgebra of $`M`$. Observe that it follows from the commutant theorem for the tensor product of von Neumann algebras that $`\mathrm{\Delta }(x)M_lM_r`$ for all $`xM`$.
Then we conclude from equation (1.4) that it is possible to define a one-parameter group $`(\tau _t)_t`$ of automorphisms of $`M_r`$ by $`\tau _t(x)=N^{it}xN^{it}`$ for all $`xM_r`$ and $`t`$. It also follows from equation (1.4) and the fact $`\mathrm{\Delta }(x)=V(x1)V^{}`$ for all $`xM`$, that we have $`\mathrm{\Delta }(\sigma _t^\psi (x))=(\sigma _t^\psi \tau _t)\mathrm{\Delta }(x)`$ for all $`xM`$ and $`t`$, which makes sense because $`\mathrm{\Delta }(x)MM_r`$. For the same reason we can write
$$M_l=\{(\iota \omega )\mathrm{\Delta }(x)xM,\omega (M_r)_{}\}^{}$$
and because $`\sigma _t^\psi \left((\iota \omega )\mathrm{\Delta }(x)\right)=(\iota \omega \tau _t)\mathrm{\Delta }(\sigma _t^\psi (x))`$ for all $`\omega (M_r)_{}`$ and $`xM`$, we get $`\sigma _t^\psi (M_l)=M_l`$ for all $`t`$. By the right invariance of $`\psi `$ it follows that the restriction $`\psi _l`$ of $`\psi `$ to $`M_l`$ is semifinite. By Takesaki’s theorem (see e.g. \[13, 10.1\]) there exists a unique normal faithful conditional expectation $`E`$ from $`M`$ to $`M_l`$ satisfying $`\psi (x)=\psi _l(E(x))`$ for all $`xM^+`$. From \[13, 10.2\] it follows that $`E(x)P=PxP`$ for all $`xM`$, where $`P`$ denotes the orthogonal projection onto the closure of $`\mathrm{\Gamma }(𝒩_\psi M_l)`$. So the range of $`P`$ contains $`\mathrm{\Gamma }\left((\iota \omega )\mathrm{\Delta }(x)\right)`$ for all $`\omega M_{}`$ and $`x𝒩_\psi `$. By the right invariant version of equation (1.1) we get that $`P=1`$. So $`E`$ is the identity map and $`M_l=M`$.
Working with the von Neumann algebraic quantum group $`(M,\chi \mathrm{\Delta })`$ we obtain $`M=M_r`$. We already proved that $`M_r`$ is the $`\sigma `$-strong closure of $`\{(\omega \iota )(V)\omega B(H)_{}\}`$ and so this concludes the proof of the proposition. ∎
Because we have proved that $`(M,\mathrm{\Delta })`$ satisfies the above density conditions, it is straightforward to translate the rest of the the proofs in to the von Neumann algebraic setting. In the following part of this section, we collect the most important results (we will not stick to the order as they appear in ).
### Uniqueness of invariant weights
An essential result is the uniqueness of left and right invariant weights. If $`\theta `$ is a normal semi-finite left invariant weight on $`(M,\mathrm{\Delta })`$, then there exists a non-negative number $`r`$ such that $`\theta =r\phi `$. A similar result holds for right invariant weights.
### The antipode and its polar decomposition
The antipode of our quantum group is defined through its polar decomposition. There exists a strongly continuous one-parameter group $`\tau `$ on $`M`$ such that $`\tau _t(x)=N^{it}xN^{it}`$ for all $`xM`$ and $`t`$. At the same time, we have a -anti-automorphism $`R`$ on $`M`$ such that $`R(x)=Ix^{}I`$ for all $`xM`$.
Then $`R^2=\iota `$, $`R`$ and $`\tau `$ commute and we define $`S=R\tau _{\frac{i}{2}}=\tau _{\frac{i}{2}}R`$. Note that these 3 properties determine the pair $`R`$, $`\tau `$ completely in terms of the map $`S`$. The map $`S:D(S)MM`$ is a $`\sigma `$-strongly closed map with $`\sigma `$-strong dense domain and range that is determined by $`(M,\mathrm{\Delta })`$ through the following so-called strong left invariance properties.
We have for all $`a,b𝒩_\phi `$ that $`(\iota \phi )(\mathrm{\Delta }(a^{})(1b))D(S)`$ and
$$S\left((\iota \phi )(\mathrm{\Delta }(a^{})(1b))\right)=(\iota \phi )((1a^{})\mathrm{\Delta }(b)).$$
(1.5)
The space $`(\iota \phi )(\mathrm{\Delta }(a^{})(1b))a,b𝒩_\phi `$ is a core for $`S`$. A similar result holds for right invariant weights (see proposition 5.24 of ). For other characterizations of $`S`$ we refer to proposition 5.33 and corollary 5.34 of .
We refer to $`S`$ as the antipode of the quantum group $`(M,\mathrm{\Delta })`$. The one-parameter group $`\tau `$ is called the scaling group of $`(M,\mathrm{\Delta })`$, the map $`R`$ is called the unitary antipode of $`(M,\mathrm{\Delta })`$.
There also exists a unique strictly positive number $`\nu `$ such that $`\phi \tau _t=\nu ^t\phi `$ for all $`t`$. We call $`\nu `$ the scaling constant of $`(M,\mathrm{\Delta })`$. In connection with this relative invariance, it is useful to define the strictly positive operator $`P`$ in $`H`$ such that $`P^{it}\mathrm{\Lambda }(x)=\nu ^{\frac{t}{2}}\mathrm{\Lambda }(\tau _t(x))`$ for all $`t`$ and $`x𝒩_\phi `$. We observe that $`\tau _t(x)=P^{it}xP^{it}`$ for all $`t`$ and $`xM`$.
### The right Haar weight and the modular element
Because we have the equation $`\chi (RR)\mathrm{\Delta }=\mathrm{\Delta }R`$, we get that $`\phi R`$ is a right invariant n.s.f. weight on $`(M,\mathrm{\Delta })`$. From now on we suppose that $`\psi =\phi R`$. Let $`\sigma ^{}`$ denote the modular group of $`\psi `$. Remember that $`\sigma _t^{}=R\sigma _tR`$ for all $`tR`$. We have that $`\phi \sigma _t^{}=\nu ^t\phi `$, $`\psi \sigma _t=\nu ^t\psi `$ and $`\psi \tau _t=\nu ^t\psi `$ for all $`t`$.
By the Radon Nikodym theorem 5.5 of , we get the existence of a unique strictly positive element $`\delta `$ affiliated to $`M`$ such that $`\sigma _t(\delta )=\nu ^t\delta `$ for all $`t`$ and $`\psi =\phi _\delta `$ (see definition 1.5 of for the precise definition of $`\phi _\delta `$). Formally we have $`\psi (x)=\phi (\delta ^{1/2}x\delta ^{1/2})`$. The element $`\delta `$ is called the modular element of $`(M,\mathrm{\Delta })`$.
We have that $`\mathrm{\Delta }(\delta )=\delta \delta `$, $`R(\delta )=\delta ^1`$ and $`\tau _t(\delta )=\delta `$ for all $`t`$.
Now we choose the GNS-construction $`(H,\iota ,\mathrm{\Gamma })`$ for $`\psi `$ such that $`\mathrm{\Gamma }=\mathrm{\Lambda }_\delta `$ (see the remarks before proposition 1.15 in for a precise definition of $`\mathrm{\Lambda }_\delta `$). We denote the modular operator of $`\psi `$ in this GNS-construction by $`\text{ }`$. Recall that $`\nu ^{\frac{i}{4}}J`$ is the modular conjugation of $`\psi `$ with respect to this same GNS-construction.
### The fundamental commutation relations
A full-fledged theory of quantum groups would be impossible without the following list of commutation relations.
1. The one-parameter groups $`\tau `$, $`\sigma `$ and $`\sigma ^{}`$ commute pairwise.
2. For all $`t`$ we have
$$\begin{array}{ccccccc}\hfill \mathrm{\Delta }\sigma _t& =& (\tau _t\sigma _t)\mathrm{\Delta }\hfill & & \hfill \mathrm{\Delta }\sigma _t^{}& =& (\sigma _t^{}\tau _t)\mathrm{\Delta }\hfill \\ \hfill \mathrm{\Delta }\tau _t& =& (\tau _t\tau _t)\mathrm{\Delta }\hfill & & \hfill \mathrm{\Delta }\tau _t& =& (\sigma _t\sigma _t^{})\mathrm{\Delta }\hfill \end{array}$$
(1.6)
3. On the Hilbert space level, we get that
$`(IJ)W`$ $`=`$ $`W^{}(IJ)`$ (1.7)
$`(N^1)W`$ $`=`$ $`W(N^1)`$ (1.8)
### The dual von Neumann algebraic quantum group
The multiplicative unitary $`W`$ is manageable in the sense of with $`P`$ as the managing positive operator (see proposition 6.10 of ). We follow more or less chapter 3 of to obtain the Haar weight on the dual von Neumann algebraic quantum group (see section 8 of for the proofs of the next results).
###### Definition 1.6.
Define $`\widehat{M}`$ to be the $`\sigma `$-strong closure of the algebra $`\{(\omega \iota )(W)\omega B(H)_{}\}`$. Then $`\widehat{M}`$ is a von Neumann algebra and there exists a unique unital normal -homomorphism $`\widehat{\mathrm{\Delta }}:\widehat{M}\widehat{M}\widehat{M}`$ such that $`\widehat{\mathrm{\Delta }}(x)=\mathrm{\Sigma }W(x1)W^{}\mathrm{\Sigma }`$ for all $`x\widehat{M}`$. The pair $`(\widehat{M},\widehat{\mathrm{\Delta }})`$ is again a von Neumann algebraic quantum group, referred to as the dual of $`(M,\mathrm{\Delta })`$.
The predual $`M_{}`$ is a Banach algebra if we define the product such that $`\omega \theta =(\omega \theta )\mathrm{\Delta }`$ for all $`\omega ,\theta M_{}`$ (of course, $`M_{}`$ should be thought of as the space of $`L^1`$-functions of $`M`$). Moreover, the map $`\lambda :M_{}\widehat{M}:\omega (\omega \iota )(W)`$ is an injective morphism of algebras.
Let us recall the construction of the dual weight $`\widehat{\phi }`$. First of all, we define
$$=\{\omega M_{}M^+:|\omega (x^{})|M\mathrm{\Lambda }(x)\text{ for all }x𝒩_\phi \}.$$
By the Riesz theorem for Hilbert spaces there exists for every $`\omega `$ a unique element $`\xi (\omega )H`$ such that $`\omega (x^{})=\xi (\omega ),\mathrm{\Lambda }(x)`$ for all $`x𝒩_\phi `$. Then $``$ is a left ideal in $`M_{}`$, the map $`H:\omega \xi (\omega )`$ is linear and $`\lambda (\eta )\xi (\omega )=\xi (\eta \omega )`$ for all $`\eta M_{}`$ and $`\omega `$.
There exists a unique $`\sigma `$-strong–norm closed linear map $`\widehat{\mathrm{\Lambda }}`$, with $`\sigma `$-strong dense domain $`D(\widehat{\mathrm{\Lambda }})\widehat{M}`$, into $`H`$ such that $`\lambda ()`$ is a $`\sigma `$-strong–norm core for $`\widehat{\mathrm{\Lambda }}`$ and $`\widehat{\mathrm{\Lambda }}(\lambda (\omega ))=\xi (\omega )`$ for all $`\omega `$. The dual weight $`\widehat{\phi }`$ is the unique n.s.f. weight on $`\widehat{M}`$ having the triple $`(H,\iota ,\widehat{\mathrm{\Lambda }})`$ as a GNS-construction. It turns out that $`\widehat{\phi }`$ is left invariant with respect to $`(\widehat{M},\widehat{\mathrm{\Delta }})`$. We denote the modular group of $`\widehat{\phi }`$ by $`\widehat{\sigma }`$.
We denote the antipode, unitary antipode and scaling group of $`(\widehat{M},\widehat{\mathrm{\Delta }})`$ by $`\widehat{S}`$, $`\widehat{R}`$ and $`\widehat{\tau }`$ respectively. The scaling constant of $`(\widehat{M},\widehat{\mathrm{\Delta }})`$ is equal to $`\nu ^1`$. Define the right invariant n.s.f. weight $`\widehat{\psi }`$ on $`(\widehat{M},\widehat{\mathrm{\Delta }})`$ as $`\widehat{\psi }=\widehat{\phi }\widehat{R}`$. The modular group of $`\widehat{\psi }`$ will be denoted by $`\widehat{\sigma }^{}`$. Denote the modular element of $`(\widehat{M},\widehat{\mathrm{\Delta }})`$ by $`\widehat{\delta }`$. Referring to the fact that $`\widehat{\psi }=\widehat{\phi }_{\widehat{\delta }}`$, we define the GNS-construction $`(H,\iota ,\widehat{\mathrm{\Gamma }})`$ such that $`\widehat{\mathrm{\Gamma }}=\widehat{\mathrm{\Lambda }}_{\widehat{\delta }}`$.
The modular operator and modular conjugation of $`\widehat{\phi }`$ with respect to $`(H,\iota ,\widehat{\mathrm{\Lambda }})`$ will be denoted by $`\widehat{}`$ and $`\widehat{J}`$ respectively. We will denote the modular operator of $`\widehat{\psi }`$ with respect to $`(H,\iota ,\widehat{\mathrm{\Gamma }})`$ by $`\widehat{\text{ }}`$.
It is also worth mentioning that $`P^{it}\widehat{\mathrm{\Lambda }}(x)=\nu ^{\frac{t}{2}}\widehat{\mathrm{\Lambda }}(\widehat{\tau }_t(x))`$ for all $`t`$ and $`x𝒩_{\widehat{\phi }}`$, which means in a sense that $`\widehat{P}=P`$.
Finally we mention that $`M\widehat{M}=`$.
### Pontryagin duality
As in the previous paragraph, we can also construct the dual $`(\widehat{M}\widehat{\text{}},\widehat{\mathrm{\Delta }}\widehat{\text{}})`$ of $`(\widehat{M},\widehat{\mathrm{\Delta }})`$. Notice that the construction of the dual depends on the choice of the GNS-construction of the left Haar weight. If we use the GNS-construction $`(H,\iota ,\widehat{\mathrm{\Lambda }})`$ for the construction of the dual $`(\widehat{M}\widehat{\text{}},\widehat{\mathrm{\Delta }}\widehat{\text{}})`$ , the Pontryagin duality theorem tells us that $`(\widehat{M}\widehat{\text{}},\widehat{\mathrm{\Delta }}\widehat{\text{}})=(M,\mathrm{\Delta })`$. We even have that $`\widehat{\phi }\widehat{\text{}}=\phi `$ and $`\widehat{\mathrm{\Lambda }}\widehat{\text{}}=\mathrm{\Lambda }`$.
Since $`(\widehat{M}\widehat{\text{}},\widehat{\mathrm{\Delta }}\widehat{\text{}})=(M,\mathrm{\Delta })`$, we get that $`\widehat{\delta }\widehat{\text{}}=\delta `$. Hence $`\widehat{\mathrm{\Gamma }}\widehat{\text{}}=\mathrm{\Lambda }_\delta =\mathrm{\Gamma }`$.
### From von Neumann algebraic to C$`^{}`$-algebraic quantum groups
In , we associated to any reduced C$`^{}`$-algebraic quantum group a von Neumann algebraic quantum group by taking the $`\sigma `$-strong closure of the underlying C$`^{}`$-algebra in the GNS-space of a left Haar weight. In the last part of this section we go the other way around by introducing a C$`^{}`$-algebraic quantum group.
To distinguish between von Neumann algebraic and C$`^{}`$-algebraic tensor products we will denote the minimal C$`^{}`$-tensor product by $`_\text{c}`$.
###### Proposition 1.7.
Define $`M_\text{c}`$ to be the norm closure of the space $`\{(\iota \omega )(W)\omega B(H)_{}\}`$ and $`\mathrm{\Delta }_\text{c}`$ to be the restriction of $`\mathrm{\Delta }`$ to $`M_\text{c}`$. Then the pair $`(M_\text{c},\mathrm{\Delta }_\text{c})`$ is a reduced C$`^{}`$-algebraic quantum group.
###### Proof.
Because $`W`$ is manageable and $`\mathrm{\Delta }_\text{c}(x)=W^{}(1x)W`$ for all $`xM_\text{c}`$, propositions 1.5 and 5.1 of imply that $`M_\text{c}`$ is a C$`^{}`$-algebra, $`\mathrm{\Delta }_\text{c}`$ is a non-degenerate -homomorphism from $`M_\text{c}`$ into the multiplier algebra of $`M_\text{c}_\text{c}M_\text{c}`$, such that $`(\mathrm{\Delta }_\text{c}_\text{c}\iota )\mathrm{\Delta }_\text{c}=(\iota _\text{c}\mathrm{\Delta }_\text{c})\mathrm{\Delta }_\text{c}`$ and both $`\mathrm{\Delta }_\text{c}(M_\text{c})(M_\text{c}1)`$ and $`\mathrm{\Delta }_\text{c}(M_\text{c})(1M_\text{c})`$ are dense in $`M_\text{c}_\text{c}M_\text{c}`$.
Now define $`\phi _\text{c}`$ and $`\psi _\text{c}`$ to be the restriction of $`\phi `$ and $`\psi `$ to $`M_\text{c}^+`$ respectively, giving you two faithful lower semi-continuous weights on $`M_\text{c}`$.
By equation (1.7) we get that $`(IJ)W(IJ)=W^{}`$, implying that $`R((\iota \omega _{v,w})(W))=(\iota \omega _{Jw,Jv})(W)`$ for all $`v,wH`$. It follows that $`R(M_\text{c})=M_\text{c}`$. Define $`R_\text{c}`$ to be the restriction of $`R`$ to $`M_\text{c}`$. Then $`R_\text{c}`$ is a -anti-automorphism of $`M_\text{c}`$ satisfying $`\chi (R_\text{c}_\text{c}R_\text{c})\mathrm{\Delta }_\text{c}=\mathrm{\Delta }_\text{c}R_\text{c}`$. It is also clear that $`\psi _\text{c}=\phi _\text{c}R_\text{c}`$.
For $`a,b𝒩_\psi `$ and $`c𝒩_\phi `$, we have that
$$(\psi \iota )(\mathrm{\Delta }(b^{}c)(a1))=R\left((\iota \phi )((1R(a))\mathrm{\Delta }(R(c)R(b)^{}))\right)=R\left((\iota \omega _{\mathrm{\Lambda }(R(c)R(b)^{}),\mathrm{\Lambda }(R(a)^{})})(W^{})\right),$$
which implies that $`M_\text{c}=[(\psi \iota )(\mathrm{\Delta }(b^{}c)(a1))a,b𝒩_\psi ,c𝒩_\phi ]`$.
We know that $`(\psi \iota )(\mathrm{\Delta }(b^{}c)(a1))𝒩_\phi `$ and thus $`(\psi \iota )(\mathrm{\Delta }(b^{}c)(a1))𝒩_{\phi _\text{c}}`$ for all $`a,b𝒩_\psi `$ and $`c𝒩_\phi `$. It follows that $`\phi _\text{c}`$ is densely defined.
Define $`\mathrm{\Lambda }_\text{c}`$ to be the restriction of $`\mathrm{\Lambda }`$ to $`𝒩_{\phi _\text{c}}`$. Equation (1.2) guarantees that $`\mathrm{\Lambda }_\text{c}(𝒩_{\phi _\text{c}})`$ is dense in $`H`$. Therefore $`(H,\iota ,\mathrm{\Lambda }_\text{c})`$ is a GNS-construction for $`\phi _\text{c}`$.
Equation (1.8) tells us that $`(N^1)W=W(N^1)`$ implying that
$$\tau _t((\iota \omega _{v,w})(W))=(\iota \omega _{^{it}v,^{it}w})(W)$$
(1.9)
for all $`v,wH`$ and $`t`$. Hence $`\tau _t(M_\text{c})=M_\text{c}`$ for all $`t`$. Define the one-parameter group $`\tau ^\text{c}`$ on $`M_\text{c}`$ by setting $`\tau _t^\text{c}=\tau _t_{M_\text{c}}`$ for all $`t`$. Notice that equation (1.9) implies that $`\tau ^\text{c}`$ is norm continuous.
Since $`\mathrm{\Delta }_\text{c}(M_\text{c})(1M_\text{c})`$ is a dense subset of $`M_\text{c}_\text{c}M_\text{c}`$, we get that $`M_\text{c}=[(\iota \omega )\mathrm{\Delta }(x)\omega B(H)_{},xM_\text{c}]`$. Equation (1.6) implies for all $`t`$, $`\omega B(H)_{}`$ and $`xM_\text{c}`$ that
$$\sigma _t((\iota \omega )\mathrm{\Delta }(x))=(\iota \omega \sigma _t^{})\mathrm{\Delta }(\tau _t^\text{c}(x)).$$
(1.10)
Therefore $`\sigma _t(M_\text{c})=M_\text{c}`$ for all $`t`$ and we can define a one parameter group $`\sigma ^\text{c}`$ on $`M_\text{c}`$ by setting $`\sigma _t^\text{c}=\sigma _t_{M_\text{c}}`$ for all $`t`$. Equation (1.10) implies that $`\sigma ^\text{c}`$ is norm continuous.
By now it is clear that $`\phi _\text{c}`$ is a KMS-weight on $`M_\text{c}`$ (in the C$`^{}`$-algebraic sense) with $`\sigma ^\text{c}`$ as its modular group. Because $`\psi =\phi R`$, we also get that $`\psi _\text{c}`$ is a KMS-weight on $`M_\text{c}`$.
Take $`\omega (M_\text{c})_+^{}`$ and $`x_{\phi _\text{c}}^+`$. Choose $`\eta B(H)_{}^+`$. On the C$`^{}`$-algebra $`M_\text{c}`$ we can make a GNS-construction for the positive functional $`\omega `$. This way we obtain a Hilbert space $`K`$, a non-degenerate representation $`\pi `$ of $`M_\text{c}`$ on $`K`$ and a (cyclic) vector $`vK`$ such that $`\omega =\omega _{v,v}\pi `$. By theorem 1.5 of we know that $`W`$ belongs to the multiplier algebra of $`M_\text{c}B_0(H)`$, where $`B_0(H)`$ denotes the C$`^{}`$-algebra of compact operators on $`H`$. Hence the unitary $`U`$ defined by $`U:=(\pi _\text{c}\iota )(W)`$ belongs to $`B(K)B(H)`$. Define $`\theta B(H)_{}^+`$ by setting $`\theta (x)=(\omega _{v,v}\eta )(U^{}(1x)U)`$ for all $`xB(H)`$. Then
$`(\eta \iota )\mathrm{\Delta }\left((\omega _\text{c}\iota )(\mathrm{\Delta }_\text{c}(x))\right)`$ $`=(\eta \iota )\left((\omega _\text{c}\iota _\text{c}\iota )((\mathrm{\Delta }_\text{c}_\text{c}\iota )\mathrm{\Delta }_\text{c}(x))\right)`$
$`=(\eta \iota )((\omega _{v,v}\iota \iota )(U_{12}^{}\mathrm{\Delta }(x)_{23}U_{12}))=(\theta \iota )\mathrm{\Delta }(x).`$
Therefore the left invariance of $`\phi `$ implies that $`(\eta \iota )\mathrm{\Delta }\left((\omega _\text{c}\iota )(\mathrm{\Delta }_\text{c}(x))\right)`$ belongs to $`_\phi ^+`$ and
$$\phi \left((\eta \iota )\mathrm{\Delta }\left((\omega _\text{c}\iota )(\mathrm{\Delta }_\text{c}(x))\right)\right)=\theta (1)\phi (x)=\omega (1)\eta (1)\phi _\text{c}(x).$$
(1.11)
Translating proposition 5.15 of to the von Neumann algebra setting, we now conclude that $`(\omega _\text{c}\iota )(\mathrm{\Delta }_\text{c}(x))`$ belongs to $`_\phi ^+`$ and therefore to $`_{\phi _\text{c}}^+`$.
Taking $`\eta B(H)_{}`$ such that $`\eta (1)=1`$, equation (1.11) and the left invariance of $`\phi `$ imply that
$$\phi _\text{c}\left((\omega _\text{c}\iota )(\mathrm{\Delta }_\text{c}(x))\right)=\phi \left((\omega _\text{c}\iota )(\mathrm{\Delta }_\text{c}(x))\right)=\phi \left((\eta \iota )\mathrm{\Delta }\left((\omega _\text{c}\iota )(\mathrm{\Delta }_\text{c}(x))\right)\right)=\omega (1)\phi _\text{c}(x).$$
So we have proven that $`\phi _\text{c}`$ is left invariant in the sense of definition 2.2 of . Because $`\chi (R_\text{c}_\text{c}R_\text{c})\mathrm{\Delta }=\mathrm{\Delta }R_\text{c}`$ and $`\psi _\text{c}=\phi _\text{c}R_\text{c}`$ we also get that $`\psi _\text{c}`$ is right invariant. From all this we conclude that $`(M_\text{c},\mathrm{\Delta }_\text{c})`$ is a reduced C$`^{}`$-algebraic quantum group. ∎
The GNS-construction $`(H,\iota ,\mathrm{\Lambda }_\text{c})`$ for $`\phi _\text{c}`$ was obtained by letting $`\mathrm{\Lambda }_\text{c}`$ be the restriction of $`\mathrm{\Lambda }`$ to $`𝒩_{\phi _\text{c}}`$. By the definitions introduced at the end of section 4 in , it is clear that this implies that $`W`$ is the multiplicative unitary of $`(M_\text{c},\mathrm{\Delta }_\text{c})`$ in this GNS-construction $`(H,\iota ,\mathrm{\Lambda }_\text{c})`$.
Since $`\sigma _t^\text{c}`$ and $`\tau _t^\text{c}`$ are restrictions of $`\sigma _t`$ and $`\tau _t`$ respectively, it is clear that $`(\tau _t^\text{c}\sigma _t^\text{c})\mathrm{\Delta }=\mathrm{\Delta }\sigma _t^\text{c}`$ for all $`t`$, and so the density conditions imply that $`\tau ^\text{c}`$ is the scaling group of $`(M_\text{c},\mathrm{\Delta }_\text{c})`$. It follows that $`\nu `$ is also the scaling constant of $`(M_\text{c},\mathrm{\Delta }_\text{c})`$. Letting $`S_\text{c}`$ denote the antipode of $`(M_\text{c},\mathrm{\Delta }_\text{c})`$, proposition 5.33 of and its von Neumann algebraic counterpart imply that $`S_\text{c}S`$. Since $`\tau ^\text{c}`$ is the scaling group of $`(M_\text{c},\mathrm{\Delta }_\text{c})`$ and $`R_\text{c}`$ was obtained by restricting $`R`$ to $`M_\text{c}`$, this implies that $`R_\text{c}`$ is the unitary antipode of $`(M_\text{c},\mathrm{\Delta }_\text{c})`$.
Following , we associate to the reduced C$`^{}`$-algebraic quantum group $`(M_\text{c},\mathrm{\Delta }_\text{c})`$ the von Neumman algebraic quantum group $`(\stackrel{~}{M}_\text{c},\stackrel{~}{\mathrm{\Delta }}_\text{c})`$ by letting $`\stackrel{~}{M}_\text{c}`$ be the $`\sigma `$-strong closure of $`M_\text{c}`$ and defining $`\stackrel{~}{\mathrm{\Delta }}_\text{c}`$ to be the unique normal -homomorphism from $`\stackrel{~}{M}_\text{c}`$ to $`\stackrel{~}{M}_\text{c}\stackrel{~}{M}_\text{c}`$ extending $`\mathrm{\Delta }_\text{c}`$. It follows from proposition 1.5 that $`(\stackrel{~}{M}_\text{c},\stackrel{~}{\mathrm{\Delta }}_\text{c})=(M,\mathrm{\Delta })`$. We get similar results for the extensions of the Haar weights, their modular groups, the scaling group, the unitary antipode and the antipode itself.
## 2. Commutation relations and related matters
In this section we establish some useful technical properties about von Neumann algebraic quantum groups that are often used when working in the operator algebra approach to quantum groups. We start of by implementing the scaling groups and unitary antipodes. Then we prove some results concerning the dual and end by formulating some commutation relations.
We still have fixed a von Neumann algebraic quantum group $`(M,\mathrm{\Delta })`$ and we use the notations introduced in section 1.
###### Proposition 2.1.
The following properties hold
$$\begin{array}{ccccccc}\tau _t(x)\hfill & =& \widehat{}^{it}x\widehat{}^{it}\text{ for all }t,xM\hfill & & R(x)\hfill & =& \widehat{J}x^{}\widehat{J}\text{ for all }xM\hfill \\ \widehat{\tau }_t(x)\hfill & =& ^{it}x^{it}\text{ for all }t,x\widehat{M}\hfill & & \widehat{R}(x)\hfill & =& Jx^{}J\text{ for all }x\widehat{M}\hfill \end{array}$$
###### Proof.
By propositions 8.17 and 8.25 of , we know that $`\widehat{R}(x)=Jx^{}J`$ for all $`x\widehat{M}`$. Therefore the Pontryagin duality theorem guarantees that also $`R(x)=\widehat{J}x^{}\widehat{J}`$ for all $`xM`$.
Choose $`xM`$. By lemma 8.8 and proposition 8.9 of we know that $`\widehat{}^{it}=P^{it}J\delta ^{it}J`$, and so we get that $`\widehat{}^{it}x\widehat{}^{it}=P^{it}J\delta ^{it}JxJ\delta ^{it}JP^{it}`$. But Tomita-Takesaki theory tells us that $`J\delta ^{it}J`$ belongs to $`M^{}`$, implying that $`\widehat{}^{it}x\widehat{}^{it}=P^{it}xP^{it}=\tau _t(x)`$.
Pontryagin duality allows us to conclude that $`\widehat{\tau }^{it}(x)=^{it}x^{it}`$ for all $`t`$ and $`x\widehat{M}`$. ∎
We have that $`(R\widehat{R})(W)=W`$ and $`(\tau _t\widehat{\tau }_t)(W)=W`$ for all $`t`$ (see the remarks before propositions 8.18 and 8.25 of ). Hence the next result.
###### Corollary 2.2.
We have the following commutation relations:
$$W(\widehat{})=(\widehat{})W\text{and}W(\widehat{J}J)=(\widehat{J}J)W^{}.$$
Notice that for the same reasons, $`W(P)=(P)W`$ and $`W(\widehat{}P)=(\widehat{}P)W`$.
In the next part, we complete the picture of the dual. For this reason, let us introduce a natural -algebra inside $`M_{}`$.
###### Definition 2.3.
Define the subspace $`M_{}^{\mathrm{}}`$ of $`M_{}`$ as
$$M_{}^{\mathrm{}}=\{\omega M_{}\theta M_{}:\theta (x)=\overline{\omega }(S(x))\text{ for all }xD(S)\}.$$
We define the antilinear mapping $`.^{}:M_{}^{\mathrm{}}M_{}^{\mathrm{}}`$ such that $`\omega ^{}(x)=\overline{\omega }(S(x))`$ for all $`\omega M_{}^{\mathrm{}}`$ and $`xD(S)`$. Then $`M_{}^{\mathrm{}}`$ is a subalgebra of $`M_{}`$ and becomes a -algebra under the operation $`.^{}`$ .
If $`xD(S)`$, then $`S(x)^{}D(S)`$ and $`S(S(x)^{})^{}=x`$ (which follows from the corresponding property for $`\tau _{\frac{i}{2}}`$). This implies that $`.^{}`$ is an involution on $`M_{}^{\mathrm{}}`$. If $`\omega `$ and $`\eta `$ are elements in $`M_{}`$ such that $`\omega S`$ and $`\eta S`$ are bounded and their unique continuous linear extensions $`(\omega S)\stackrel{~}{\text{}}`$ and $`(\eta S)\stackrel{~}{\text{}}`$ belong to $`M_{}`$, then $`(\omega \eta )S`$ is bounded and $`(\eta S)\stackrel{~}{\text{}}(\omega S)\stackrel{~}{\text{}}`$ is its unique continuous linear extension (cfr lemma 5.25 of ). From this, it follows that $`M_{}^{\mathrm{}}`$ is a subalgebra of $`M_{}`$ and that $`.^{}`$ is antimultiplicative. Notice that, since $`S`$ can be unbounded, $`M_{}^{\mathrm{}}`$ can be strictly smaller than $`M_{}`$.
For $`\omega B(H)_{}`$, the element $`(\iota \omega )(W)`$ belongs to $`D(S)`$ and $`S((\iota \omega )(W))=(\iota \omega )(W^{})`$. The space $`\{(\iota \omega )(W)\omega B(H)_{}\}`$ is moreover a $`\sigma `$-strong core for $`S`$. (see proposition 8.3 of ). We use this characterization of $`S`$ to prove the next result.
###### Proposition 2.4.
The following holds :
* $`M_{}^{\mathrm{}}=\{\omega M_{}\theta M_{}:\lambda (\omega )^{}=\lambda (\theta )\}`$.
* $`\lambda (\omega )^{}=\lambda (\omega ^{})`$ for all $`\omega M_{}^{\mathrm{}}`$.
###### Proof.
Take $`\omega M_{}`$. Then we have for all $`\eta B(H)_{}`$ that
$$\overline{\omega }\left(S((\iota \eta )(W))\right)=\overline{\omega }((\iota \eta )(W^{}))=\eta ((\omega \iota )(W)^{})=\eta (\lambda (\omega )^{}).$$
(2.1)
If $`\omega M_{}^{\mathrm{}}`$ then the formula above implies for all $`\eta B(H)_{}`$ that
$$\eta (\lambda (\omega ^{}))=\eta ((\omega ^{}\iota )(W))=\omega ^{}((\iota \eta )(W))=\eta (\lambda (\omega )^{}),$$
and hence $`\lambda (\omega ^{})=\lambda (\omega )^{}`$.
Now suppose that there exists $`\theta M_{}`$ such that $`\lambda (\omega )^{}=\lambda (\theta )`$. By formula (2.1) above, we get for all $`\eta B(H)_{}`$ that
$$\overline{\omega }\left(S((\iota \eta )(W))\right)=\eta ((\theta \iota )(W))=\theta ((\iota \eta )(W)).$$
Because such elements $`(\iota \eta )(W)`$ form a $`\sigma `$-strong core for $`S`$, we get $`\overline{\omega }(S(x))=\theta (x)`$ for all $`xD(S)`$. So $`\omega `$ belongs to $`M_{}^{\mathrm{}}`$. ∎
It is easy to prove that $`M_{}^{\mathrm{}}`$ is dense in $`M_{}`$ implying that the -algebra $`\lambda (M_{}^{\mathrm{}})`$ is $`\sigma `$-strong dense in $`\widehat{M}`$. For later purposes, we will need a result which gives a little bit more information.
###### Lemma 2.5.
The spaces $`M_{}^{\mathrm{}}`$ and $`(M_{}^{\mathrm{}})^{}`$ are dense in $`M_{}`$ and $`\lambda (M_{}^{\mathrm{}})`$ is a $`\sigma `$-strong–norm core for $`\widehat{\mathrm{\Lambda }}`$.
###### Proof.
Consider $`\omega `$. For every $`n`$ and $`z`$, we define $`\omega (n,z)M_{}`$ as
$$\omega (n,z)=\frac{n}{\sqrt{\pi }}\mathrm{exp}(n^2(t+z)^2)\omega \tau _t𝑑t.$$
So we have for $`xD(S)`$ that $`xD(\tau _{\frac{i}{2}})`$ and thus
$`\overline{\omega (n,z)}(S(x))`$ $`=`$ $`{\displaystyle \frac{n}{\sqrt{\pi }}}{\displaystyle \mathrm{exp}(n^2(t+\overline{z})^2)\overline{\omega }\left(\tau _t(S(x))\right)𝑑t}`$
$`=`$ $`{\displaystyle \frac{n}{\sqrt{\pi }}}{\displaystyle \mathrm{exp}(n^2(t+\overline{z})^2)\overline{\omega }\left(R(\tau _{t\frac{i}{2}}(x))\right)𝑑t}`$
$`=`$ $`{\displaystyle \frac{n}{\sqrt{\pi }}}{\displaystyle \mathrm{exp}(n^2(t+\frac{i}{2}+\overline{z})^2)\overline{\omega }\left(R(\tau _t(x))\right)𝑑t},`$
from which we conclude that $`\omega (n,z)M_{}^{\mathrm{}}`$ and
$$\omega (n,z)^{}=\frac{n}{\sqrt{\pi }}\mathrm{exp}(n^2(t+\frac{i}{2}+\overline{z})^2)\overline{\omega }R\tau _t𝑑t.$$
(2.2)
It is easy to check that for every $`t`$ we have $`\omega \tau _t`$ and $`\xi (\omega \tau _t)=\nu ^{\frac{t}{2}}P^{it}\xi (\omega )`$. Therefore the closedness of the mapping $`\eta \xi (\eta )`$ implies that $`\omega (n,z)`$ and
$$\xi (\omega (n,z))=\frac{n}{\sqrt{\pi }}\mathrm{exp}(n^2(t+z)^2)\nu ^{\frac{t}{2}}P^{it}\xi (\omega )𝑑t.$$
(2.3)
* Let $`\omega `$. Then we have for every $`n`$ that $`\omega (n,0)M_{}^{\mathrm{}}`$. Clearly, $`(\omega (n,0))_{n=1}^{\mathrm{}}`$ converges to $`\omega `$. Equation (2.3) implies that $`\left(\xi (\omega (n,0))\right)_{n=1}^{\mathrm{}}`$ converges to $`\xi (\omega )`$. In other words, $`\left(\widehat{\mathrm{\Lambda }}(\lambda (\omega (n,0)))\right)_{n=1}^{\mathrm{}}`$ converges to $`\widehat{\mathrm{\Lambda }}(\lambda (\omega ))`$. Since $``$ is dense in $`M_{}`$ and $`\lambda ()`$ is a core for $`\widehat{\mathrm{\Lambda }}`$, we conclude that $`M_{}^{\mathrm{}}`$ is dense in $`M_{}`$ and that $`\lambda (M_{}^{\mathrm{}})`$ is a $`\sigma `$-strong–norm core for $`\widehat{\mathrm{\Lambda }}`$.
* Let $`\omega `$. Then we have for every $`n`$ that $`\omega (n,\frac{i}{2})M_{}^{\mathrm{}}`$ and
$$\omega (n,\frac{i}{2})^{}=\frac{n}{\sqrt{\pi }}\mathrm{exp}(n^2t^2)\overline{\omega }R\tau _t𝑑t$$
by equation (2.2). So $`(\omega (n,\frac{i}{2})^{})_{n=1}^{\mathrm{}}`$ converges to $`\overline{\omega }R`$. From this all, we conclude that $`(M_{}^{\mathrm{}})^{}`$ is dense in $`M_{}`$.
###### Proposition 2.6.
Define $`^{\mathrm{}}=\{xM_{}^{\mathrm{}}x^{}\}`$. Then $`^{\mathrm{}}`$ is a -subalgebra of $`M_{}^{\mathrm{}}`$ such that $`^{\mathrm{}}`$ is dense in $`M_{}`$ and $`\lambda (^{\mathrm{}})`$ is a $`\sigma `$-strong–norm core for $`\widehat{\mathrm{\Lambda }}`$.
###### Proof.
It is clear that $`^{\mathrm{}}`$ is a -subalgebra of $`M_{}^{\mathrm{}}`$. Because $``$ is a left ideal in $`M_{}`$, we get that $`(M_{}^{\mathrm{}})^{}(M_{}^{\mathrm{}})^{\mathrm{}}`$. Thus in order to prove that $`^{\mathrm{}}`$ is dense in $`M_{}`$, it is by the previous lemma enough to prove that $`(M_{})^2`$ is dense in $`M_{}`$. But we have for all $`vH`$ with $`v=1`$, $`w_1,w_2H`$ and $`xM`$ that
$$\mathrm{\Delta }(x)W^{}(vw_1),W^{}(vw_2)=xw_1,w_2,$$
which easily implies that $`(M_{})^2`$ is dense in $`M_{}`$. Hence $`^{\mathrm{}}`$ is dense in $`M_{}`$.
Since $`IM_{}^{\mathrm{}}`$ is dense in $`M_{}`$, $`1`$ belongs to the $`\sigma `$-strong closure of $`\lambda (IM_{}^{\mathrm{}})^{}`$. Combining this with the fact that $`\lambda (IM_{}^{\mathrm{}})`$ is a $`\sigma `$-strong–norm core for $`\widehat{\mathrm{\Lambda }}`$ and the inclusion $`\lambda (M_{}^{\mathrm{}})^{}\lambda (M_{}^{\mathrm{}})\lambda (^{\mathrm{}})`$, we conclude that $`\lambda (^{\mathrm{}})`$ is a $`\sigma `$-strong–norm core for $`\widehat{\mathrm{\Lambda }}`$. ∎
Let us connect the modular objects of $`\widehat{\phi }`$ to objects already constructed on the level of $`(M,\mathrm{\Delta })`$.
We know that the operators $`P`$ and $`J\delta J`$ strongly commute and that $`\widehat{}^{it}=P^{it}J\delta ^{it}J`$ for $`t`$ (see lemma 8.8 and proposition 8.9 of ). Also notice that this implies for every $`a𝒩_\phi `$ that $`\tau _t(a)\delta ^{it}`$ belongs to $`𝒩_\phi `$ and $`\widehat{}^{it}\mathrm{\Lambda }(a)=\mathrm{\Lambda }(\tau _t(a)\delta ^{it})`$.
Put $`\widehat{T}=\widehat{J}\widehat{}^{\frac{1}{2}}`$. So $`\widehat{\mathrm{\Lambda }}(𝒩_{\widehat{\phi }}𝒩_{\widehat{\phi }}^{})`$ is a core for $`\widehat{T}`$ and $`\widehat{T}\widehat{\mathrm{\Lambda }}(x)=\widehat{\mathrm{\Lambda }}(x^{})`$ for all $`x𝒩_{\widehat{\phi }}𝒩_{\widehat{\phi }}^{}`$.
###### Lemma 2.7.
The set $`\widehat{\mathrm{\Lambda }}(\lambda (^{\mathrm{}}))`$ is a core for $`\widehat{T}`$.
###### Proof.
Since $`\widehat{T}=\widehat{J}\widehat{}^{\frac{1}{2}}`$, the definition of $`\widehat{T}`$ gives clearly that $`\widehat{\mathrm{\Lambda }}(\lambda (^{\mathrm{}}))D(\widehat{}^{\frac{1}{2}})`$. From proposition 2.6, we know that $`\widehat{\mathrm{\Lambda }}(\lambda (^{\mathrm{}}))`$ is a dense subspace in $`H`$.
We now use the notation $`\rho _t`$ as it was introduced in notation 8.7 of . For every $`\omega M_{}`$ we denote by $`\rho _t(\omega )`$ the element in $`M_{}`$ defined by $`\rho _t(\omega )(x)=\omega (\delta ^{it}\tau _t(x))`$. Then $`\rho _t()=`$ and $`\xi (\rho _t(\omega ))=\widehat{}^{it}\xi (\omega )`$ for all $`\omega `$ and $`t`$. If $`\omega M_{}^{\mathrm{}}`$ and $`t`$, it is not so difficult to check that $`\rho _t(\omega )M_{}^{\mathrm{}}`$ and $`\rho _t(\omega )^{}=\rho _t(\omega ^{})`$. It follows that $`\rho _t(^{\mathrm{}})=^{\mathrm{}}`$ for all $`t`$, hence $`\widehat{\sigma }_t(\lambda (^{\mathrm{}}))=\lambda (\rho _t(^{\mathrm{}}))=\lambda (^{\mathrm{}})`$ for all $`t`$.
Therefore $`\widehat{}^{it}\widehat{\mathrm{\Lambda }}(\lambda (^{\mathrm{}}))=\widehat{\mathrm{\Lambda }}(\lambda (^{\mathrm{}}))`$ for all $`t`$. We conclude from all this that $`\widehat{\mathrm{\Lambda }}(\lambda (^{\mathrm{}}))`$ is a core for $`\widehat{}^{\frac{1}{2}}`$ (see e.g. corollary 1.21 of ) and the lemma follows. ∎
###### Proposition 2.8.
Consider $`x𝒩_\phi D(S^1)`$ such that $`S^1(x)^{}𝒩_\phi `$. Then $`\mathrm{\Lambda }(x)D(\widehat{T}^{})`$ and $`\widehat{T}^{}\mathrm{\Lambda }(x)=\mathrm{\Lambda }(S^1(x)^{})`$.
###### Proof.
Choose $`\theta ^{\mathrm{}}`$. Then
$`\widehat{T}\widehat{\mathrm{\Lambda }}(\lambda (\theta )),\mathrm{\Lambda }(x)=\widehat{\mathrm{\Lambda }}(\lambda (\theta )^{}),\mathrm{\Lambda }(x)=\widehat{\mathrm{\Lambda }}(\lambda (\theta ^{})),\mathrm{\Lambda }(x)=\xi (\theta ^{}),\mathrm{\Lambda }(x).`$
Therefore the definition of $`\xi (\theta ^{})`$ and $`\theta ^{}`$ imply that
$`\widehat{T}\widehat{\mathrm{\Lambda }}(\lambda (\theta )),\mathrm{\Lambda }(x)=\theta ^{}(x^{})=\overline{\theta }(S(x^{}))=\overline{\theta (S^1(x))}=\overline{\xi (\theta ),\mathrm{\Lambda }(S^1(x)^{})}`$
$`=\mathrm{\Lambda }(S^1(x)^{}),\xi (\theta )=\mathrm{\Lambda }(S^1(x)^{}),\widehat{\mathrm{\Lambda }}(\lambda (\theta )).`$
Thus the previous lemma implies that $`\mathrm{\Lambda }(x)`$ belongs to $`D(\widehat{T}^{})`$ and $`\widehat{T}^{}\mathrm{\Lambda }(x)=\mathrm{\Lambda }(S^1(x)^{})`$. ∎
This proposition allows us to establish easily a connection between $`G`$ and $`\widehat{T}`$. Recall that the operators $`G`$, $`N`$ and $`I`$ were introduced in proposition 1.3 and notation 1.4.
###### Corollary 2.9.
We have that $`\widehat{T}^{}=G`$, $`\widehat{}=N^1`$ and $`\widehat{J}=I`$.
###### Proof.
Using proposition 1.3 and the strong left invariance of $`\psi `$ (see proposition 5.24 of ), the previous result implies easily $`G\widehat{T}^{}`$.
Define the subspace $`C`$ of $`D(G)`$ as $`C=\mathrm{\Lambda }\left((\psi \iota )(\mathrm{\Delta }(y^{})(x1))\right)x,y𝒩_\phi ^{}𝒩_\psi `$.
Let $`t`$. Remember that $`\widehat{}^{it}\mathrm{\Lambda }(a)=\mathrm{\Lambda }(\tau _t(a)\delta ^{it})`$ for all $`a𝒩_\phi `$.
Choose $`x,y𝒩_\phi ^{}𝒩_\psi `$. Then $`\delta ^{it}\tau _t(x)`$ and $`\delta ^{it}\tau _t(y)`$ belong to $`𝒩_\phi ^{}𝒩_\psi `$ and
$$\tau _t((\psi \iota )(\mathrm{\Delta }(y^{})(x1)))\delta ^{it}=\nu ^t(\psi \iota )(\mathrm{\Delta }((\delta ^{it}\tau _t(y))^{})(\delta ^{it}\tau _t(x)1)).$$
Therefore the element $`\widehat{}^{it}\mathrm{\Lambda }\left((\psi \iota )(\mathrm{\Delta }(y^{})(x1))\right)=\mathrm{\Lambda }(\tau _t((\psi \iota )(\mathrm{\Delta }(y^{})(x1))\left)\delta ^{it}\right)`$ belongs to $`C`$.
We conclude that $`C`$ is a dense subspace of $`D(\widehat{}^{\frac{1}{2}})`$, invariant under the family of operators $`\widehat{}^{it}`$ $`(t)`$. It follows that $`C`$ is a core for $`\widehat{}^{\frac{1}{2}}`$ and thus a core for $`\widehat{T}^{}=\widehat{J}\widehat{}^{\frac{1}{2}}`$. Combining this with the fact that $`G\widehat{T}^{}`$, we conclude that $`G=\widehat{T}^{}`$. Now the uniqueness of the polar decomposition implies that $`\widehat{}=N^1`$ and $`\widehat{J}=I`$. ∎
Combining the previous corollary with proposition 1.3 we get the following.
###### Corollary 2.10.
The set
$$\{\mathrm{\Lambda }(x)x𝒩_\phi D(S^1)\text{ such that }S^1(x)^{}𝒩_\phi \}$$
is a core for $`\widehat{T}^{}`$.
Recall that we introduced the GNS-construcion $`(H,\iota ,\mathrm{\Gamma })`$ for $`\psi `$ by considering $`\psi `$ as $`\phi _\delta `$ and setting $`\mathrm{\Gamma }=\mathrm{\Lambda }_\delta `$. But $`\psi `$ is by definition equal to $`\phi R`$. It turns out that $`\widehat{J}`$ connects both pictures of $`\psi `$:
###### Proposition 2.11.
We have for all $`x𝒩_\psi `$ that $`\widehat{J}\mathrm{\Gamma }(x)=\mathrm{\Lambda }(R(x)^{})`$.
###### Proof.
Define the anti-unitary $`U:HH`$ such that $`U\mathrm{\Gamma }(x)=\mathrm{\Lambda }(R(x)^{})`$ for $`x𝒩_\psi `$. Choose $`a𝒩_\phi `$ such that $`aD(S^1)`$ and $`S^1(a)^{}𝒩_\phi `$.
For $`n`$, we define $`e_nM`$ such that $`e_n=\frac{n}{\sqrt{\pi }}\mathrm{exp}(n^2t^2)\delta ^{it}𝑑t`$. Remember that $`e_n`$ is analytic with respect to $`\sigma `$ and $`\sigma ^{}`$, implying that $`𝒩_\phi e_n𝒩_\phi `$ and $`𝒩_\psi e_n𝒩_\psi `$
Since $`\tau _s(\delta )=\delta `$ we see that $`\tau _s(e_n)=e_n`$ for $`s`$, hence $`e_nD(\tau _{\frac{i}{2}})`$ and $`\tau _{\frac{i}{2}}(e_n)=e_n`$. By assumption, $`aD(\tau _{\frac{i}{2}})`$, so $`ae_nD(\tau _{\frac{i}{2}})`$ and $`\tau _{\frac{i}{2}}(ae_n)=\tau _{\frac{i}{2}}(a)e_n`$. Hence $`\tau _{\frac{i}{2}}(ae_n)\delta ^{\frac{1}{2}}`$ is a bounded operator and its closure equals $`\tau _{\frac{i}{2}}(a)(\delta ^{\frac{1}{2}}e_n)`$.
Define the strongly continuous one-parameter group $`\kappa `$ of isometries of $`M`$ such that $`\kappa _t(x)=\tau _t(x)\delta ^{it}`$ for $`xM`$ and $`t`$. The discussion above implies (see e.g. proposition 4.9 of ) that $`ae_nD(\kappa _{\frac{i}{2}})`$ and $`\kappa _{\frac{i}{2}}(ae_n)=\tau _{\frac{i}{2}}(a)(\delta ^{\frac{1}{2}}e_n)`$.
By assumption $`R(\tau _{\frac{i}{2}}(a))^{}=S^1(a)^{}𝒩_\phi `$, implying that $`\tau _{\frac{i}{2}}(a)`$ belongs to $`𝒩_\psi `$. So we see that $`\kappa _{\frac{i}{2}}(ae_n)\delta ^{\frac{1}{2}}`$ is a bounded operator and that its closure equals $`\tau _{\frac{i}{2}}(a)e_n𝒩_\psi `$. Since $`\mathrm{\Lambda }=\mathrm{\Gamma }_{\delta ^1}`$, this implies that $`\kappa _{\frac{i}{2}}(ae_n)𝒩_\phi `$ and
$$\mathrm{\Lambda }(\kappa _{\frac{i}{2}}(ae_n))=\mathrm{\Gamma }(\kappa _{\frac{i}{2}}(ae_n)\delta ^{\frac{1}{2}})=\mathrm{\Gamma }(\tau _{\frac{i}{2}}(a)e_n).$$
We know that we have for every $`x𝒩_\phi `$ that $`\kappa _t(x)𝒩_\phi `$ and $`\mathrm{\Lambda }(\kappa _t(x))=\widehat{}^{it}\mathrm{\Lambda }(x)`$. Since $`ae_n𝒩_\phi `$ and $`\kappa _{\frac{i}{2}}(ae_n)𝒩_\phi `$, we conclude (see e.g. proposition 4.4 of ) that $`\mathrm{\Lambda }(ae_n)D(\widehat{}^{\frac{1}{2}})`$ and
$$\widehat{}^{\frac{1}{2}}\mathrm{\Lambda }(ae_n)=\mathrm{\Lambda }(\kappa _{\frac{i}{2}}(ae_n))=\mathrm{\Gamma }(\tau _{\frac{i}{2}}(a)e_n).$$
Since $`(\mathrm{\Lambda }(ae_n))_{n=1}^{\mathrm{}}`$ converges to $`\mathrm{\Lambda }(a)`$ and $`(\mathrm{\Gamma }(\tau _{\frac{i}{2}}(a)e_n))_{n=1}^{\mathrm{}}`$ converges to $`\mathrm{\Gamma }(\tau _{\frac{i}{2}}(a))`$, the closedness of $`\widehat{}^{\frac{1}{2}}`$ implies that $`\mathrm{\Lambda }(a)D(\widehat{}^{\frac{1}{2}})`$ and
$$\widehat{}^{\frac{1}{2}}\mathrm{\Lambda }(a)=\mathrm{\Gamma }(\tau _{\frac{i}{2}}(a)).$$
Consequently
$$U\widehat{}^{\frac{1}{2}}\mathrm{\Lambda }(a)=U\mathrm{\Gamma }(\tau _{\frac{i}{2}}(a))=\mathrm{\Lambda }(R(\tau _{\frac{i}{2}}(a))^{})=\mathrm{\Lambda }(S^1(a)^{})=\widehat{T}^{}\mathrm{\Lambda }(a)=\widehat{J}\widehat{}^{\frac{1}{2}}\mathrm{\Lambda }(a).$$
Since such elements $`\mathrm{\Lambda }(a)`$ form a core for $`\widehat{}^{\frac{1}{2}}=\widehat{J}\widehat{T}^{}`$, such elements $`\widehat{}^{\frac{1}{2}}\mathrm{\Lambda }(a)`$ form a dense subspace of $`H`$. Therefore $`\widehat{J}=U`$ and we are done. ∎
The equality in the next corollary is a slight adaptation of corollary 3.6.2.(iv) of .
###### Corollary 2.12.
We have that $`\widehat{J}J=\nu ^{\frac{i}{4}}J\widehat{J}`$.
###### Proof.
As already mentioned in section 1, $`J^{}:=\nu ^{\frac{i}{4}}J`$ is the modular conjugation of $`\psi `$ in the GNS-construction $`(H,\iota ,\mathrm{\Gamma })`$. Choose $`x𝒩_\psi D(\sigma _{\frac{i}{2}}^{})`$. Since $`\sigma _tR=R\sigma _t^{}`$ for all $`t`$, we get that $`R(x)D(\sigma _{\frac{i}{2}})`$ and $`\sigma _{\frac{i}{2}}(R(x))=R(\sigma _{\frac{i}{2}}^{}(x))`$. Hence $`R(x)^{}𝒩_\phi D(\sigma _{\frac{i}{2}})`$ and $`\sigma _{\frac{i}{2}}(R(x)^{})=R(\sigma _{\frac{i}{2}}^{}(x))^{}`$. Combining this with the previous proposition and the definition of the modular conjugation, we get that
$`\widehat{J}J\mathrm{\Gamma }(x)`$ $`=`$ $`\nu ^{\frac{i}{4}}\widehat{J}J^{}\mathrm{\Gamma }(x)=\nu ^{\frac{i}{4}}\widehat{J}\mathrm{\Gamma }(\sigma _{\frac{i}{2}}^{}(x)^{})=\nu ^{\frac{i}{4}}\mathrm{\Lambda }(R(\sigma _{\frac{i}{2}}^{}(x)^{})^{})`$
$`=`$ $`\nu ^{\frac{i}{4}}\mathrm{\Lambda }(\sigma _{\frac{i}{2}}(R(x)^{})^{})=\nu ^{\frac{i}{4}}J\mathrm{\Lambda }(R(x)^{})=\nu ^{\frac{i}{4}}J\widehat{J}\mathrm{\Gamma }(x).`$
Therefore $`\widehat{J}J=\nu ^{\frac{i}{4}}J\widehat{J}`$. ∎
Recall that we denoted by $`\text{ }`$ the modular operator of $`\psi `$ in the GNS-construction $`(H,\iota ,\mathrm{\Gamma })`$ and by $`\widehat{\text{ }}`$ the modular operator of $`\widehat{\psi }`$ in the GNS-construction $`(H,\iota ,\widehat{\mathrm{\Gamma }})`$.
###### Proposition 2.13.
For all $`s,t`$ we have the following commutation relations.
$`\widehat{}^{it}^{is}=\nu ^{ist}^{is}\widehat{}^{it}\text{and}\widehat{\text{ }}^{it}\text{ }^{is}=\nu ^{ist}\text{ }^{is}\widehat{\text{ }}^{it}`$ (2.4)
$`\widehat{}^{it}\text{ }^{is}=\nu ^{ist}\text{ }^{is}\widehat{}^{it}\text{and}^{is}\text{ }^{it}=\text{ }^{it}^{is}`$ (2.5)
$`\widehat{J}\widehat{J}=\text{ },JJ=^1\text{and}J\text{ }J=\text{ }^1`$ (2.6)
$`\widehat{J}P\widehat{J}=P^1\text{and}\widehat{J}\delta \widehat{J}=\delta ^1`$ (2.7)
$`P^{is}^{it}=^{it}P^{is}\text{and}P^{is}\text{ }^{it}=\text{ }^{it}P^{is}`$ (2.8)
$`P^{is}\delta ^{it}=\delta ^{it}P^{is}`$ (2.9)
$`^{is}\delta ^{it}=\nu ^{ist}\delta ^{it}^{is}\text{and}\text{ }^{is}\delta ^{it}=\nu ^{ist}\delta ^{it}\text{ }^{is}`$ (2.10)
$`\widehat{}^{is}\delta ^{it}=\delta ^{it}\widehat{}^{is}\text{and}\widehat{\text{ }}^{is}\delta ^{it}=\delta ^{it}\widehat{\text{ }}^{is}`$ (2.11)
All commutation relations remain true if we remove the $`\widehat{}`$ of $``$,$`\text{ }`$ and $`J`$ if there is one, add a $`\widehat{}`$ to $``$,$`\text{ }`$ and $`J`$ if there is not one, replace $`\nu `$ by $`\nu ^1`$, replace $`\delta `$ by $`\widehat{\delta }`$ and leave $`P`$ unchanged.
###### Proof.
It is easy to check that $`\text{ }^{it}\mathrm{\Lambda }(x)=\nu ^{\frac{t}{2}}\mathrm{\Lambda }(\sigma _t^{}(x))`$ for all $`x𝒩_\phi `$. We already mentioned that $`\widehat{}^{it}\mathrm{\Lambda }(x)=\mathrm{\Lambda }(\tau _t(x)\delta ^{it})`$ for all $`x𝒩_\phi `$ and by definition we have $`^{it}\mathrm{\Lambda }(x)=\mathrm{\Lambda }(\sigma _t(x))`$ and $`P^{it}\mathrm{\Lambda }(x)=\nu ^{\frac{t}{2}}\mathrm{\Lambda }(\tau _t(x))`$ for all $`x𝒩_\phi `$. Because $`\tau _t(\delta )=\delta `$ for all $`t`$ it is easy to verify that $`P^{it}\mathrm{\Gamma }(x)=\nu ^{\frac{t}{2}}\mathrm{\Gamma }(\tau _t(x))`$ for all $`x𝒩_\psi `$ and by definition we have $`\text{ }^{it}\mathrm{\Gamma }(x)=\mathrm{\Gamma }(\sigma _t^{}(x))`$ for all $`x𝒩_\psi `$.
Using that all three one-parametergroups $`\sigma `$, $`\sigma ^{}`$ and $`\tau `$ commute and that $`\sigma _t(\delta ^{is})=\sigma _t^{}(\delta ^{is})=\nu ^{ist}\delta ^{is}`$ and $`\tau _t(\delta ^{is})=\delta ^{is}`$ for all $`s,t`$, it is straightforward to check the first equality in equation (2.4), equation (2.5) and equation (2.8) by applying the operators to an element $`\mathrm{\Lambda }(x)`$ with $`x𝒩_\phi `$. Using proposition 2.11 and the fact that $`\sigma _tR=R\sigma _t^{}`$ we can check the equalities $`\widehat{J}^{it}\widehat{J}=\text{ }^{it}`$ and $`\widehat{J}P^{it}\widehat{J}=P^{it}`$ on a vector $`\mathrm{\Gamma }(x)`$ when $`x𝒩_\psi `$. This gives the first equalities of equations (2.6) and (2.7), and the rest of equation (2.6) follows from modular theory, because $`\nu ^{\frac{i}{4}}J`$ is the modular conjugation of $`\psi `$ in the GNS-construction $`(H,\iota ,\mathrm{\Gamma })`$. By the biduality theorem we also get $`J\widehat{}J=\widehat{\text{ }}`$ and $`JPJ=P^1`$. Because for all $`t`$ we have $`\widehat{}^{it}=P^{it}J\delta ^{it}J`$, we get $`\widehat{\text{ }}^{it}=P^{it}\delta ^{it}`$. This implies that $`\widehat{\text{ }}^{it}\mathrm{\Lambda }(x)=\nu ^{\frac{t}{2}}\mathrm{\Lambda }(\delta ^{it}\tau _t(x))`$ for all $`x𝒩_\phi `$ and we can check then the second equality in equation (2.4) on a vector $`\mathrm{\Lambda }(x)`$ with $`x𝒩_\phi `$.
Because $`R(x)=\widehat{J}x^{}\widehat{J}`$ for all $`xM`$ and $`R(\delta )=\delta ^1`$, we get the second equality of equation (2.7). Equations (2.9) and (2.10) follow because $`P^{is},^{is}`$ and $`\text{ }^{is}`$ implement respectively $`\tau _s,\sigma _s`$ and $`\sigma _s^{}`$ on $`M`$. Also $`\widehat{}^{is}`$ implements $`\tau _s`$ on $`M`$ and this gives the first equality of equation (2.11). Because we already saw that $`\widehat{\text{ }}^{is}=P^{is}\delta ^{is}`$ the second equality of equation (2.11) follows immediately from equation (2.9).
By the biduality theorem we can indeed perform the operation stated in the proposition, because $`P^{it}\widehat{\mathrm{\Lambda }}(y)=\nu ^{\frac{t}{2}}\widehat{\mathrm{\Lambda }}(\widehat{\tau }_t(y))`$ for all $`y𝒩_{\widehat{\phi }}`$ and so in a sense $`\widehat{P}=P`$. ∎
Recall that we introduced the GNS-maps $`\mathrm{\Lambda }`$, $`\mathrm{\Gamma }`$, $`\widehat{\mathrm{\Lambda }}`$ and $`\widehat{\mathrm{\Gamma }}`$ for the weights $`\phi `$, $`\psi `$, $`\widehat{\phi }`$ and $`\widehat{\psi }`$ respectively. Using this we will define now three new multiplicative unitaries on $`HH`$ and relate them with our multiplicative unitary $`W`$ we used all the time. Recall that we already mentioned $`V`$ in section 1.
###### Definition 2.14.
Applying the von Neumann algebraic counterpart of theorem 3.16 in (and its right invariant version) to $`(M,\mathrm{\Delta })`$ and $`(\widehat{M},\widehat{\mathrm{\Delta }})`$, one can define the unitaries $`V`$, $`\widehat{W}`$ and $`\widehat{V}`$ on $`HH`$ by the following formulas.
$`V(\mathrm{\Gamma }(x)\mathrm{\Gamma }(y))`$ $`=(\mathrm{\Gamma }\mathrm{\Gamma })(\mathrm{\Delta }(x)(1y))\text{for all}x,y𝒩_\psi `$
$`\widehat{W}^{}(\widehat{\mathrm{\Lambda }}(x)\widehat{\mathrm{\Lambda }}(y))`$ $`=(\widehat{\mathrm{\Lambda }}\widehat{\mathrm{\Lambda }})(\widehat{\mathrm{\Delta }}(y)(x1))\text{for all}x,y𝒩_{\widehat{\phi }}`$
$`\widehat{V}(\widehat{\mathrm{\Gamma }}(x)\widehat{\mathrm{\Gamma }}(y))`$ $`=(\widehat{\mathrm{\Gamma }}\widehat{\mathrm{\Gamma }})(\widehat{\mathrm{\Delta }}(x)(1y))\text{for all}x,y𝒩_{\widehat{\psi }}.`$
Observe that all the unitaries $`W`$, $`V`$, $`\widehat{W}`$ and $`\widehat{V}`$ satisfy the pentagonal equation
$$W_{12}W_{13}W_{23}=W_{23}W_{12}.$$
This could be checked directly, but it also follows from the pentagonal equation for $`W`$ and the formulas appearing in proposition 2.15.
Almost by definition we have the following.
$$\mathrm{\Delta }(x)=W^{}(1x)W=V(x1)V^{}\text{for all}xM\text{and}$$
$$\widehat{\mathrm{\Delta }}(y)=\widehat{W}^{}(1y)\widehat{W}=\widehat{V}(y1)\widehat{V}^{}\text{for all}y\widehat{M}.$$
The relation between all these multiplicative unitaries is given in the next proposition.
###### Proposition 2.15.
We have the following formulas.
$`\widehat{W}`$ $`=\mathrm{\Sigma }W^{}\mathrm{\Sigma }`$
$`V`$ $`=(\widehat{J}\widehat{J})\mathrm{\Sigma }W^{}\mathrm{\Sigma }(\widehat{J}\widehat{J})`$
$`\widehat{V}`$ $`=(JJ)W(JJ).`$
So we have $`\widehat{W}\widehat{M}M`$, $`V\widehat{M}^{}M`$ and $`\widehat{V}M^{}\widehat{M}`$.
###### Proof.
The first equality follows from proposition 8.16 in . Combining lemma 8.26 in and our proposition 2.11 we get the second equality. Dualizing this we get $`\widehat{V}=(JJ)\mathrm{\Sigma }\widehat{W}^{}\mathrm{\Sigma }(JJ)`$ and this gives the third equality after applying the first one.
The final statement follows from the fact that $`WM\widehat{M}`$, the previous formulas and the equalities $`JMJ=M^{}`$, $`\widehat{J}\widehat{M}\widehat{J}=\widehat{M}^{}`$, $`\widehat{J}M\widehat{J}=M`$ and $`J\widehat{M}J=\widehat{M}`$. The first two of these equalities follow from modular theory and the last two from proposition 2.1. ∎
## 3. A stronger form of left invariance.
In this section we want to prove some stronger form of left invariance of the Haar weight $`\phi `$. We want to show that $`(\iota \iota \phi )(\iota \mathrm{\Delta })(X)=(\iota \phi )(X)1`$ for any positive element $`XNM`$ and any von Neumann algebra $`N`$. The same formula is stated in for Kac algebras, but not proved. The first proof for this formula in the Kac algebra case was given bij Zsidó in (see also remark 18.23 in ). Unfortunately the proof of Zsidó does not work in the case of an arbitrary von Neumann algebraic quantum group, where possibly $`\tau _t\iota `$.
In our definition of a von Neumann algebraic quantum group we assumed the existence of invariant weights. The notion of left invariance we use, is in fact the weakest form of left invariance that one can assume, namely $`\phi \left((\omega \iota )\mathrm{\Delta }(x)\right)=\phi (x)\omega (1)`$ for all $`\omega M_{}^+`$ and $`x_\phi ^+`$. As a special case of the next proposition we will get the strongest form of left invariance, namely $`(\iota \phi )\mathrm{\Delta }(x)=\phi (x)1`$ for all $`xM^+`$. Some result in between was already proved in proposition 5.15 of , and will be used in the proof of the proposition.
When $`N`$ is a von Neumann algebra we denote by $`N^+\text{ext}`$ the extended positive part of $`N`$ as was already mentioned in the introduction. In the proof of the next proposition we denote by $`,`$ the composition of elements in $`N^+\text{ext}`$ and $`N_{}^+`$.
###### Proposition 3.1.
Let $`N`$ be a von Neumann algebra and $`X(NM)^+`$. Then we have
$$(\iota \iota \phi )(\iota \mathrm{\Delta })(X)=(\iota \phi )(X)1.$$
Here both sides of the equation make sense in $`(NM)^+\text{ext}`$. In particular we get
$$(\iota \phi )\mathrm{\Delta }(x)=\phi (x)1$$
for all $`xM^+`$.
###### Proof.
We will prove the proposition for the dual von Neumann algebraic quantum group $`(\widehat{M},\widehat{\mathrm{\Delta }})`$. Because of the biduality theorem this proves the stated result. We also represent $`N`$ on a Hilbert space $`K`$ and then it is enough to prove the proposition in case $`N=B(K)`$.
Recall that we introduced the multiplicative unitary $`\widehat{V}`$ in definition 2.14. Then define for every $`zB(KH)^+`$ the element $`T(z)B(KH)^+\text{ext}`$ by the following formula, which makes sense because $`\widehat{V}B(H)\widehat{M}`$.
$$T(z)=(\iota \iota \widehat{\phi })\left((1\widehat{V})(z1)(1\widehat{V}^{})\right).$$
When $`\eta KH`$ we denote by $`P_\eta `$ the positive rank one operator defined by $`P_\eta (\xi )=\xi ,\eta \eta `$. Let now $`\eta KH`$ and suppose $`\eta =1`$. Choose an orthonormal basis $`(e_i)_{iI}`$ of $`KH`$ such that $`\eta =e_i`$ for some $`iI`$.
Choose $`\xi KH`$. Then we have
$`T(P_\eta ),\omega _\xi `$ $`=\widehat{\phi }\left((\omega _\xi \iota )\left((1\widehat{V})(P_\eta 1)(1\widehat{V}^{})\right)\right)`$
$`={\displaystyle \underset{iI}{}}\widehat{\phi }\left(\left((\omega _{\xi ,e_i}\iota )((P_\eta 1)(1\widehat{V}^{}))\right)^{}(\omega _{\xi ,e_i}\iota )((P_\eta 1)(1\widehat{V}^{}))\right)`$
$`=\widehat{\phi }\left(\left((\omega _{\xi ,\eta }\iota )(1\widehat{V}^{})\right)^{}(\omega _{\xi ,\eta }\iota )(1\widehat{V}^{})\right).`$
In proposition 2.15 we saw that $`\widehat{V}^{}=(JJ)W^{}(JJ)=(w^{}1)W(w1)`$ where $`w=\widehat{J}J`$. The last equality follows from corollary 2.2. So it follows that
$$T(P_\eta ),\omega _\xi =\widehat{\phi }\left(\left((\omega _{(1w)\xi ,(1w)\eta }\iota )(1W)\right)^{}(\omega _{(1w)\xi ,(1w)\eta }\iota )(1W)\right).$$
In remark 8.31 of we saw that for $`\omega M_{}`$ one has $`(\omega \iota )(W)𝒩_{\widehat{\phi }}`$ if and only if $`\omega `$. So it follows that $`T(P_\eta ),\omega _\xi <\mathrm{}`$ if and only if
$$\omega _{(1w)\xi ,(1w)\eta }(1)$$
(3.1)
and in that case
$$T(P_\eta ),\omega _\xi =\xi (\omega _{(1w)\xi ,(1w)\eta }(1))^2.$$
Suppose that $`uM`$ is a unitary. And suppose that formula (3.1) is valid. We claim that
$$\omega _{(1w)(1JuJ)\xi ,(1w)\eta }(1)$$
and
$$\xi \left(\omega _{(1w)(1JuJ)\xi ,(1w)\eta }(1)\right)=R(u^{})\xi \left(\omega _{(1w)\xi ,(1w)\eta }(1)\right).$$
For this choose $`x𝒩_\phi `$ and make the following computation:
$`\omega _{(1w)(1JuJ)\xi ,(1w)\eta }`$ $`(1x^{})`$
$`=(1x^{})(1\widehat{J}uJ)\xi ,(1w)\eta `$
$`=(1x^{})(1R(u^{}))(1w)\xi ,(1w)\eta `$
$`=\xi \left(\omega _{(1w)\xi ,(1w)\eta }(1)\right),R(u)\mathrm{\Lambda }(x)`$
$`=R(u^{})\xi \left(\omega _{(1w)\xi ,(1w)\eta }(1)\right),\mathrm{\Lambda }(x)`$
From this follows our claim.
But then we get for every $`\xi KH`$ and every unitary $`uM`$ that
$$T(P_\eta ),\omega _\xi =T(P_\eta ),\omega _{(1JuJ)\xi }.$$
From this we may conclude that $`T(P_\eta )(B(K)M)^+\text{ext}`$, for all $`\eta KH`$. Let now $`zB(KH)^+`$. Let $`(e_i)_{iI}`$ again be an orthonormal basis for $`KH`$. Then
$$z=\underset{iI}{}z^{1/2}P_{e_i}z^{1/2}=\underset{iI}{}P_{z^{1/2}e_i}.$$
By lower semicontinuity of $`T`$ we can conclude that $`T(z)(B(K)M)^+\text{ext}`$.
Let now $`X\left(B(K)\widehat{M}\right)^+`$. Then
$$T(X)=(\iota \iota \widehat{\phi })(\iota \widehat{\mathrm{\Delta }})(X)$$
and this clearly belongs to $`(B(K)\widehat{M})^+\text{ext}`$. But it also belongs to $`(B(K)M)^+\text{ext}`$ by the result in the previous paragraph. Let
$$T(X)=\mathrm{}(1e)+_0^{\mathrm{}}\lambda 𝑑e_\lambda $$
be the unique spectral decomposition of $`T(X)`$, considered as an element of $`B(KH)^+\text{ext}`$. Then
$$e,e_\lambda \left(B(K)M\right)\left(B(K)\widehat{M}\right)=B(K)$$
because $`M\widehat{M}=`$. So take $`f,f_\lambda B(K)`$ such that $`e=f1`$ and $`e_\lambda =f_\lambda 1`$. Then define the element $`SB(K)^+\text{ext}`$ by
$$S=\mathrm{}(1f)+_0^{\mathrm{}}\lambda 𝑑f_\lambda .$$
Then we get that
$$(\iota \iota \widehat{\phi })(\iota \widehat{\mathrm{\Delta }})(X)=S1.$$
(3.2)
Let us now suppose first that $`K=`$. This will prove the special case stated in the proposition. Then $`X\widehat{M}^+`$ and $`S`$ will be a scalar. So we get a $`\lambda [0,+\mathrm{}]`$ such that
$$(\iota \widehat{\phi })\widehat{\mathrm{\Delta }}(X)=\lambda \mathrm{\hspace{0.17em}1}.$$
Now there are two possibilities.
* Either there exists a $`\omega \widehat{M}_{}^+`$ with $`\omega 0`$ such that $`(\omega \iota )\widehat{\mathrm{\Delta }}(X)_{\widehat{\phi }}^+`$. Then $`\lambda <+\mathrm{}`$ because
$$\lambda \omega (1)=\widehat{\phi }\left((\omega \iota )\widehat{\mathrm{\Delta }}(X)\right)<\mathrm{}.$$
But then also
$$\widehat{\phi }\left((\mu \iota )\widehat{\mathrm{\Delta }}(X)\right)=\lambda \mu (1)<\mathrm{}$$
for all $`\mu \widehat{M}_{}^+`$, and so $`(\mu \iota )\widehat{\mathrm{\Delta }}(X)_{\widehat{\phi }}^+`$ for all $`\mu \widehat{M}_{}^+`$. Then it follows from proposition 5.15 in that $`X_{\widehat{\phi }}^+`$ and so $`\lambda =\widehat{\phi }(X)`$ because of left invariance.
* Either we have $`\widehat{\phi }\left((\omega \iota )\widehat{\mathrm{\Delta }}(X)\right)=+\mathrm{}`$ for all $`\omega \widehat{M}_{}^+\{0\}`$. This means that $`\lambda =+\mathrm{}`$. Because of left invariance we cannot have $`X_{\widehat{\phi }}^+`$ and so $`\widehat{\phi }(X)=+\mathrm{}`$. Again $`\lambda =\widehat{\phi }(X)`$.
In both cases we arrive at $`(\iota \widehat{\phi })\widehat{\mathrm{\Delta }}(X)=\widehat{\phi }(X)\mathrm{\hspace{0.17em}1}`$.
Now we return to the general case. Let $`\omega B(K)_{}^+`$ and $`\mu \widehat{M}_{}^+`$. Then we apply $`\omega \mu `$ to equation (3.2). This gives us
$`S,\omega \mu (1)=\widehat{\phi }\left((\omega \mu \iota )(\iota \widehat{\mathrm{\Delta }})(X)\right)`$ $`=\widehat{\phi }\left((\mu \iota )\widehat{\mathrm{\Delta }}((\omega \iota )(X))\right)`$
$`=\mu (1)\widehat{\phi }\left((\omega \iota )(X)\right)=\mu (1)(\iota \widehat{\phi })(X),\omega .`$
In this computation we used the special case of the proposition proved above. So it follows that $`S=(\iota \widehat{\phi })(X)`$ and this gives what we wanted to prove. ∎
## 4. The opposite and the commutant von Neumann algebraic quantum group
Given a von Neumann algebraic quantum group $`(M,\mathrm{\Delta })`$, represented standardly such that $`(H,\iota ,\mathrm{\Lambda })`$ is a GNS-construction for the left Haar weight $`\phi `$, we can define two new von Neumann algebraic quantum groups called the opposite von Neumann algebraic quantum group $`(M,\mathrm{\Delta })\text{op}`$ and the commutant von Neumann algebraic quantum group $`(M,\mathrm{\Delta })^{}`$. With the notations introduced before we give the following definition.
###### Definition 4.1.
The underlying von Neumann algebra of the opposite von Neumann algebraic quantum group $`(M,\mathrm{\Delta })\text{op}`$ is again $`M`$ and the comultiplication $`\mathrm{\Delta }\text{op}`$ is given by $`\mathrm{\Delta }\text{op}(x)=\chi \mathrm{\Delta }(x)`$ for all $`xM`$.
The underlying von Neumann algebra of the commutant von Neumann algebraic quantum group $`(M,\mathrm{\Delta })^{}`$ is given by $`M^{}`$ and the comultiplication $`\mathrm{\Delta }^{}`$ is defined by $`\mathrm{\Delta }^{}(x)=(JJ)\mathrm{\Delta }(JxJ)(JJ)`$ for all $`xM^{}`$.
It is easy to see that $`(M,\mathrm{\Delta })\text{op}`$ and $`(M,\mathrm{\Delta })^{}`$ are again von Neumann algebraic quantum groups. We will now give canonical choices for the left invariant weights and their GNS-construction. As a left invariant weight on $`(M,\mathrm{\Delta })\text{op}`$ we take $`\psi `$, with GNS-construction $`(H,\iota ,\mathrm{\Gamma })`$. On $`M^{}`$ we define the weight $`\phi ^{}`$ by $`\phi ^{}(x)=\phi (JxJ)`$ for all $`x(M^{})^+`$. Then $`\phi ^{}`$ is a left invariant weight on $`(M,\mathrm{\Delta })^{}`$, with GNS-construction $`(H,\iota ,\mathrm{\Lambda }^{})`$, where $`\mathrm{\Lambda }^{}(x)=J\mathrm{\Lambda }(JxJ)`$ for all $`x𝒩_\phi ^{}`$.
Given these GNS-constructions we can define the multiplicative unitaries $`W\text{op}`$ and $`W^{}`$ associated to $`(M,\mathrm{\Delta })\text{op}`$ and $`(M,\mathrm{\Delta })^{}`$ and it is clear that, using definition 2.14 and proposition 2.15, they are given by
$$W\text{op}=\mathrm{\Sigma }V^{}\mathrm{\Sigma }\text{and}W^{}=(JJ)W(JJ)=\widehat{V}.$$
It is also clear that the unitary antipode $`R\text{op}`$ of $`(M,\mathrm{\Delta })\text{op}`$ equals $`R`$ and the unitary antipode $`R^{}`$ of $`(M,\mathrm{\Delta })^{}`$ is given by $`R^{}(x)=JR(JxJ)J`$ for all $`xM^{}`$. So the canonical right invariant weights on $`(M,\mathrm{\Delta })\text{op}`$ and $`(M,\mathrm{\Delta })^{}`$ are $`\phi `$ and $`\psi ^{}`$. Then the modular elements $`\delta \text{op}`$ and $`\delta ^{}`$ are given by
$$\delta \text{op}=\delta ^1\text{and}\delta ^{}=J\delta J.$$
One also checks easily that $`\tau _t\text{op}`$ equals $`\tau _t`$ and $`\tau _t^{}(x)=J\tau _t(JxJ)J`$ for all $`t`$ and $`xM^{}`$.
Defining the unitary $`w=\widehat{J}J=\nu ^{i/4}J\widehat{J}`$ it is easy to see that $`\mathrm{\Phi }:MM^{}:\mathrm{\Phi }(x)=wxw^{}`$ gives an isomorphism between the von Neumann algebraic quantum groups $`(M,\mathrm{\Delta })`$ and $`(M,\mathrm{\Delta })^{}\text{op}`$. To prove this we only have to observe that $`R(x)=\widehat{J}x^{}\widehat{J}`$ for all $`xM`$ and $`(RR)\mathrm{\Delta }(x)=\mathrm{\Delta }\text{op}(R(x))`$ for all $`xM`$.
We conclude this section with the following formulas.
###### Proposition 4.2.
With the notations introduced above we have:
$`(M,\mathrm{\Delta })\text{op}\widehat{\text{}}`$ $`=(M,\mathrm{\Delta })\widehat{\text{}}^{}`$
$`(M,\mathrm{\Delta })^{}\widehat{\text{}}`$ $`=(M,\mathrm{\Delta })\widehat{\text{}}\text{op}`$
$`(M,\mathrm{\Delta })^{}\text{op}`$ $`=(M,\mathrm{\Delta })\text{op}^{}.`$
###### Proof.
Because $`W\text{op}=\mathrm{\Sigma }V^{}\mathrm{\Sigma }`$, the von Neumann algebra underlying $`(M,\mathrm{\Delta })\text{op}\widehat{\text{}}`$ is given by
$$\{(\omega \iota )(W\text{op})\omega B(H)_{}\}^{\prime \prime }=\{(\iota \omega )(V^{})\omega B(H)_{}\}^{\prime \prime }=\widehat{M}^{}.$$
The last equality follows from proposition 2.15. Further we have for every $`x\widehat{M}^{}`$ that
$$\mathrm{\Delta }\text{op}\widehat{\text{}}(x)=\mathrm{\Sigma }W\text{op}(x1)(W\text{op})^{}\mathrm{\Sigma }=V^{}(1x)V.$$
Because $`V=(\widehat{J}\widehat{J})\mathrm{\Sigma }W^{}\mathrm{\Sigma }(\widehat{J}\widehat{J})`$, this gives
$$\mathrm{\Delta }\text{op}\widehat{\text{}}(x)=\mathrm{\Sigma }(\widehat{J}\widehat{J})W(\widehat{J}x\widehat{J}1)W^{}(\widehat{J}\widehat{J})\mathrm{\Sigma }=(\widehat{J}\widehat{J})\widehat{\mathrm{\Delta }}(\widehat{J}x\widehat{J})(\widehat{J}\widehat{J})=\widehat{\mathrm{\Delta }}^{}(x).$$
This gives $`(M,\mathrm{\Delta })\text{op}\widehat{\text{}}=(M,\mathrm{\Delta })\widehat{\text{}}^{}`$.
Applying this last formula to $`(M,\mathrm{\Delta })\widehat{\text{}}`$ and using the biduality theorem we get $`(M,\mathrm{\Delta })\widehat{\text{}}\text{op}\widehat{\text{}}=(M,\mathrm{\Delta })^{}`$. Taking the dual and using once again the biduality theorem this gives our second result $`(M,\mathrm{\Delta })\widehat{\text{}}\text{op}=(M,\mathrm{\Delta })^{}\widehat{\text{}}`$.
To compute $`(M,\mathrm{\Delta })\text{op}^{}`$ we have to observe once again that the modular conjugation $`J^{}`$ of the left invariant weight $`\psi `$ on $`(M,\mathrm{\Delta })\text{op}`$ is given by $`J^{}=\nu ^{i/4}J`$. Then it is clear that $`(M,\mathrm{\Delta })\text{op}^{}=(M,\mathrm{\Delta })^{}\text{op}`$. ∎
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# Does Special Relativity Have Limits of Applicability in the Domain of Very Large Energies?
## I Introduction
The special relativity theory (SR) is one of the physics theories that are compose the base of the modern physics, it has well experimental foundation in the large area of the reached velocities and energies, is the working theory of a modern physics and widely use in a science and technique. Nevertheless, as any physical theory created by men, SR has a boundaries of applicability (the inertial systems of reference). As far as we know a problem of energy boundaries SR for any moving body was not analyzed in detail. At the same time the tends of energy of a moving body when it’s velocity reaches the value of speed of light to infinity calls doubts in applicability SR in this area of energies, as the occurrence of infinity in the physical theories always testifies about their interior deficiencies. The purpose of the paper is presenting the theory of almost inertial systems in the space with multifractal time (I believe that time and space of our world may have multifractal characteristics) in which the motions of bodies with arbitrary velocities are possible and there are no energy infinity. At the same time if the energies reached by a body, as a result of it’s motion, are smaller than $`E_010^3sec^{1/2}`$ $`t^{1/2}`$ ($`t`$-time of acceleration of a particle up to such energies, $`E_0`$\- a rest energy) all results of the theory coincide with results of SR and thus do not contradict experiment. The differences between our theory and SR appear only if energy of moving body exceeds the energy $`E_0`$$`10^3`$$`sec^{1/2}t^{1/2}`$. The theory is based on using, but as a good approach, of a principle of the constancy speed of light, an invariance of modified Galilean and Lorentz transformation laws. This theory is not a generalization of SR, because any SR generalizations in the domain of validity SR (inertial systems) are absurd. This theory describes relative movements only in the almost inertial systems, and thus does not contradict SR, thou coincide with it for case of inertial systems. For constructing such theory it is necessary to refuse from the rigorous realization of the SR postulates: a homogeneous of space and time, the constancy of speed of light, the Galilean relativity principle. In an inhomogeneous space and time, if the inhomogeneous are small, the motion of bodies will be almost inertial, and velocity of light is almost stationary value. This paper contains the example of the theory based on the time and the space with fractional dimensions (FD) . This theory use for description of time and space characteristics the ideas of fractal geometry. The values FD are a little bit distinguished from the integer dimensions (the theory of multifractal time and space is given in . In fractal theory the motions of particles with arbitrary velocities are permissible, the speed of light is almost independent from the velocity of lights sources. For example, on the surface of Earth the differences of value of speed of light under change moving direction $`v`$ by $`v`$ consist $`2v/c10^6t`$. The theory almost coincides with SR, for velocities which are lesser than speed of light but does not includes singularities at $`v=c`$.
## II Multifractal time
Following , we will consider both time and space as the initial real material fields existing in the world and generating all other physical fields. Assume that every of them consists of a continuous, but not differentiable bounded set of small elements (elementary intervals, further treated as ”points”). Consider the set of small time elements $`S_t`$. Let time be defined on multifractal subsets of such elements, defined on certain measure carrier $`^n`$. Each element of these subsets (or ”points”) is characterized by the fractional (fractal) dimension (FD) $`d_t(𝐫(t),t)`$ and for different elements FD are different. In this case the classical mathematical calculus or fractional (say, Riemann - Liouville) calculus can not be applied to describe a small changes of a continuous function of physical values $`f(t)`$, defined on time subsets $`S_t`$, because the fractional exponent depends on the coordinates and time. Therefore, we have to introduce integral functionals (both left-sided and right-sided) which are suitable to describe the dynamics of functions defined on multifractal sets (see ). Actually, these functionals are simple and natural generalization of the Riemann-Liouville fractional derivatives and integrals:
$$D_{+,t}^df(t)=\left(\frac{d}{dt}\right)^n_a^t\frac{f(t^{})dt^{}}{\mathrm{\Gamma }(nd(t^{}))(tt^{})^{d(t^{})n+1}}$$
(1)
$$D_{,t}^df(t)=(1)^n\left(\frac{d}{dt}\right)^n_t^b\frac{f(t^{})dt^{}}{\mathrm{\Gamma }(nd(t^{}))(t^{}t)^{d(t^{})n+1}}$$
(2)
where $`\mathrm{\Gamma }(x)`$ is Euler’s gamma function, and $`a`$ and $`b`$ are some constants from $`[0,\mathrm{})`$. In these definitions, as usually, $`n=\{d\}+1`$ , where $`\{d\}`$ is the integer part of $`d`$ if $`d0`$ (i.e. $`n1d<n`$) and $`n=0`$ for $`d<0`$. If $`d=const`$, the generalized fractional derivatives (GFD) (1)-(2) coincide with the Riemann - Liouville fractional derivatives ($`d0`$) or fractional integrals ($`d<0`$). When $`d=n+\epsilon (t),\epsilon (t)0`$, GFD can be represented by means of integer derivatives and integrals. For $`n=1`$, that is, $`d=1+\epsilon `$, $`\left|\epsilon \right|<<1`$ it is possible to obtain:
$$D_{+,t}^{1+\epsilon }f(t)\frac{}{t}f(t)+a\frac{}{t}\left[\epsilon (r(t),t)f(t)\right]$$
(3)
where $`a`$ is constant and defined by the choice of the rules of regularization of integrals (1)-(2) (for more detailed see ). The selection of the rule of regularization that gives a real additives for usual derivative in (3) yield $`a=0.5`$ for $`d<1`$ and $`a=1.077`$ for $`d>1`$ . The functions under integral sign in (1)-(2) we consider as the generalized functions defined on the set of the finite functions . The notions of GFD, similar to (1)-(2), can also be defined and for the space variables $`𝐫`$.
The definitions of GFD (2)-(2) are formal until the connections between fractal dimensions of time $`d_t(𝐫(t),t)`$ and certain characteristics of physical fields (say, potentials $`\mathrm{\Phi }_i(𝐫(t),t),i=1,2,..)`$ or densities of Lagrangians $`L_i`$) are determined. Following , we define this connection by the relation
$$d_t(𝐫(t),t)=1+\underset{i}{}\beta _iL_i(\mathrm{\Phi }_i(𝐫(t),t))$$
(4)
where $`L_i`$ are densities of energy of physical fields, $`\beta _i`$ are dimensional constants with physical dimension of $`[L_i]^1`$ (it is worth to choose $`\beta _i^{}`$ in the form $`\beta _i^{}=a^1\beta _i`$ for the sake of independence from regularization constant). The definition of time as the system of subsets and definition the FD $`d`$ (see 4) connects the value of fractional (fractal) dimension $`d_t(r(t),t)`$ with each time instant $`t`$. The latter depends both on time $`t`$ and coordinates $`𝐫`$. If $`d_t=1`$ (an absence of physical fields) the set of time has topological dimension equal to unity. The multifractal model of time allows, as will be shown below, to consider the divergence of energy of masses moving with speed of light in the SR theory as the result of the requirement of rigorous validity of the laws pointed out in the beginning of this paper in the presence of physical fields (in our theory there are approximate fulfillment of these laws).
## III The principle of the speed of light invariance
Because of the non-uniformity of time in our multifractal model, the speed of light, just as in the general relativity theory, depends on potentials of physical fields that define the fractal dimensions of time $`d_t(𝐫(t),t)`$ (see (4)). If fractal dimension $`d_t(𝐫(t),t)`$ is close enough to unity ($`d_t(r(t),t)=1+\epsilon ,\left|\epsilon \right|<<1`$), the difference of the speed of light in moving (with velocity $`v`$ along the $`x`$ axis) and fixed frame of reference will be small. In the systems that move with respect to each other with almost constant velocity (stationary velocities do not exist in the mathematical theory based on the definitions of GFD (1) - (2)) the speed of light can not be taken as a fundamental constant. In the multifractal time theory the principle of the speed of light invariance can be considered only as approximate. But if $`\epsilon `$ is small, it allows to consider a nonlinear coordinates transformations from the fixed frame to the moving frame (replacing the Galilean transformations in non-uniform time and space), as close to linear (weakly nonlinear) transformations and, thus, makes it possible to preserve the conservation laws, and all the invariant’s of the Minkowski space, as the approximate laws. Then the way of reasoning and argumentation accepted in SR theory (see for example, ()) can also remains valid. Designating the coordinates in the moving and fixed frames of reference through $`x^{}`$ and $`x`$, accordingly, we write down
$`x^{}`$ $`=`$ $`\alpha (t,x)[xv(x,t)t(x(t),t]`$ (5)
$`x`$ $`=`$ $`\alpha ^{}(t,x)[x^{}+v^{}(x^{}(t^{}),t^{}),t^{}(x^{}(t^{}),t^{})`$ (6)
In (5) $`\alpha \alpha ^{}`$ and the velocities $`v^{}`$ and $`v`$ (as well as $`t`$ and $`t^{}`$) are not equal (it follows from the inhomogeneous of multifractal time). Place clocks in origins of both the frame of references and let the light signal be emitted in the moment, when the origins of the fixed and moving frames coincide in space and time at the instant $`t{}_{}{}^{1}=t=0`$ and in points $`x^{}=x=0`$. The propagation of light in moving and fixed frames of reference is then determined by equations
$$x^{}=c^{}t^{}x=ct$$
(7)
These equations characterize the propagation of light in both of the frames of reference at every moment. Due to the time inhomogeneous $`c^{}c`$, but since $`|\epsilon <<1|`$ the difference between velocities of light in the two frames of reference will be small. For this case we can neglect by the differences between $`\alpha ^{}`$ and $`\alpha `$ and, for different frames of reference write the expressions for velocities of light (using (3) to define velocity (denote $`f(t)=x,dx/dt=c_0`$)). Thus we obtain
$$c=D_{+,t}^{1+\epsilon }x=c_0(1\epsilon )\frac{d\epsilon }{dt}x$$
(8)
$$c^{}=D_{+,t}^{1+\epsilon ^{}}x^{}=c_0(1\epsilon )\frac{d\epsilon }{dt}x^{}$$
(9)
$$c_1=c_0(1\epsilon )\frac{d\epsilon }{dt}x^{}$$
(10)
$$c_1^{}=c_0(1\epsilon )\frac{d\epsilon }{dt}x$$
(11)
The equalities (10) and (11) appear in our model of multifractal time as the result of the statement, that in this model all the frames of reference are absolute frames of reference (because of the real nature of the time field) and the speed of light depends on the state of frames: if the frame of reference is a moving or a fixed one, if the object under consideration in this frame moves or not. This dependence disappears only when $`\epsilon =0`$. Before substitution the relations (5) in the equalities (8) - (11) (with $`\alpha ^{}\alpha `$) it is necessary to find out how $`d\epsilon /dt`$ depends on $`\alpha `$. Using for this purpose equation (4) we obtain:
$$\frac{d\epsilon }{dt}=\frac{d\epsilon }{d𝐫}𝐯\underset{i}{}\beta _i(𝐅_i𝐯+\frac{L_i}{t})$$
(12)
where $`𝐅_i=\frac{dL_i}{d𝐫}`$. As the forces for moving frames of reference are proportional to $`\alpha `$ we get (for the case when there is no explicit dependence of $`L_i`$ on time)
$$\frac{d\epsilon }{dt}\underset{i}{}\beta _i𝐅_{0i}𝐯\alpha $$
(13)
where $`F_{0i}`$ are the corresponding forces at zero velocity. Multiplying (8) - (11) on the corresponding times $`t,t^{},t_1,t_1^{}`$ yields the following expressions
$$c^{}t^{}=c_0t\left[1+\frac{v_i\beta _iF_{0i}}{c_0}\alpha ^2ct(1\frac{v}{c})\right]$$
(14)
$$ct=c_0t^{}\left[1+\frac{v_i\beta _iF_{0i}}{c_0}\alpha ^2ct(1\frac{v}{c})\right]$$
(15)
$$c_1^{}t_1^{}=c_0t_1\left[1\frac{v_i\beta _iF_{0i}}{c_0}\alpha ^2ct(1\frac{v}{c})\right]$$
(16)
$$c_1t_1=c_0t_1^{}\left[1\frac{v_i\beta _iF_{0i}}{c_0}\alpha ^2ct(1\frac{v}{c})\right]$$
(17)
Since in our model the motion and frames of reference are absolute, the times $`t_1`$ and $`t_1^{}`$ correspond to the cases, when the moving and fixed frames of reference exchange their roles - the moving one becomes fixed and vice versa. These times coincide only when $`\epsilon =0`$. The times in square brackets, as well as the velocities, are taken to equal, because the terms containing them are small as compared to unity. The principle of invariance of the velocity of light for transition between the moving and fixed frames of reference in multifractal time model is approximate (though quite natural, because the frames of reference are absolute frames of reference). Taking into account (5), the relations (14) - (17) take the forms
$$c^{}t^{}=c\alpha t(1\frac{v}{c}),c_1^{}t_1^{}=c\alpha t_1(1\frac{v}{c})$$
(18)
$$ct=c\alpha t^{}(1+\frac{v}{c}),c_1t_1=c\alpha t^{}(1+\frac{v}{c})$$
(19)
Once again we note, that the four equations for $`c_1^{}t_1^{}`$ and $`c_1t_1`$, instead of the two equations in special relativity, appear as the consequences of the absolute character of the motion and frames of reference in the model of multifractal time. In the right-hand side of (18) - (19) the dependence of velocity of light on fractal dimensions of time is not taken into account (just as in the equations (14) - (17)). Actually, this dependence leads to pretty unwieldy expressions. But if we retain only the terms that depend on $`\beta =\sqrt{|1v^2/c^2|}`$ or $`a_0`$ and neglect non-essential terms containing the products $`\beta \alpha _0`$, utilizing (14) - (17) after the multiplication of the four equalities (18) - (19), we receive the following equation for $`\alpha `$ (it satisfies to all four equations):
$$4a_0^4\beta ^4\alpha ^84a_0^2\alpha ^4+1=\beta ^4\alpha ^4+4a_0^4\beta ^4\alpha ^8$$
(20)
where
$$\beta =\sqrt{\left|1\frac{v^2}{c^2}\right|}$$
(21)
$$a_0=\underset{i}{}\beta _iF_{0i}\frac{v}{c}ct$$
(22)
From (20) follows
$$\alpha _1\beta _{}^1=\frac{1}{\sqrt[4]{\beta ^4+4a_0^2}}$$
(23)
The solutions $`\alpha _{2,3,4}`$ are given by $`\alpha _2=\alpha _1,\alpha _{3,4}=\pm i\alpha `$. Applicability of above obtained results is restricted by requirement $`|\epsilon |1`$
## IV Lorentz transformations and transformations of length <br>and time in multifractal time model
The Lorentz transformations, as well as transformations of coordinate frames of reference, in the multifractal model of time are nonlinear due to the dependence of the fractional dimensions of time $`d_t(𝐫,t)`$ on coordinates and time. Since the nonlinear corrections to Lorentz transformation rules are very small for $`\epsilon 1`$, we shall take into account only the corrections that eliminate the singularity at the velocity $`v=c`$. It yields in the replacement of the factor $`\beta ^1`$ in Lorentz transformations by the modified factor $`\alpha =1/\beta ^{}`$ given by (23). The Lorentz transformation rules (for the motion along the $`x`$ axis) take the form
$$x^{}=\frac{1}{\beta ^{}}(xvt),t^{}\frac{1}{\beta ^{}}(tx\frac{v}{c^2})$$
(24)
In the equations (23) and (24) the velocities $`v`$ and $`c`$ weakly depend on $`x`$ and $`t`$ and their contribution to the singular terms are small. Hence, we can neglect by this dependence. The transformations from fixed system to moving system are almost orthogonal (for $`\epsilon 1`$ ), and the squares of almost four-dimensional the energy-momentum vectors of Minkowski space vary under the coordinates transformations very slightly (i.e. they are almost invariant). Then it is possible to neglect the correction terms of order about $`O(\epsilon ,\dot{\epsilon })`$, which, for not equal to infinity variables, are very small too. From (23) - (24) the possibility of arbitrary velocity motion of bodies with nonzero rest mass follows. With the corrections of to order $`O(\epsilon ,\dot{\epsilon })`$ in nonsingular terms being neglected, the momentum and energy of a body with a nonzero rest mass in the frame of reference moving along the $`x`$ axis ($`E_0=m_0c^2)`$ equal to
$$p=\frac{1}{\beta ^{}}m_0v=\frac{m_0v}{\sqrt[4]{\beta ^4+4a_0^2}},E=E_0\sqrt{\frac{v^2c^2}{\sqrt{\beta ^4+4a_0^2}}+1}$$
(25)
The energy of such a body reaches its maximal value at $`v=c`$ and is equal then $`E_{v=c}E_0/\sqrt{2\alpha _0}`$. When $`v\mathrm{}`$ the energy is finite an tends to $`E_0\sqrt{2}`$. For $`vc`$ the total energy of a body is represented by the expression
$$E\frac{E_0}{\sqrt[4]{\beta ^4+4a_0^2}}=mc^2,m=\frac{m_0}{\beta ^{}}$$
(26)
For $`vc`$, total energy, defined by (25), is given by
$`m=\beta ^1m_0\sqrt{\text{1}+\beta ^2+\sqrt{\beta ^44a_0^2}}`$ (27)
If we take into account only the gravitational field of Earth (here, as in (),the gravitational field is a real field) and neglect by the influences of all the other fields), the parameter $`a_0(t)`$ can be estimated as $`a_0=r_0R^3x_Ect`$, where $`r_0`$ is the gravitational radius of Earth, $`r`$ is the distance from the Earth’s surface to its center ($`\epsilon =0.5\beta _g\mathrm{\Phi }_g,\beta _g=2c^2,x_Er_0,v=c`$). For energy maximum we get $`E_{max}E_010^3t^{0.5}sec^{0.5}`$. If we take into account only the constant electric field with electric strange $`E`$ then parameter $`a_0`$ has the value:$`a_0=\frac{eE}{Mc^2}ct`$ where $`Mc^2`$ is the rest energy of charges originated the strange $`E`$ and $`t`$ is the time of acceleration of particle. Contraction of lengths and the retardation of time in moving frames of reference in the model of multifractal time are also have several peculiarities. Let $`l`$ and $`t`$ be the length and time interval in a fixed frame of reference. In a moving frame
$$l^{}=\beta ^{}l,t^{}=\beta ^{}t$$
(28)
Thus, there exist the maximal contraction of length when the body’s velocity equals the speed of light, but length is not equal to zero. With the further increasing of velocity (if it is possible to fulfill some requirements for a motion in this region with constant velocity without radiating), the length of a body begins to grow and at infinitely large velocity is also infinite. The retardation of time, from the point of view of the observer in the fixed system (maximal retardation equals to $`t^{}=t\sqrt{2a_0}`$) is replaced, with the further increase of velocity over the speed of light, by acceleration of a flow of time ($`t0`$ when $`v\mathrm{}`$).
The rule for velocities transformation keeps its form, but $`\beta `$ is replaced by $`\beta ^{}`$
$$u_x=\frac{u_x^{}+v}{1+\frac{u_x^{}v}{c^2}},u_y=\frac{u_y^{}\beta ^{}}{1+\frac{u_y^{}v}{c^2}},u_z=\frac{u_z^{}\beta ^{}}{1+\frac{u_z^{}v}{c^2}}$$
(29)
Since there is no law that prohibits velocities greater than that of light, the velocities in (29) can also exceed the speed of light. The electrodynamics of moving media in the model of multifractal time can be obtained, in most cases, by the substitution $`\beta \beta ^{}`$.
## V Newton equation for relativistic particle
The Newton relativistic equation for particle with velocities $`vc`$ has form
$$\frac{}{t}(m𝐯)=\frac{}{t}\left(\frac{m_0𝐯}{\sqrt[4]{(1\frac{v_2}{c_2})^2+4a_0^2}}\right)=e𝐄$$
(30)
For $`vc`$ the mass $`m`$ in equation (30) is determined by (27). If $`𝐄=const.`$ (30) gives for $`v=c`$ possibility to find (neglecting by radiation of charge) minimum for the time $`t_0`$ that is necessary for receiving by particles the velocity equal $`c`$
$$\frac{m_0c}{\sqrt{\frac{2eEct_0}{Mc^2}}}=eEt_0$$
(31)
or
$$t_0=\sqrt[3]{\frac{Mc^2}{2E_0}}\frac{E_0}{eEc}$$
(32)
The $`t_0`$ defined by (32) gives only an order of the value $`t`$ that is necessary for receiving by particle velocity $`v=c`$ . The maximum energy at $`v=c`$ may be written now as (we introduce the value $`\alpha <1`$ for describing the radiation energy losses )
$$E_{max}=E_0\sqrt[3]{\frac{Mc^2}{2E_0\alpha }}$$
(33)
The values of $`\alpha `$ and electric strange $`E`$ are determined by construction of accelerators and conditions of their work regimes.
## VI Conclusions
The theory of relative motions in almost inertial systems based on the multifractal time theory and constructed in this paper gives in the new describing for characteristics (energy, momentum, mass and so on) of moving bodies. The main results are:a) the possibility of moving with arbitrary velocities without appearance of infinitum energy and imaginary mass; b) existence of maximum energy if $`v=c`$; c) possibility of experimental verification the main results of the theory. This theory describes open systems (theory of open systems see in ). The theory coincides with SR after transition to inertial systems (if neglect by the fractional dimensions of time) or almost coincides (the differences are non-essential) for velocities $`v<c`$. The movement of bodies with velocities that exceed the speed of light is accompanied by a series of physical effect’s which can be found experimentally (these effects was considered in the separate papers () in more details). It is necessary for verification of the theory to receive the particle’s with energy $`E\frac{E_0}{\sqrt{\frac{2eE}{Mc^2}ct_0}}`$ where $`t_0`$ defined by (32). If take into account the radiation losses of energy, it will be enough to receive at the intervals of time acceleration of the particle about one second the energies of order $`EE_010^3`$.
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# Global well-posedness below energy space for the 1D Zakharov system
## 0 Introduction
Consider the Cauchy problem for the (1+1)-dimensional Zakharov system
$`iu_t+u_{xx}`$ $`=`$ $`nu`$ (1)
$`n_{tt}n_{xx}`$ $`=`$ $`(|u|^2)_{xx}`$ (2)
$`u(0)=u_0,n(0)`$ $`=`$ $`n_0,n_t(0)=n_1`$ (3)
where $`u`$ is a complex-valued and $`n`$ a real-valued function defined for $`(x,t)𝐑\times 𝐑^+`$.
The main result shows global well-posedness of the problem for rough data
$$(u_0,n_0,n_1)H^{s,2}(𝐑)\times L^2(𝐑)\times \dot{H}^{1,2}(𝐑)\text{with}1>s>9/10$$
without any smallness assumption. The same result for $`s=1`$ is a direct consequence of the local well-posedness shown by and the conservation laws satisfied for solutions of (1),(2),(3), namely conservation of $`u(t)`$ and
$$E(u,n):=u_x(t)^2+1/2(n(t)^2+V(t)^2)+_{\mathrm{}}^+\mathrm{}n(t)|u(t)|^2𝑑x$$
where $`V_x=n_t`$.
Local well-posedness for $`s>9/10`$ follows from so that the problem is to show that the local solution exists globally in time.
Local and global well-posedness in dimension 2+1 and 3+1 for finite energy solutions was shown in .
Our proof uses Bougain’s ideas who introduces a general method to show global well-posedness for some types of nonlinear evolution equations for data with less regularity than needed for an application of the conservation laws directly. He applied it to the (2+1)- and (3+1)-dimensional Schrödinger equation ,. Later it was also used for other model equations ,,,.
This paper is organized as follows. In section 1 the needed estimates for the nonlinearities in the $`X^{s,b}`$–spaces introduced by Bourgain and the $`Y^s`$–spaces introduced by Ginibre, Tsutsumi and Velo are given along the lines of .
For an equation of the form
$$iu_t+\varphi (i_x)u=0$$
(4)
where $`\varphi `$ is a measurable real-valued function, let $`X^{s,b}`$ be the completion of $`S(𝐑^2)`$ with respect to
$`f_{X^{s,b}}`$ $`:=`$ $`e^{it\varphi (i_x)}f_{H_t^b(𝐑,H_x^s(𝐑))}`$
$`=`$ $`<\xi >^s<\tau >^b(e^{it\varphi (i_x)}f(x,t))_{L_{x,t}^2}`$
$`=`$ $`<\xi >^s<\tau +\varphi (\xi )>^b\widehat{f}(\xi ,\tau )_{L_{\xi ,\tau }^2}`$
$`\dot{X}^{s,b}`$ is defined similarly by replacing $`<\xi >^s`$ by $`|\xi |^s`$.
Similarly let $`Y^s`$ be the completion of $`S(𝐑^2)`$ with respect to
$`f_{Y^s}`$ $`:=`$ $`<\xi >^s<\tau >^1(e^{it\varphi (i_x)}f(x,t))_{L_\xi ^2L_\tau ^1}`$
$`=`$ $`<\xi >^s<\tau +\varphi (\xi )>^1\widehat{f}(\xi ,\tau )_{L_\xi ^2L_\tau ^1}`$
In our case we shall use these spaces for the phase functions $`\varphi (\xi )=\xi ^2`$ and $`\varphi (\xi )=\pm |\xi |`$.
We also have to use the norms in $`X^{s,b}(I)`$ for a given time interval I defined as
$$f_{X^{s,b}(I)}=\underset{\stackrel{~}{f}_{|I}=f}{inf}\stackrel{~}{f}_{X^{s,b}}\text{and similarly}f_{Y^s(I)}=\underset{\stackrel{~}{f}_{|I}=f}{inf}\stackrel{~}{f}_{Y^s}$$
In section 2 we transform the system in a standard way into a first order system for $`(u,n_+,n_{})`$. For details we again refer to . Using Bourgain’s ideas we split the datum $`u_0`$ into a sum $`u_{01}+u_{02}`$ where the low frequency part $`u_{01}`$ is regular and has large $`H^1`$ \- norm whereas the high frequency part $`u_{02}`$ is just in $`H^s`$ with small $`L^2`$ \- norm.
In section 3 the solution $`(\stackrel{~}{u},\stackrel{~}{n_+},\stackrel{~}{n_{}})`$ of the problem with data $`(u_{01},n_{0+},n_0)`$ is further investigated on a suitable time interval I depending on s using the energy bounds.
In section 4 we consider the system fulfilled by $`(v,m_\pm )=(u\stackrel{~}{u},n_\pm \stackrel{~}{n_\pm })`$ with data $`(u_{02},0,0)`$ and construct a solution in the same time interval I, thus we have a solution of the original problem on I. The inhomogeneous part $`w(t)`$ of $`v(t)`$ is shown to belong to $`H^{1,2}(𝐑)`$, thus is smoother than the homogeneous part $`e^{it_x^2}u_{02}`$ which is just in $`H^{s,2}(𝐑)`$.
In section 5 we show that this process can be iterated to construct a solution on any time interval $`[0,T]`$. What one does is to construct a solution on time intervals of equal length $`|I|`$. One takes as new initial data at time $`|I|`$ the triple $`(\stackrel{~}{u}(|I|)+w(|I|),\stackrel{~}{n_\pm }(|I|)+m_\pm (|I|))`$ and repeats the argument in $`[|I|,2|I|]`$. Of course in each step the involved norms have to be controlled in order to be able to choose intervals of equal length.
We collect some elementary facts about the spaces $`X^{s,b}`$ and $`Y^s`$. If $`u`$ is a solution of (4) with $`u(0)=f`$ we have for $`b0`$ :
$$\psi _1u_{X^{s,b}}cf_{H_x^s}$$
(5)
If $`v`$ is the solution of the problem
$$iv_t+\varphi (i_x)u=F,v(0)=0$$
we have for $`b^{}+1b0b^{}>1/2`$:
$$\psi _\delta v_{X^{s,b}}c\delta ^{1+b^{}b}F_{X^{s,b^{}}}$$
(6)
and if $`b^{}+1b0b^{}`$ :
$$\psi _\delta v_{X^{s,b}}c(\delta ^{1+b^{}b}F_{X^{s,b^{}}}+\delta ^{\frac{1}{2}b}F_{Y^s})$$
(7)
(for a proof see , Lemma 2.1).
Here the cut-off function $`\psi `$ is in $`C_0^{\mathrm{}}(𝐑)`$ with $`supp\psi (2,2)`$ , $`\psi 1`$ on $`[1,1],`$ $`\psi (t)=\psi (t)`$ , $`\psi (t)0`$ , $`\psi _\delta (t):=\psi (t/\delta )`$ if $`0<\delta 1`$.
We have $`X^{s,b}(I)C^0(I,H^s(𝐑))`$ if $`b>1/2`$ , $`I𝐑`$.
Moreover if $`w(t)=_0^te^{i(ts)_x^2}F(s)𝑑s`$ and $`f(s)=e^{is_x^2}F(s)`$ we have by , Lemma 2.2, especially (2.35)
$`w(t)_{L^2(𝐑)}={\displaystyle _0^t}e^{is_x^2}F(s)𝑑s_{L^2(𝐑)}={\displaystyle _0^t}f(s)𝑑s_{L^2(𝐑)}`$ (8)
$`c<\tau >^1\widehat{f}(\tau )_{L_\xi ^2L_\tau ^1}=c<\tau >^1(e^{is_x^2}F(s))_{L_\xi ^2L_\tau ^1}=cF_{Y^0(I)}`$
if $`t[0,|I|]`$ with $`|I|1`$.
Thus if $`FY^0(I)`$ we have $`wC^0(I,L^2(𝐑))`$ and if $`|I|1`$
$$w_{L^{\mathrm{}}(I,L^2(𝐑))}cF_{Y^0(I)}$$
Similarly if $`FY^1(I)`$ we have $`wC^0(I,H^{1,2}(𝐑))`$ and if $`|I|1`$
$$w_{L^{\mathrm{}}(I,H^{1,2}(𝐑))}cF_{Y^1(I)}$$
(9)
Next we need an interpolation property for the spaces $`X^{s,b}(I)`$. It is well-known that
$$H_t^b(𝐑,H_x^s(𝐑))=(H_{t_0}^{b_0}(𝐑,H_x^{s_0}(𝐑)),H_{t_1}^{b_1}(𝐑,H_x^{s_1}(𝐑)))_{[\theta ]}$$
(10)
where
$$0\theta 1,b=(1\theta )b_0+\theta b_1,s=(1\theta )s_0+\theta s_1$$
(11)
This also holds true if $`𝐑`$ is replaced by $`I`$ because the restriction operator from $`H_t^b(𝐑,H_x^s(𝐑))`$ onto $`H_t^b(I,H_x^s(𝐑))`$ is a retraction with a corresponding coretraction (extension). We refer to here. The following interpolation property is a consequence:
$$X^{s,b}(I)=(X^{s_0,b_0}(I),X^{s_1,b_1}(I))_{[\theta ]}$$
(12)
with $`s,b,\theta `$ as above. One only has to remark that
$$V_\varphi :H_t^b(I,H_x^s(𝐑))X^{s,b}(I)$$
defined by
$$V_\varphi f(x,t):=e^{it\varphi (i_x)}f(x,t),tI$$
is an isometric isomorphism.
Finally we have the following consequence of the Strichartz inequalities in the case of the (1+1)-dimensional Schrödinger equation $`\varphi (\xi )=\xi ^2`$ :
$$f_{L_t^q(𝐑,L_x^r(𝐑))}cf_{X^{0,b}}$$
(13)
and
$$f_{X^{0,b}}f_{L_t^q^{}(𝐑,L_x^r^{}(𝐑))}$$
(14)
where $`1/q+1/q^{}=1,\mathrm{\hspace{0.17em}1}/r+1/r^{}=1`$ and
$$b_0>1/2,\mathrm{\hspace{0.17em}0}bb_0,\mathrm{\hspace{0.17em}1}/2\eta 1,\mathrm{\hspace{0.17em}2}/q=1\eta b/b_0,\mathrm{\hspace{0.17em}1}/21/r=(1\eta )b/b_0$$
(15)
See , Lemma 2.4 with $`\nu `$ = 1 , plus duality.
We use the following notation for $`\lambda 𝐑`$:
$`<\lambda >:=(1+\lambda ^2)^{1/2},[\lambda ]_+:=\lambda `$ if $`\lambda >0`$ , $`=ϵ`$ if $`\lambda =0`$ , $`=0`$ if $`\lambda <0`$.
Acknowledgement: I am grateful to A. Grünrock for many helpful discussions and to T. Tao for informing the author of a gap in an earlier version of this paper.
## 1 Nonlinear estimates
Our aim here is to estimate the nonlinearities $`f=n_\pm u`$ in $`X^{s,a_1}`$ for given $`n_\pm X^{l,a}`$ and $`uX^{k,a_2}`$ for suitable $`s,l,k,a_1,a,a_2`$ and also in $`Y^s`$. We estimate $`\widehat{f}(\xi _1^{},\tau _1)=(\widehat{n_\pm }\widehat{u})(\xi _1^{},\tau _1)`$ in terms of $`\widehat{n_\pm }(\xi ,\tau )`$ and $`\widehat{u}(\xi _2^{},\tau _2),`$ where $`\xi =\xi _1^{}\xi _2^{},`$ $`\tau =\tau _1\tau _2`$. We also introduce the variables $`\sigma _1=\tau _1+\xi _{1}^{}{}_{}{}^{2}`$ , $`\sigma _2=\tau _2+\xi _{2}^{}{}_{}{}^{2}`$ , $`\sigma =\tau \pm |\xi |`$ so that
$$z:=\xi _{1}^{}{}_{}{}^{2}\xi _{2}^{}{}_{}{}^{2}|\xi |=\sigma _1\sigma _2\sigma $$
(16)
Define $`\widehat{v_2}=<\xi _2^{}>^k<\sigma _2>^{a_2}\widehat{u},\widehat{v}=<\xi >^l<\sigma >^a\widehat{n_\pm }`$ so that $`u_{X^{k,a_2}}=v_2_2`$ , $`n_\pm _{X^{l,a}}=v_2`$ .
In order to estimate $`f`$ in $`X^{s,a_1}`$ we take its scalar product with a function in $`X^{s,a_1}`$ with Fourier transform $`<\xi _1^{}>^s<\sigma _1>^{a_1}\widehat{v_1}`$ , $`v_1L^2`$. In the sequel we want to show an estimate of the type
$$|S|cv_2v_1_2v_2_2$$
(17)
where
$$S:=\frac{|\widehat{v}\widehat{v_1}\widehat{v_2}|<\xi _1^{}>^s}{<\sigma >^a<\sigma _1>^{a_1}<\sigma _2>^{a_2}<\xi _2^{}>^k<\xi >^l}𝑑\xi _1^{}𝑑\xi _2^{}𝑑\tau _1𝑑\tau _2$$
(18)
This directly gives the desired estimate
$$n_\pm u_{X^{s,a_1}}cn_\pm _{X^{l,a}}u_{X^{k,a_2}}$$
(19)
###### Proposition 1.1
The estimate (19) holds under the following conditions:
$`k,l0`$ (20)
$`sk<\mathrm{min}(2a{\displaystyle \frac{1}{2}},2a_1{\displaystyle \frac{1}{2}},2(a+a_1){\displaystyle \frac{3}{2}})`$ (21)
$`sl2a_1`$ (22)
$`a,a_1,a_2>1/4,a_11/2`$ (23)
$`a+a_1,a+a_2,a_1+a_2>3/4`$ (24)
$`k+a_1,k+a_2>1/2,k+a_1+a_2>1`$ (25)
Remark 1: We simplify (16) in the following way:
if (16) holds with the minus-sign and if $`\xi _1^{}\xi _2^{}`$ (resp. $`\xi _1^{}\xi _2^{}`$) we have
$$z=\xi _{2}^{}{}_{}{}^{2}\xi _{2}^{}{}_{}{}^{2}|\xi _1^{}\xi _2^{}|=(\xi _1^{}\frac{1}{2})^2(\xi _2^{}\frac{1}{2})^2=\xi _1^2\xi _2^2$$
where $`\xi _i:=\xi _i^{}1/2`$.
Thus the region $`\xi _1^{}\xi _2^{}`$ (resp. $`\xi _1^{}\xi _2^{}`$) of $`S`$ is majorized by
$$\overline{S}=\frac{|\widehat{v}(\xi ,\tau )\widehat{v_1}(\xi _1\pm \frac{1}{2},\tau _1)\widehat{v_2}(\xi _2\pm \frac{1}{2},\tau _2)|<\xi _1>^s}{<\sigma >^a<\sigma _1>^{a_1}<\sigma _2>^{a_2}<\xi _2>^k<\xi >^l}𝑑\xi _1𝑑\xi _2𝑑\tau _1𝑑\tau _2$$
(26)
where now
$`z=\xi _1^2\xi _2^2=\sigma _1\sigma _2\sigma ,\xi =\xi _1\xi _2,\tau =\tau _1\tau _2`$
$`\sigma _i=\tau _i+(\xi _i\pm 1/2)^2,\sigma =\tau +|\xi |=\tau +|\xi _1\xi _2|`$
Also the plus-sign in (16) can be treated similarly by again defining $`\xi _i=\xi _i^{}\pm 1/2`$. If one wants to estimate $`\overline{S}`$ by $`v_2v_1_2v_2_2`$ the variables $`\xi _i`$ and $`\xi _i\pm 1/2`$ of $`\widehat{v_i}`$ are completely equivalent, thus we do not distinguish between them.
Remark 2: We use the following application of Schwarz’ inequality: in order to estimate
$$I=|\widehat{v}(\zeta )\widehat{v_1}(\zeta _1)\widehat{v_2}(\zeta _2)K(\zeta _1,\zeta _2)|𝑑\zeta _1𝑑\zeta _2$$
where $`\zeta =(\xi ,\tau ),\zeta _i=(\xi _i,\tau _i),\zeta =\zeta _1\zeta _2`$ one has
$`|I|^2`$ $``$ $`v_2^2{\displaystyle \left(|\widehat{v_1}(\zeta +\zeta _2)\widehat{v_2}(\zeta _2)K(\zeta +\zeta _2,\zeta _2)|𝑑\zeta _2\right)^2𝑑\zeta }`$
$``$ $`v_2^2\left(\underset{\zeta }{sup}{\displaystyle |K(\zeta +\zeta _2,\zeta _2)|^2𝑑\zeta _2}\right){\displaystyle |\widehat{v_1}(\zeta +\zeta _2)\widehat{v_2}(\zeta _2)|^2𝑑\zeta _2𝑑\zeta }`$
$`=`$ $`C^2v_2^2v_1_2^2v_2_2^2`$
with
$$C^2=\underset{\zeta }{sup}|K(\zeta _1,\zeta _2)|^2𝑑\zeta _2$$
(29)
where the integral runs over $`\zeta _2`$ (or $`\zeta _1`$) for fixed $`\zeta `$.
Similar estimates hold by circularly permuting the variables $`\zeta ,\zeta _1,\zeta _2`$.
Proof of Prop. 1.1: We estimate (26) in several subregions.
Region a: $`|\xi _1|2|\xi _2|`$
Case aa: $`\sigma _1`$ dominant, i.e. $`|\sigma _1||\sigma |`$ , $`|\sigma _1||\sigma _2|`$.
We show according to (1),(29):
$`C_1^2:=\underset{\xi _1,\sigma _1}{sup}<\sigma _1>^{2a_1}<\xi _1>^{2s}{\displaystyle \underset{\xi _1,\sigma _1\text{fixed}}{}}<\sigma >^{2a}<\sigma _2>^{2a_2}<\xi _2>^{2k}<\xi >^{2l}𝑑\xi _2𝑑\sigma _2`$
$`<\mathrm{}`$
We have
$`{\displaystyle <\sigma >^{2a}<\sigma _2>^{2a_2}𝑑\sigma _2}`$ $`=`$ $`{\displaystyle <\sigma _1\sigma _2\xi _1^2+\xi _2^2>^{2a}<\sigma _2>^{2a_2}𝑑\sigma _2}`$
$``$ $`c<\xi _1^2\xi _2^2\sigma _1>^{\alpha _1}`$
by (1) and , Lemma 4.2 with $`\alpha _1:=2\mathrm{min}(a,a_2)[12\mathrm{max}(a,a_2)]_+`$ if $`a+a_2>1/2`$ which holds by (23).
Thus
$$C_1^2c\underset{\xi _1,\sigma _1}{sup}<\sigma _1>^{2a_1}<\xi _1>^{2s}<\xi _2>^{2k}<\xi _1^2\xi _2^2\sigma _1>^{\alpha _1}<\xi >^{2l}𝑑\xi _2$$
Now using $`|\xi ||\xi _1||\xi _2|\frac{1}{2}|\xi _1|`$ , thus $`<\xi >^{2l}c<\xi _1>^{2l}`$ , and substituting $`y=\xi _2^2`$ , $`dy=2|y|^{1/2}d\xi _2`$ , $`|y|\frac{1}{4}\xi _1^2`$ we get
$`C_1^2`$ $``$ $`c\underset{\xi _1,\sigma _1}{sup}<\sigma _1>^{2a_1}<\xi _1>^{2(sl)}{\displaystyle _{\mathrm{}}^+\mathrm{}}|y|^{1/2}<y>^k`$
$`\chi _{\{|y|\frac{1}{4}\xi _1^2\}}<y(\xi _1^2\sigma _1)>^{\alpha _1}dy`$
According to , Lemma 4.1 the integral takes its maximum at $`\xi _1^2=\sigma _1`$ , so that
$$C_1^2c\underset{\xi _1,\sigma _1}{sup}<\sigma _1>^{2a_1}<\xi _1>^{2(sl)}_{\mathrm{}}^+\mathrm{}|y|^{1/2}<y>^{(k+\alpha _1)}𝑑y$$
The integral converges if $`k+\alpha _1>1/2`$. By definition $`\alpha _1=2a`$ or $`2a_2`$ or $`2(a+a_2)1`$ up to an $`ϵ`$-term. Now $`k+2a>1/2`$ by (20),(23), similarly $`k+2a_2>1/2,`$ and $`k+2(a+a_2)1>1/2`$ by (20),(24), thus $`k+\alpha _1>1/2`$. Moreover $`\frac{3}{4}\xi _1^2\xi _1^2\xi _2^2=\sigma _1\sigma _2\sigma 3|\sigma _1|`$ so that
$$C_1^2c\underset{\xi _1}{sup}<\xi _1>^{4a_1+2(sl)}<\mathrm{}$$
because $`4a_1+2(sl)0`$ by (22).
Case ab: $`\sigma _2`$ dominant, i.e. $`|\sigma _2||\sigma |`$ , $`|\sigma _2||\sigma _1|`$.
By (1),(29) we have to show
$`C_2^2:=\underset{\xi _2,\sigma _2}{sup}<\sigma _2>^{2a_2}<\xi _2>^{2k}{\displaystyle \underset{\sigma _2,\xi _2\text{fixed}}{}}<\xi _1>^{2s}<\sigma >^{2a}<\sigma _1>^{2a_1}<\xi >^{2l}𝑑\xi _1𝑑\sigma _1`$
$`<\mathrm{}`$
We have
$$C_2^2c\underset{\xi _2,\sigma _2}{sup}<\sigma _2>^{2a_2}<\xi _2>^{2k}\underset{\sigma _2,\xi _2\text{fixed}}{}<\xi _1>^{2(sl)}<\sigma >^{2a}<\sigma _1>^{2a_1}𝑑\xi _1𝑑\sigma _1$$
Substituting $`\xi _1`$ by $`z`$ with fixed $`\xi _2`$ gives $`z=\xi _1^2\xi _2^2`$ , $`dz=2\xi _1d\xi _1`$ and by (1) $`3\xi _2^2z=\sigma _1\sigma _2\sigma 3|\sigma _2|`$. Thus
$`C_2^2`$ $`c`$ $`\underset{\xi _2,\sigma _2}{sup}<\sigma _2>^{2a_2}<\xi _2>^{2k}`$
$`{\displaystyle _{3\xi _2^2}^{3|\sigma _2|}}<z>^{sl}|z|^{1/2}\left({\displaystyle <\sigma _1(\sigma _2+z)>^{2a}<\sigma _1>^{2a_1}𝑑\sigma _1}\right)𝑑z`$
Now the inner integral is estimated using , Lemma 4.2 by $`c<\sigma _2+z>^{\alpha _2}`$, if $`a+a_1>1/2`$ (which follows from (23)) with $`\alpha _2:=2\mathrm{min}(a,a_1)[12\mathrm{max}(a,a_1)]_+`$. This gives
$`C_2^2`$ $``$ $`c\underset{\xi _2,\sigma _2}{sup}<\sigma _2>^{2a_2}<\xi _2>^{2k}{\displaystyle _{3\xi _2^2}^{3|\sigma _2|}}<z>^{sl}|z|^{1/2}<z+\sigma _2>^{\alpha _2}𝑑z`$
$``$ $`c\underset{\sigma _2}{sup}<\sigma _2>^{2a_2}{\displaystyle _0^{3|\sigma _2|}}<z>^{sl}|z|^{1/2}<z+\sigma _2>^{\alpha _2}𝑑z`$
We split up the integral into the parts $`0z\frac{1}{2}|\sigma _2|`$ and $`\frac{1}{2}|\sigma _2|z3|\sigma _2|`$. Assuming w.l.o.g. $`sl0`$ we estimate the first part by $`c<\sigma _2>^{sl+\frac{1}{2}\alpha _2}`$ and the second part by $`c<\sigma _2>^{sl\frac{1}{2}+[1\alpha _2]_+}`$ which is the larger one. Thus
$$C_2^2c\underset{\sigma _2}{sup}<\sigma _2>^{2a_2+sl\frac{1}{2}+[1\alpha _2]_+}<\mathrm{}$$
if
$$sl2a_2+\frac{1}{2}[1\alpha _2]_+$$
(30)
Now $`\alpha _2=2a`$ or $`2a_1`$ or $`2(a+a_1)1`$ and we have
$`sl`$ $``$ $`2a_1<2a_1+2a_2{\displaystyle \frac{1}{2}}\text{by (}\text{22}\text{),(}\text{23}\text{)}`$
$`sl`$ $``$ $`2a_1<2a_1+[2(a+a_2){\displaystyle \frac{3}{2}}]=2a_2+{\displaystyle \frac{1}{2}}+2(a+a_1)2\text{by (}\text{22}\text{),(}\text{23}\text{),(}\text{24}\text{)}`$
$`sl`$ $``$ $`1<1+[2(a_2+a){\displaystyle \frac{3}{2}}]=2a_2+2a{\displaystyle \frac{1}{2}}\text{by (}\text{22}\text{),(}\text{23}\text{),(}\text{24}\text{)}`$
$`sl`$ $``$ $`1<2a_2+{\displaystyle \frac{1}{2}}\text{by (}\text{22}\text{),(}\text{23}\text{)}`$
Thus (30) is satisfied.
Case ac: $`\sigma `$ dominant, i.e. $`|\sigma ||\sigma _1|`$ , $`|\sigma ||\sigma _2|`$
We have to show that
$`C^2:=\underset{\xi ,\sigma }{sup}<\sigma >^{2a}<\xi >^{2l}{\displaystyle \underset{\sigma ,\xi \text{fixed}}{}}<\xi _2>^{2k}<\sigma _1>^{2a_1}<\sigma _2>^{2a_2}<\xi _1>^{2s}𝑑\xi _2𝑑\sigma _2`$
$`<\mathrm{}`$
Using $`|\xi |\frac{1}{2}|\xi _1|`$ and $`|\xi ||\xi _1|+|\xi _2|\frac{3}{2}|\xi _1|`$, thus
$$\frac{1}{3}|\xi |^2\frac{3}{4}\xi _1^2\xi _1^2\xi _2^2=z=\sigma _1\sigma _2\sigma 3|\sigma |$$
we get
$$C^2c\underset{\xi ,\sigma }{sup}<\xi >^{4a+2(sl)}\underset{\sigma ,\xi \text{fixed}}{}<\xi _2>^{2k}<\sigma _1>^{2a_1}<\sigma _2>^{2a_2}𝑑\xi _2𝑑\sigma _2$$
Substituting $`\xi _2`$ by $`z`$ for fixed $`\xi `$ gives $`z=\xi _1^2\xi _2^2=(\xi _1+\xi _2)\xi =(\xi +2\xi _2)\xi `$ , $`\frac{dz}{d\xi _2}=2\xi `$ and $`z\xi ^2=(\xi _1^2\xi _2^2)\xi ^2=(\xi _1+\xi _2)\xi (\xi _1\xi _2)\xi =2\xi _2\xi `$ thus $`\xi _2=\frac{z\xi ^2}{2\xi }`$ leads to
$$C^2c\underset{\xi ,\sigma }{sup}<\xi >^{4a+2(sl)}|\xi |^1\underset{0}{\overset{3\xi ^2}{}}<\frac{z\xi ^2}{2\xi }>^{2k}(<\sigma _1>^{2a_1}<\sigma _2>^{2a_2}𝑑\sigma _2)𝑑z$$
We used here
$$z=\xi _1^2\xi _2^2=\xi (\xi _1+\xi _2)\frac{3}{2}|\xi ||\xi _1|3\xi ^2$$
(31)
Now
$`{\displaystyle <\sigma _1>^{2a_1}<\sigma _2>^{2a_2}𝑑\sigma _2}`$ $`=`$ $`{\displaystyle <\sigma _2(z\sigma )>^{2a_1}<\sigma _2>^{2a_2}𝑑\sigma _2}`$
$``$ $`c<z+\sigma >^\alpha `$
by , Lemma 4.2 with $`\alpha =2\mathrm{min}(a_1,a_2)[12\mathrm{max}(a_1,a_2)]_+`$ using $`a_1+a_2>1/2`$ which holds by (23).
Substitute $`y=z\xi ^2`$ and use $`|y||z|+\xi ^24\xi ^2`$ by (31) to conclude
$`C^2`$ $``$ $`c\underset{\xi ,\sigma }{sup}<\xi >^{4a+2(sl)}|\xi |^1{\displaystyle _{\mathrm{}}^{\mathrm{}}}\chi _{\{|y|4\xi ^2\}}<{\displaystyle \frac{y}{2|\xi |}}>^{2k}<y+\xi ^2+\sigma >^\alpha 𝑑y`$
$`=`$ $`c\underset{\xi }{sup}<\xi >^{4a+2(sl)}|\xi |^1{\displaystyle _0^{4\xi ^2}}<{\displaystyle \frac{y}{2|\xi |}}>^{2k}<y>^\alpha 𝑑y`$
by , Lemma 4.1.
If $`|\xi |1`$ we directly get
$$C^2c\underset{|\xi |1}{sup}|\xi |^1_0^{4\xi ^2}<\frac{y}{2|\xi |}>^{2k}𝑑yc\underset{|\xi |1}{sup}|\xi |^1|\xi |^2<\mathrm{}$$
If $`|\xi |1`$ we get
$`C^2`$ $``$ $`c\underset{|\xi |1}{sup}|\xi |^{4a+2(sl)1}\left({\displaystyle _0^{|\xi |}}<y>^\alpha 𝑑y+{\displaystyle _{|\xi |}^{4\xi ^2}}|y|^{2k\alpha }𝑑y|\xi |^{2k}\right)`$
$``$ $`c\underset{|\xi |1}{sup}(|\xi |^{4a+2(sl)1+[1\alpha ]_+}+|\xi |^{4a+2(sl)1+2k}|\xi |^{2k\alpha +1})`$
$``$ $`c\underset{|\xi |1}{sup}|\xi |^{4a+2(sl)1+[1\alpha ]_+}<\mathrm{}`$
if
$$2k+\alpha >1$$
(32)
and
$$4a+2(sl)1+[1\alpha ]_+0$$
(33)
By the definition of $`\alpha `$ (32) holds if $`2k+2a_1>1`$ and $`2k+2a_2>1`$ and $`2k+2(a_1+a_2)1>1`$ which follows from (25).
In order to show (33) we use $`\alpha =2a_1`$ or $`2a_2`$ or $`2(a_1+a_2)1`$ and get $`2(sl)2<4a+1`$ and $`2(sl)2<2(a+a_1)+2a=4a+2a_1`$ and $`2(sl)2<2(a+a_2)+2a=4a+2a_2`$ and $`2(sl)2<2(a+a_1)+2(a+a_2)1`$ by (22),(23),(24), which implies (33).
Region b: $`|\xi _1|2|\xi _2|`$
We show
$`C_2^2:=\underset{\xi _2,\sigma _2}{sup}<\sigma _2>^{2a_2}<\xi _2>^{2k}{\displaystyle \underset{\sigma _2,\xi _2\text{fixed}}{}}<\xi _1>^{2s}<\sigma >^{2a}<\sigma _1>^{2a_1}<\xi >^{2l}𝑑\xi _1𝑑\sigma _1`$
$`<\mathrm{}`$
We have
$$C_2^2c\underset{\xi _2,\sigma _2}{sup}<\xi _2>^{2k}\underset{\sigma _2,\xi _2\text{fixed}}{}<\xi _1>^{2s}<\sigma >^{2a}<\sigma _1>^{2a_1}𝑑\xi _1𝑑\sigma _1$$
Now by (1) and , Lemma 4.2:
$`{\displaystyle <\sigma >^{2a}<\sigma _1>^{2a_1}𝑑\sigma _1}`$
$`={\displaystyle <\sigma _1(\sigma _2+\xi _1^2\xi _2^2)>^{2a}<\sigma _1>^{2a_1}𝑑\sigma _1}`$
$`c<\sigma _2+\xi _1^2\xi _2^2>^{\alpha _2}`$
with $`\alpha _2`$ as above.
The substitution $`y=\xi _1^2,`$ $`dy=2|y|^{1/2}d\xi _1,`$ $`y4\xi _2^2`$ gives
$`C_2^2`$ $``$ $`c\underset{\xi _2,\sigma _2}{sup}{\displaystyle <\xi _1>^{2(sk)}<\sigma _2+\xi _1^2\xi _2^2>^{\alpha _2}𝑑\xi _1}`$
$``$ $`c\underset{\xi _2,\sigma _2}{sup}{\displaystyle _0^{4\xi _2^2}}<y>^{sk}<y(\xi _2^2\sigma _2)>^{\alpha _2}|y|^{1/2}𝑑y`$
$``$ $`c\underset{\xi _2,\sigma _2}{sup}{\displaystyle _{\mathrm{}}^+\mathrm{}}<y>^{sk}|y|^{1/2}\chi _{\{|y|4\xi _2^2\}}<y(\xi _2^2\sigma _2)>^{\alpha _2}𝑑y`$
$``$ $`c{\displaystyle _0^{\mathrm{}}}<y>^{sk}|y|^{1/2}<y>^{\alpha _2}𝑑y`$
by use of , Lemma 4.1 (remark that $`sk1/2`$ by (21),(23)).
Thus $`C_2^2<\mathrm{}`$ provided $`sk<\alpha _2\frac{1}{2}`$ which follows from (21) and completes the proof.
Our next aim is to give a similar estimate for $`f=n_\pm u`$ in $`Y^s`$. We first integrate $`<\sigma _1>^1\widehat{f}`$ over $`\tau _1`$ and take the scalar product with a function in $`H_x^k(𝐑)`$ with Fourier transform $`<\xi _1>^k\widehat{w_1}`$ , $`w_1L_x^2(𝐑)`$. We show that an estimate of the type
$$|\stackrel{~}{S}|cv_2w_1_2v_2_2$$
holds, where
$$\stackrel{~}{S}:=\frac{|\widehat{v}\widehat{w_1}\widehat{v_2}|<\xi _1^{}>^s}{<\sigma >^a<\sigma _1><\sigma _2>^{a_2}<\xi _2^{}>^k<\xi >^l}𝑑\xi _1^{}𝑑\xi _2^{}\tau _1𝑑\tau _2$$
and the notation is the same as for $`S`$ before.
This directly gives the estimates
$$n_\pm u_{Y^s}cn_\pm _{X^{l,a}}u_{X^{k,a_2}}$$
(34)
###### Proposition 1.2
The estimate (34) holds under the following conditions:
$`k>0,l0`$ (35)
$`sk<\mathrm{min}(2a{\displaystyle \frac{1}{2}},{\displaystyle \frac{1}{2}})`$ (36)
$`sl1`$ (37)
$`a,a_2>1/4`$ (38)
$`a+a_2>3/4`$ (39)
$`k+a_2>1/2`$ (40)
Remark: The same remarks as for Prop. 1.1 apply. Thus we estimate
$$\overline{\stackrel{~}{S}}=\frac{|\widehat{v}(\xi ,\tau )\widehat{w_1}(\xi _1\pm 1/2)\widehat{v_2}(\xi _2\pm 1/2,\tau _2)|<\xi _1>^s}{<\sigma >^a<\sigma _1><\sigma _2>^{a_2}<\xi _2>^k<\xi >^l}𝑑\xi _1𝑑\xi _2𝑑\tau _1𝑑\tau _2$$
(41)
where again (1) holds with the same notation again.
Proof of Prop.1.2: The proof works along the same lines as the foregoing one. Choose $`a_1`$ such that
$$1/2>a_1>\mathrm{min}(\frac{1}{4},\frac{3}{4}a,\frac{3}{4}a_2,\frac{5}{4}aa_2,\frac{1}{2}k,1a_2k,\frac{sk}{2}+\frac{1}{4},\frac{sk}{2}a+\frac{3}{4})$$
(42)
This and (35)-(40) imply the conditions (20),(21),(23),(24),(25) and
$$a+a_1+a_2>5/4$$
(43)
We consider the regions of integration similarly as in the proof of Prop.1.1 with the minor variation that case aa means $`|\sigma _1|4|\sigma _2|`$ , $`|\sigma _1|4|\sigma |`$ , case ab is given by $`|\sigma _2|\frac{1}{4}|\sigma _1|`$ , $`|\sigma _2||\sigma |`$ , and case ac by $`|\sigma |\frac{1}{4}|\sigma _1|`$ , $`|\sigma ||\sigma _2|`$.
Case aa: $`|\xi _1|2|\xi _2|`$
Choose
$$\widehat{v_1}(\xi _1\pm 1/2,\tau _1):=<\sigma _1>^{1/2}\widehat{w_1}(\xi _1\pm 1/2)\chi _{\{\frac{1}{2}|\sigma _1|\xi _1^22|\sigma _1|\}}$$
Now in the considered region we have $`z=\xi _1^2\xi _2^2\frac{3}{4}\xi _1^2`$ and
$$z=\sigma _1\sigma _2\sigma =|\sigma _1\sigma _2\sigma ||\sigma _1||\sigma _2||\sigma ||\sigma _1|\frac{1}{4}|\sigma _1|\frac{1}{4}|\sigma _1|=\frac{1}{2}|\sigma _1|$$
Moreover
$$\frac{4}{3}z=\frac{4}{3}(\sigma _1\sigma _2\sigma )\frac{4}{3}(|\sigma _1|+|\sigma _2|+|\sigma |)\frac{4}{3}(|\sigma _1|+\frac{1}{4}|\sigma _1|+\frac{1}{4}|\sigma _1|)=2|\sigma _1|$$
altogether
$$\frac{1}{2}|\sigma _1|z\xi _1^2\frac{4}{3}z2|\sigma _1|$$
This means that in the case at hand $`\chi 1`$ , so that
$$\overline{\stackrel{~}{S}}=\frac{|\widehat{v}(\xi ,\tau )\widehat{v_1}(\xi _1\pm 1/2,\tau _1)\widehat{v_2}(\xi _2\pm 1/2,\tau _2)|<\xi _1>^s}{<\sigma >^a<\sigma _1>^{1/2}<\sigma _2>^{a_2}<\xi _2>^k<\xi >^l}𝑑\xi _1𝑑\xi _2𝑑\tau _1𝑑\tau _2$$
where the integral runs over the region aa. This is exactly the term treated in case aa of Prop. 1.1 with $`a_1`$ replaced by $`1/2`$. Since all the assumptions of Prop. 1.1 are satisfied with $`a_1=1/2`$ under our assumptions (35) - (40) we get from that proof an estimate by $`cv_2v_1_2v_2_2`$. Because we want to have an estimate by $`cv_2w_1_2v_2_2`$, the only thing to be checked is $`v_1_2cw_1_2`$. But this is true, namely
$`v_1_2^2`$ $`=`$ $`{\displaystyle \underset{\frac{1}{2}|\sigma _1|\xi _1^22|\sigma _1|}{}<\tau _1+(\xi _1\pm 1/2)^2>^1|\widehat{w_1}(\xi _1\pm 1/2)|^2𝑑\tau _1𝑑\xi _1}`$
$`=`$ $`{\displaystyle (\underset{\frac{1}{2}\xi _1^2|\sigma _1|2\xi _1^2}{}<\sigma _1>^1𝑑\sigma _1)|\widehat{w_1}(\xi _1\pm 1/2)|^2𝑑\xi _1}`$
Now the inner integral is bounded independently of $`\xi _1`$ by some logarithm. Thus $`v_1_2cw_1_2`$. This concludes the proof of case aa.
In all other cases we define $`\widehat{v_1}:=<\sigma _1>^{a_11}\widehat{w_1}`$ with $`a_1`$ as above. Then we simply have
$$v_1_2^2=<\sigma _1>^{2(a_11)}|\widehat{w_1}(\xi _1)|^2𝑑\sigma _1𝑑\xi _1cw_1_2^2$$
because $`a_1<1/2`$ by our choice (42).
Now with this choice $`\overline{\stackrel{~}{S}}`$ reduces to $`\overline{S}`$ so that it only remains to check that the old estimates of Prop. 1.1 in all other cases (with the modified cases aa,ab,ac as above) remain true. According to the remarks at the beginning of this proof we have to avoid using (22) in its strong form, namely (37) and (43) should suffice.
Case ab: We have $`z=\sigma _1\sigma _2\sigma |\sigma _1|+|\sigma _2|+|\sigma |4|\sigma _2|+|\sigma _2|+|\sigma _2|=6|\sigma _2|`$ instead of $`3|\sigma _2|`$ which appears as upper limit of integration and makes no essential difference. We have to check (30) without use of (22): $`sl1<2a_1+2a_21/2`$ by (37)(24), $`sl1<2a_2+1/2+2(a+a_1)2`$ by (37)(43), $`sl1<2a_2+2a1/2`$ by (37)(39), $`sl1<2a_2+1/2`$ by (37)(38).
Case ac: We have $`z=\sigma _1\sigma _2\sigma |\sigma _1|+|\sigma _2|+|\sigma |4|\sigma |+|\sigma |+|\sigma |=6|\sigma |`$ instead of $`3|\sigma |`$ which makes no essential difference. In order to check (33) we only used (37) instead of (22).
Region b: remains unchanged.
The proof is complete.
## 2 Energy bounds and decomposition of data
The system (1),(2),(3) is now transformed into a system of first order in $`t`$ in the usual way. Defining $`A=\frac{d^2}{dx^2}`$ and
$$n_\pm :=n\pm iA^{1/2}n_t$$
(44)
we have
$`n`$ $`=`$ $`{\displaystyle \frac{1}{2}}(n_++n_{})`$ (45)
$`2iA^{1/2}n_t`$ $`=`$ $`n_+n_{}`$ (46)
and the equivalent problem reads as follows
$`iu_t+u_{xx}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(n_++n_{})u`$ (47)
$`in_{\pm t}A^{1/2}n_\pm `$ $`=`$ $`\pm A^{1/2}(|u|^2)_{xx}`$
with initial data
$$u(0)=u_0,n_\pm (0)=n_0\pm iA^{1/2}n_1$$
(48)
The standard conservation laws for the original system are: conservation of the $`L^2`$-norm $`u(t)=:M(u)=M`$ and the energy
$$E:=E(u,n,n_t):=u_x(t)^2+1/2(n(t)^2+A^{1/2}n_t(t)^2)+_{\mathrm{}}^+\mathrm{}n(t)|u(t)|^2𝑑x$$
Now we have by Gagliardo-Nirenberg and $`L^2`$-conservation:
$`\left|{\displaystyle n|u|^2𝑑x}\right|`$ $``$ $`{\displaystyle \frac{1}{4}}{\displaystyle n^2𝑑x}+c{\displaystyle |u|^4𝑑x}{\displaystyle \frac{1}{4}}n^2+cu_xu^3`$
$``$ $`{\displaystyle \frac{1}{4}}(n^2+u_x^2)+cu^6={\displaystyle \frac{1}{4}}(n^2+u_x^2)+cu_0^6`$
This implies
$$u_x(t)^2+\frac{1}{2}(n(t)^2+A^{1/2}n_t(t)^2)E+\frac{1}{4}(n(t)^2+u_x(t)^2)+c_1M^6$$
(49)
consequently
$`u_x(t)^2`$ $``$ $`{\displaystyle \frac{4}{3}}(E+c_1M^6)`$ (50)
$`n(t)^2+A^{1/2}n_t(t)^2`$ $``$ $`4(E+c_1M^6)`$ (51)
We also have
$`E(u,n,n_t)`$ $``$ $`u_x(t)^2+{\displaystyle \frac{1}{2}}(n(t)^2+A^{1/2}n_t(t)^2)+\left|{\displaystyle _{\mathrm{}}^{\mathrm{}}}n(t)|u(t)|^2𝑑x\right|`$ (52)
$``$ $`{\displaystyle \frac{5}{4}}u_x(t)^2+{\displaystyle \frac{3}{4}}(n(t)^2+A^{1/2}n_t(t)^2)+c_1u(t)^6`$
These estimates together with $`L^2`$ \- conservation of $`u`$ and local well-posedness for data in $`H^{1,2}\times L^2\times \dot{H}^{1,2}`$ which is given by implies directly also global well-posedness for these data.
Let now data be given with
$$u_0H^{s,2}(𝐑),n_0L^2(𝐑),n_1\dot{H}^{1,2}(𝐑),1>s>9/10$$
and decompose for $`N1`$ :
$$u_0=u_{01}+u_{02}$$
where
$`u_{01}`$ $`=`$ $`^1(\chi _{\{|\xi |N\}}\widehat{u_0}(\xi ))={\displaystyle _{|\xi |N}}e^{ix\xi }\widehat{u_0}(\xi )𝑑\xi `$
$`u_{02}`$ $`=`$ $`^1(\chi _{\{|\xi |N\}}\widehat{u_0}(\xi ))={\displaystyle _{|\xi |N}}e^{ix\xi }\widehat{u_0}(\xi )𝑑\xi `$
One easily shows that
$`u_{01}_{H^{l,2}}`$ $``$ $`cN^{ls}u_0_{H^{s,2}}\text{f}orls`$
$`u_{01}_{L^2}`$ $``$ $`u_0_{L^2}`$
$`u_{02}_{H^{l,2}}`$ $``$ $`cN^{ls}u_0_{H^{s,2}}\text{f}orls`$
$`u_{02}_{L^2}`$ $``$ $`cN^su_0_{H^{s,2}}`$
Thus we have the following global bounds for the solution $`(\stackrel{~}{u},\stackrel{~}{n})`$ of (1),(2) with data $`(u_{01},n_0,n_1)`$ by (52):
$$E(\stackrel{~}{u},\stackrel{~}{n},\stackrel{~}{n}_t)\frac{5}{4}u_{01_x}^2+\frac{3}{4}(n_0^2+A^{1/2}n_1^2)+c_1u_{01}^6\overline{c}N^{2(1s)}$$
(53)
and thus by (50),(51) and $`L^2`$-conservation of $`\stackrel{~}{u}`$:
$`\stackrel{~}{u}_x(t)^2+A^{1/2}\stackrel{~}{n}_t(t)^2+\stackrel{~}{n}(t)^2\widehat{c}N^{2(1s)}`$ (54)
$`\stackrel{~}{u}(t)M`$ (55)
The corresponding global solution $`(\stackrel{~}{u},\stackrel{~}{n}_\pm )`$ of (47) with data $`(u_{01},n_{0+},n_0)`$ therefore fulfills:
$`\stackrel{~}{u}_x(t)\widehat{c}N^{1s}`$ (56)
$`\stackrel{~}{n}_\pm (t)\widehat{c}N^{1s}`$ (57)
$`\stackrel{~}{u}(t)M`$ (58)
where $`\widehat{c}`$ depends essentially only on $`\overline{c}`$ (the initial energy) and $`M`$ on the initial $`L^2`$-norm of $`\stackrel{~}{u}`$.
## 3 Further bounds for the regular part
In order to give further estimates of $`(\stackrel{~}{u},\stackrel{~}{n}_\pm )`$ we consider the system of integral equations which belongs to problem (47) with data $`(u_{01},n_+(0),n_{}(0))`$:
$`\stackrel{~}{u}(t)`$ $`=`$ $`e^{it_x^2}u_{01}i{\displaystyle _0^t}e^{i(ts)_x^2}{\displaystyle \frac{1}{2}}(\stackrel{~}{n}_+(s)+\stackrel{~}{n}_{}(s))\stackrel{~}{u}(s)𝑑s`$ (59)
$`\stackrel{~}{n}_\pm (t)`$ $`=`$ $`e^{itA^{1/2}}n_{0\pm }i{\displaystyle _0^t}e^{i(ts)A^{1/2}}A^{1/2}(|\stackrel{~}{u}(s)|^2)_{xx}𝑑s`$
We always assume $`tI=[0,|I|]`$. In this case we could, whenever helpful, place a factor $`\psi _1(t)`$ in front of the first terms on the r.h.sides and $`\psi _{|I|}(t)`$ in front of any of the integrals in (59),(3) without changing the equations at all. Here $`\psi C_0^{\mathrm{}}(𝐑)`$ is a non-negative cut-off function with $`\psi (t)=0`$ if $`|t|2`$ , $`\psi (t)=1`$ if $`|t|1`$ and $`\psi _\delta :=\psi (t/\delta )`$.
Important remark: Here and in the following section the constants denoted by $`c`$ or $`c_0`$ depend essentially only on $`\overline{c}`$ in (53) (and therefore on $`E(\stackrel{~}{u},\stackrel{~}{n},\stackrel{~}{n_t})`$ ) and on $`M`$ (in (58)).
The energy estimate (55),(56) gives
$$\stackrel{~}{u}_{X^{1,0}(I)}=\stackrel{~}{u}_{L^2(I,H^{1,2}(𝐑))}\stackrel{~}{u}_{L^{\mathrm{}}(I,H^{1,2}(𝐑))}|I|^{1/2}cN^{1s}|I|^{1/2}$$
(60)
By (57),(59),(5),(6) and (14):
(in the sequel $`a\pm `$ denotes a number slightly larger resp. smaller than $`a`$)
$`\stackrel{~}{u}_{X^{0,\frac{1}{2}+}(I)}`$
$``$ $`cu_{01}_{L^2(𝐑)}+c(\stackrel{~}{n}_+\stackrel{~}{u}_{X^{0,\frac{1}{2}+}(I)}+\stackrel{~}{n}_{}\stackrel{~}{u}_{X^{0,\frac{1}{2}+}(I)})`$
$``$ $`c+c(\stackrel{~}{n}_+\stackrel{~}{u}_{L^{\frac{6}{5}+}(I,L^{\frac{6}{5}+}(𝐑))}+\stackrel{~}{n}_{}\stackrel{~}{u}_{L^{\frac{6}{5}+}(I,L^{\frac{6}{5}+}(𝐑))})`$
$``$ $`c+c\left({\displaystyle _I}\stackrel{~}{n}_+_{L_x^2}^{\frac{6}{5}+}\stackrel{~}{u}_{L_x^{3+}}^{\frac{6}{5}+}𝑑t+{\displaystyle _I}\stackrel{~}{n}_{}_{L_x^2}^{\frac{6}{5}+}\stackrel{~}{u}_{L_x^{3+}}^{\frac{6}{5}+}𝑑t\right)^{(\frac{6}{5}+)^1}`$
$``$ $`c+c(\stackrel{~}{n}_+_{L^{\mathrm{}}(I,L^2(𝐑))}+\stackrel{~}{n}_{}_{L^{\mathrm{}}(I,L^2(𝐑))})\left({\displaystyle _I}\stackrel{~}{u}_{L_x^2}^1\stackrel{~}{u}_x_{L_x^2}^{\frac{1}{5}+}𝑑t\right)^{(\frac{6}{5}+)^1}`$
$``$ $`c+c(\stackrel{~}{n}_+_{L^{\mathrm{}}(I,L^2(𝐑))}+\stackrel{~}{n}_{}_{L^{\mathrm{}}(I,L^2(𝐑))})\stackrel{~}{u}_{L_t^{\mathrm{}}(I,L_x^2)}^{\frac{5}{6}}\stackrel{~}{u}_x_{L_t^{\mathrm{}}(I,L_x^2)}^{\frac{1}{6}+}|I|^{\frac{5}{6}}`$
$``$ $`c+c(\stackrel{~}{n}_+_{L^{\mathrm{}}(I,L^2(𝐑))}+\stackrel{~}{n}_{}_{L^{\mathrm{}}(I,L^2(𝐑))})N^{\frac{1s}{6}+}|I|^{\frac{5}{6}}`$
$``$ $`c+cN^{1s}N^{\frac{1s}{6}+}N^{\frac{5}{6}4(1s)+}`$
$``$ $`2c`$
for $`N`$ sufficiently large. We here also used Gagliardo-Nirenberg and (55),(56) and assumed $`|I|cN^{4(1s)}`$.
Next we estimate $`\stackrel{~}{n}_\pm _{X^{0,\frac{1}{2}+}(I)}`$ in terms of $`\stackrel{~}{u}_{X^{0,\frac{1}{2}+}(I)}`$. From the integral equation (3),(5),(6) we have
$$\stackrel{~}{n}_\pm _{X^{0,\frac{1}{2}+}(I)}cn_{0\pm }_{L^2(𝐑)}+cA^{1/2}(|\stackrel{~}{u}|^2)_{xx}_{X^{0,\frac{3}{8}}(I)}|I|^{\frac{1}{8}}$$
(62)
By , Lemma 4.4 (with $`k=1/4,l=0,c=3/8,b_1=3/8`$) we have
$$A^{1/2}(|\stackrel{~}{u}|^2)_{xx}_{X^{0,\frac{3}{8}}(I)}c\stackrel{~}{u}_{X^{\frac{1}{4},\frac{3}{8}}(I)}^2c\stackrel{~}{u}_{X^{0,\frac{1}{2}+}(I)}^{\frac{3}{2}}\stackrel{~}{u}_{X^{1,0}(I)}^{\frac{1}{2}}$$
(63)
where the last estimate follows by interpolation from (12). From (60),(3),(62),(63) we get
$$\stackrel{~}{n}_\pm _{X^{0,\frac{1}{2}+}(I)}cn_{0\pm }_{L^2(𝐑)}+c\stackrel{~}{u}_{X^{0,\frac{1}{2}+}(I)}^{\frac{3}{2}}N^{\frac{1s}{2}}|I|^{\frac{3}{8}}c(N^{1s}+N^{\frac{1s}{2}})$$
(64)
Thus we have
###### Lemma 3.1
Let $`|I|N^{4(1s)}`$ and $`n_{0\pm }_{L^2(𝐑)}cN^{1s}`$. Then we have
$`\stackrel{~}{n}_\pm _{X^{0,\frac{1}{2}+}(I)}`$ $``$ $`cN^{1s}`$ (65)
$`\stackrel{~}{u}_{X^{0,\frac{1}{2}+}(I)}`$ $``$ $`c`$ (66)
The next step is an estimate of $`\stackrel{~}{u}_{X^{1,\frac{1}{2}}(I)}`$.
From the integral equation (59) we get by (5),(7):
$`\stackrel{~}{u}_{X^{1,\frac{1}{2}}(I)}`$ $``$ $`cu_{01}_{H^{1,2}(𝐑)}+c(\stackrel{~}{n}_+\stackrel{~}{u}_{X^{1,\frac{1}{2}}(I)}+\stackrel{~}{n}_{}\stackrel{~}{u}_{X^{1,\frac{1}{2}}(I)})`$
$`+c(\stackrel{~}{n}_+\stackrel{~}{u}_{Y^1(I)}+\stackrel{~}{n}_{}\stackrel{~}{u}_{Y^1(I)})`$
$``$ $`cN^{1s}+c(\stackrel{~}{n}_+_{X^{0,\frac{1}{2}+}(I)}+\stackrel{~}{n}_{}_{X^{0,\frac{1}{2}+}(I)})\stackrel{~}{u}_{X^{\frac{1}{2}+,\frac{1}{4}+}(I)}`$
$``$ $`cN^{1s}+c(\stackrel{~}{n}_+_{X^{0,\frac{1}{2}+}(I)}+\stackrel{~}{n}_{}_{X^{0,\frac{1}{2}+}(I)})\stackrel{~}{u}_{X^{0,\frac{1}{2}+}(I)}^{\frac{1}{2}}\stackrel{~}{u}_{X^{1,0}(I)}^{\frac{1}{2}+}`$
$``$ $`cN^{1s}+c(\stackrel{~}{n}_+_{X^{0,\frac{1}{2}+}(I)}+\stackrel{~}{n}_{}_{X^{0,\frac{1}{2}+}(I)})\stackrel{~}{u}_{X^{0,\frac{1}{2}+}(I)}^{\frac{1}{2}}N^{\frac{1s}{2}+}|I|^{\frac{1}{4}+}`$
Here we used Prop. 1.1 and Prop. 1.2 (with $`s=1,l=0,k=\frac{1}{2}+,a_1=\frac{1}{2},a=\frac{1}{2}+,a_2=\frac{1}{4}+`$), an interpolation argument and (60). Consequently we get
###### Lemma 3.2
If $`|I|N^{4(1s)}`$ and $`n_{0\pm }_{L^2(𝐑)}cN^{1s}`$, the following estimate holds
$$\stackrel{~}{u}_{X^{1,\frac{1}{2}}(I)}cN^{1s}$$
(68)
Proof: follows immediately from Lemma 3.1 and (3):
$$\stackrel{~}{u}_{X^{1,\frac{1}{2}}(I)}cN^{1s}+cN^{1s}N^{\frac{1s}{2}+}N^{4(1s)\frac{1}{4}+}cN^{1s}$$
Remark: If the data fulfill the conditions
$`u_{01}_{L^2(𝐑)}`$ $``$ $`c`$ (69)
$`u_{01_x}_{L^2(𝐑)}`$ $``$ $`cN^{1s}`$ (70)
$`n_{0+}_{L^2(𝐑)}+n_0_{L^2(𝐑)}`$ $``$ $`cN^{1s}`$ (71)
then the following estimates hold on $`|I|N^{4(1s)}`$:
$`\stackrel{~}{n}_\pm _{X^{0,\frac{1}{2}+}(I)}cN^{1s}`$ (72)
$`\stackrel{~}{u}_{X^{1,\frac{1}{2}}(I)}cN^{1s}`$ (73)
$`\stackrel{~}{u}_{X^{0,\frac{1}{2}+}(I)}c`$ (74)
This follows from Lemma 3.1, 3.2.
Also the estimates (56),(57),(58) hold under these assumptions.
## 4 The part with rough data
Let $`(u,n_+,n_{})`$ be a solution of (47) with data $`(u_0,n_{0+},n_0)`$ and $`(\stackrel{~}{u},\stackrel{~}{n}_+,\stackrel{~}{n}_{})`$ the solution with data $`(u_{01},n_{0+},n_0)`$.
Define $`v:=u\stackrel{~}{u},m_\pm :=n_\pm \stackrel{~}{n}_\pm .`$ Then $`(v,m_+,m_{})`$ fulfills
$`iv_t+v_{xx}`$ $`=`$ $`iu_t+u_{xx}i\stackrel{~}{u}_t\stackrel{~}{u}_{xx}={\displaystyle \frac{1}{2}}(n_++n_{})u{\displaystyle \frac{1}{2}}(\stackrel{~}{n}_++\stackrel{~}{n}_{})\stackrel{~}{u}`$ (75)
$`=`$ $`{\displaystyle \frac{1}{2}}(\stackrel{~}{n}_++m_++\stackrel{~}{n}_{}+m_{})(\stackrel{~}{u}+v){\displaystyle \frac{1}{2}}(\stackrel{~}{n}_++\stackrel{~}{n}_{})\stackrel{~}{u}`$
$`=`$ $`{\displaystyle \frac{1}{2}}(\stackrel{~}{n}_++\stackrel{~}{n}_{})v+{\displaystyle \frac{1}{2}}(m_++m_{})v+{\displaystyle \frac{1}{2}}(m_++m_{})\stackrel{~}{u}`$
$`=`$ $`F_1+F_2+F_3=:F`$
and
$`im_{\pm t}A^{1/2}m_\pm =in_{\pm t}i\stackrel{~}{n}_{\pm t}A^{1/2}n_\pm \pm A^{1/2}\stackrel{~}{n}_\pm `$
$`=\pm A^{1/2}(|u|^2)_{xx}A^{1/2}(|\stackrel{~}{u}|^2)_{xx}`$
$`=\pm A^{1/2}((\stackrel{~}{u}+v)(\overline{\stackrel{~}{u}}+\overline{v}))_{xx}A^{1/2}(\stackrel{~}{u}\overline{\stackrel{~}{u}})_{xx}`$
$`=\pm A^{1/2}(\stackrel{~}{u}\overline{v})_{xx}\pm A^{1/2}(|v|^2)_{xx}\pm A^{1/2}(v\overline{\stackrel{~}{u}})_{xx}`$
$`=:G_1+G_2+G_3`$
$`=:G`$ (76)
Furthermore
$`v(0)=u(0)\stackrel{~}{u}(0)=u_0u_{01}=u_{02}`$ (78)
$`m_\pm (0)=n_\pm (0)\stackrel{~}{n}_\pm (0)=0`$
The corresponding system of integral equations reads as follows
$`v(t)`$ $`=`$ $`e^{it_x^2}u_{02}i{\displaystyle _0^t}e^{i(ts)_x^2}F(s)𝑑s`$ (79)
$`m_\pm (t)`$ $`=`$ $`i{\displaystyle _0^t}e^{i(ts)A^{1/2}}G(s)𝑑s`$ (80)
Here we have $`u_{02}H^{s,2}(𝐑)`$ with
$`u_{02}_{H^{s,2}}cu_0_{H^{s,2}}c`$ (81)
$`u_{02}_{L^2}cN^su_0_{H^{s,2}}cN^s`$ (82)
We construct a solution of (79),(80) in some time interval $`I`$ by the standard contraction mapping principle. We define
$$w(t):=i_0^te^{i(ts)_x^2}F(s)ds,z_\pm (t):=\text{r.h.s. of (}\text{80}\text{)}$$
(83)
and a mapping $`S=(S_0,S_+,S_{})`$ by
$`(S_0v)(t)`$ $`:=`$ $`e^{it_x^2}u_{02}+w(t)`$
$`(S_\pm m_\pm )(t)`$ $`:=`$ $`z_\pm (t)`$
###### Proposition 4.1
For $`9/10<s<1`$ and given data $`u_{02}H^{s,2}(𝐑)`$ with (81),(82) and $`u_{01},n_{0\pm }`$ as in (69),(70),(71) the system of integral equations (79),(80) has a unique solution $`(v,m_\pm )X^{s,\frac{1}{2}+}(I)\times X^{0,\frac{1}{2}+}(I)`$ in the same interval $`I`$ with $`|I|=N^{4(1s)\delta }`$, $`\delta >0`$ of the preceding section, which fulfills
$`v_{X^{0,\frac{1}{2}+}(I)}`$ $``$ $`cN^s`$ (84)
$`v_{X^{s,\frac{1}{2}+}(I)}`$ $``$ $`c`$ (85)
$`m_\pm _{X^{\frac{1}{2},\frac{1}{2}+}(I)}`$ $``$ $`cN^s`$ (86)
$`m_\pm _{X^{0,\frac{1}{2}+}(I)}`$ $``$ $`cN^{\frac{1}{2}\frac{1}{4}s\frac{\delta }{4}+}`$ (87)
Proof: We want to use Banach’s fixed point theorem in the set $`Z`$, where
$`Z`$ $`:=`$ $`\{v_{X^{0,\frac{1}{2}+ϵ}(I)}c_0N^s,v_{X^{s,\frac{1}{2}+ϵ}(I)}c_0`$
$`m_\pm _{X^{\frac{1}{2},\frac{1}{2}+ϵ}(I)}c_0N^s,m_\pm _{X^{0,\frac{1}{2}+ϵ}(I)}c_0N^{\frac{1}{2}\frac{1}{4}s\frac{\delta }{4}+}\}`$
with its natural metric, $`c_0`$ chosen below.
Now take any $`(v,m_+,m_{})Z`$. In order to show $`(S_0v,S_+m_+,S_{}m_{})Z`$ we estimate $`S_0v_{X^{s,\frac{1}{2}+ϵ}(I)}`$ first.
We have by Prop. 1.1 (with $`k=s\frac{1}{2}+8ϵ,l=0,a_1=\frac{1}{2}2ϵ,a=a_2=\frac{1}{2}+ϵ`$) and interpolation
$`\stackrel{~}{n}_\pm v_{X^{s,\frac{1}{2}+2ϵ}(I)}`$ $``$ $`c\stackrel{~}{n}_\pm _{X^{0,\frac{1}{2}+ϵ}(I)}v_{X^{s\frac{1}{2}+8ϵ,\frac{1}{2}+ϵ}(I)}`$
$``$ $`c\stackrel{~}{n}_\pm _{X^{0,\frac{1}{2}+ϵ}(I)}v_{X^{0,\frac{1}{2}+ϵ}(I)}^{\frac{1}{2s}}v_{X^{s,\frac{1}{2}+ϵ}(I)}^{1\frac{1}{2s}+}`$
$``$ $`cN^{1s}N^{\frac{1}{2}+}cN^{\frac{1}{2}s+}`$
$`m_\pm v_{X^{s,\frac{1}{2}+2ϵ}(I)}`$ $``$ $`cm_\pm _{X^{0,\frac{1}{2}+ϵ}(I)}v_{X^{s\frac{1}{2}+8ϵ,\frac{1}{2}+ϵ}(I)}`$
$``$ $`cm_\pm _{X^{0,\frac{1}{2}+ϵ}(I)}v_{X^{0,\frac{1}{2}+ϵ}(I)}^{\frac{1}{2s}}v_{X^{s,\frac{1}{2}+ϵ}(I)}^{1\frac{1}{2s}+}`$
$``$ $`cN^{\frac{1}{2}\frac{1}{4}s+}N^{\frac{1}{2}+}=cN^{1\frac{1}{4}s+}`$
$`m_\pm \stackrel{~}{u}_{X^{s,\frac{1}{2}+2ϵ}(I)}`$ $``$ $`cm_\pm _{X^{0,\frac{1}{2}+ϵ}(I)}\stackrel{~}{u}_{X^{s\frac{1}{2}+8ϵ,\frac{1}{2}+ϵ}(I)}`$
$``$ $`cm_\pm _{X^{0,\frac{1}{2}+ϵ}(I)}\stackrel{~}{u}_{X^{0,\frac{1}{2}+ϵ}(I)}^{\frac{3}{2}s}\stackrel{~}{u}_{X^{1,\frac{1}{2}+ϵ}(I)}^{s\frac{1}{2}+}`$
$``$ $`cN^{\frac{1}{2}\frac{1}{4}s}N^{(1s)(s\frac{1}{2})+}=cN^{s(s\frac{1}{2})1+\frac{3}{4}s+}`$
We have $`1>s>9/10`$, so the exponents of $`N`$ are negative. Thus with $`\gamma (s)>0`$:
$$F_{X^{s,\frac{1}{2}+2ϵ}(I)}cN^{\gamma (s)}$$
and therefore with $`c_02cu_{02}_{H^{s,2}(𝐑)}`$ we have:
$$S_0v_{X^{s,\frac{1}{2}+ϵ}(I)}cu_{02}_{H^{s,2}(𝐑)}+cF_{X^{s,\frac{1}{2}+2ϵ}(I)}\frac{c_0}{2}+cN^{\gamma (s)}c_0$$
if $`N`$ is sufficiently large.
Next we estimate $`S_0v_{X^{0,\frac{1}{2}+ϵ}(I)}`$. We use ,Lemma 4.3 (with $`k=0`$ , $`l=\frac{1}{2}`$ , $`c_1=\frac{1}{2}`$ , $`b=\frac{1}{2}+`$ , $`b_1=\frac{1}{2}+`$ and with $`k=l=0`$ , $`c_1=\frac{1}{4}+`$ , $`b=b_1=\frac{1}{2}+`$ ) and get:
$`S_0v_{X^{0,\frac{1}{2}+ϵ}(I)}`$
$``$ $`cu_{02}_{L^2(𝐑)}`$
$`+c(m_\pm v_{X^{0,\frac{1}{2}+2ϵ}(I)}|I|^ϵ+\stackrel{~}{n}_\pm v_{X^{0,\frac{1}{4}}(I)}|I|^{\frac{1}{4}}+m_\pm \stackrel{~}{u}_{X^{0,\frac{1}{2}+2ϵ}(I)}|I|^ϵ)`$
$``$ $`c_1N^s+c(m_\pm _{X^{\frac{1}{2},\frac{1}{2}+}(I)}v_{X^{0,\frac{1}{2}+}(I)}|I|^ϵ+\stackrel{~}{n}_\pm _{X^{0,\frac{1}{2}+}(I)}v_{X^{0,\frac{1}{2}+}(I)}|I|^{\frac{1}{4}}`$
$`+m_\pm _{X^{\frac{1}{2},\frac{1}{2}+}(I)}\stackrel{~}{u}_{X^{0,\frac{1}{2}+}(I)}|I|^ϵ)`$
$``$ $`c_1N^s+c(N^sN^s|I|^ϵ+N^{1s}N^s|I|^{\frac{1}{4}}+N^s|I|^ϵ)`$
$``$ $`c(N^s|I|^ϵ+N^{12s}|I|^{\frac{1}{4}})+c_1N^s`$
$``$ $`c_1N^s+cN^s|I|^ϵ+cN^{12s(1s)\frac{\delta }{4}+}`$
$``$ $`N^s({\displaystyle \frac{c_0}{2}}+cN^{4(1s)ϵ}+cN^{\frac{\delta }{4}+})`$
$``$ $`c_0N^s`$
from the definition of $`Sv`$,(82),(5) and (6), where $`c_02c_1`$ , $`ϵ>0`$ small and $`N`$ sufficiently large.
Next we treat $`S_\pm m_\pm _{X^{0,\frac{1}{2}+ϵ}(I)}`$. By , Lemma 4.4 (with $`l=0`$ , $`k=\frac{1}{4}`$ , $`c=\frac{1}{4}+`$ , $`b_1=\frac{1}{2}`$):
$`(\stackrel{~}{u}\overline{v})_x_{X^{0,\frac{1}{4}}(I)}+(|v|^2)_x_{X^{0,\frac{1}{4}}(I)}+(v\overline{\stackrel{~}{u}})_x_{X^{0,\frac{1}{4}}(I)}`$
$`c(\stackrel{~}{u}_{X^{\frac{1}{4},\frac{1}{2}}(I)}v_{X^{\frac{1}{4},\frac{1}{2}}(I)}+v_{X^{\frac{1}{4},\frac{1}{2}}(I)}^2)`$
$`c(\stackrel{~}{u}_{X^{0,\frac{1}{2}}(I)}^{\frac{3}{4}}\stackrel{~}{u}_{X^{1,\frac{1}{2}}(I)}^{\frac{1}{4}+}v_{X^{0,\frac{1}{2}+}(I)}^{1\frac{1}{4s}}v_{X^{s,\frac{1}{2}+}(I)}^{\frac{1}{4s}}+v_{X^{0,\frac{1}{2}+}(I)}^{2(1\frac{1}{4s})}v_{X^{s,\frac{1}{2}+}(I)}^{\frac{1}{2s}})`$
$`c(N^{\frac{1s}{4}+}N^{s(1\frac{1}{4s})}+N^{2s(1\frac{1}{4s})})`$
$`=c(N^{\frac{1}{2}\frac{5}{4}s+}+N^{\frac{1}{2}2s})`$
$`cN^{\frac{1}{2}\frac{5}{4}s+}`$
Thus
$`S_\pm m_\pm _{X^{0,\frac{1}{2}+}(I)}={\displaystyle _0^t}e^{i(ts)A^{1/2}}G(s)𝑑s_{X^{0,\frac{1}{2}+}(I)}cG_{X^{0,\frac{1}{4}}(I)}|I|^{\frac{1}{4}}`$
$`cN^{\frac{1}{2}\frac{5}{4}s+}|I|^{\frac{1}{4}}=cN^{\frac{1}{2}\frac{5}{4}s(1s)\frac{\delta }{4}+}=cN^{\frac{1}{2}\frac{1}{4}s\frac{\delta }{4}+}c_0N^{\frac{1}{2}\frac{1}{4}s\frac{\delta }{4}+}`$ (88)
if $`N`$ is chosen sufficiently large.
Finally we treat $`S_\pm m_\pm _{X^{\frac{1}{2},\frac{1}{2}+ϵ}(I)}`$. By , Lemma 4.4 (with $`l=\frac{1}{2}`$ , $`k=0`$ , $`b_1=\frac{1}{2}+`$ , $`c=\frac{1}{4}+`$):
$`G_{X^{\frac{1}{2},\frac{1}{4}}(I)}`$ $``$ $`(\stackrel{~}{u}\overline{v})_x_{X^{\frac{1}{2},\frac{1}{4}}(I)}+(|v|^2)_x_{X^{\frac{1}{2},\frac{1}{4}}(I)}+(v\overline{\stackrel{~}{u}})_x_{X^{\frac{1}{2},\frac{1}{4}}(I)}`$
$``$ $`c(\stackrel{~}{u}_{X^{0,\frac{1}{2}+}(I)}v_{X^{0,\frac{1}{2}+}(I)}+v_{X^{0,\frac{1}{2}+}(I)}^2)`$
$``$ $`c(N^s+N^{2s})`$
$``$ $`cN^s`$
Thus
$`S_\pm m_\pm _{X^{\frac{1}{2},\frac{1}{2}+ϵ}(I)}={\displaystyle _0^t}e^{i(ts)A^{1/2}}G(s)𝑑s_{X^{\frac{1}{2},\frac{1}{2}+ϵ}(I)}`$ (89)
$`cG_{X^{\frac{1}{2},\frac{1}{4}}(I)}|I|^{\frac{1}{4}}cN^s|I|^{\frac{1}{4}}c_0N^s`$
for $`N`$ sufficiently large.
Summarizing we have shown that $`S`$ maps $`Z`$ into itself.
The contraction property uses exactly the same type of estimates and its proof is therefore omitted. Thus the proposition is proved.
The next estimates show that the nonlinear part $`w`$ of (79), defined by (83), behaves better than the linear part.
We use (9) in connection with Prop. 1.2 (with $`s=1,l=0,k=\frac{1}{2}+,a=a_2=\frac{1}{2}`$). This gives
$`w_{L^{\mathrm{}}(I,H^{1,2}(𝐑))}`$ (90)
$`c\left((m_++m_{})v_{Y^1(I)}+(\stackrel{~}{n}_++\stackrel{~}{n}_{})v_{Y^1(I)}+(m_++m_{})\stackrel{~}{u}_{Y^1(I)}\right)`$
$`c(m_++m_{}_{X^{0,\frac{1}{2}}(I)}v_{X^{\frac{1}{2}+,\frac{1}{2}}(I)}+\stackrel{~}{n}_++\stackrel{~}{n}_{}_{X^{0,\frac{1}{2}}(I)}v_{X^{\frac{1}{2}+,\frac{1}{2}}(I)}`$
$`+m_++m_{}_{X^{0,\frac{1}{2}}(I)}\stackrel{~}{u}_{X^{\frac{1}{2}+,\frac{1}{2}}(I)})`$
$`c(m_++m_{}_{X^{0,\frac{1}{2}}(I)}v_{X^{0,\frac{1}{2}}(I)}^{1\frac{1}{2s}}v_{X^{s,\frac{1}{2}}(I)}^{\frac{1}{2s}+}`$
$`+\stackrel{~}{n}_++\stackrel{~}{n}_{}_{X^{0,\frac{1}{2}}(I)}v_{X^{0,\frac{1}{2}}(I)}^{1\frac{1}{2s}}v_{X^{s,\frac{1}{2}}(I)}^{\frac{1}{2s}+}`$
$`+m_++m_{}_{X^{0,\frac{1}{2}}(I)}\stackrel{~}{u}_{X^{0,\frac{1}{2}}(I)}^{\frac{1}{2}}\stackrel{~}{u}_{X^{1,\frac{1}{2}}(I)}^{\frac{1}{2}+})`$
$`c(N^{\frac{1}{2}\frac{1}{4}s}N^{s(1\frac{1}{2s})+}+N^{1s}N^{s(1\frac{1}{2s})+}+N^{\frac{1}{2}\frac{1}{4}s+}N^{\frac{1s}{2}+})`$
$`=c(N^{\frac{5}{4}s}+N^{\frac{3}{2}2s+}+N^{\frac{3}{4}s})cN^{\frac{3}{2}2s+}`$
because $`\frac{3}{2}2s>\frac{3}{4}s`$ , by (88),(84),(85),(65),(66),(68), using $`|I|=N^{4(1s)}`$. The next step is to estimate $`w_{L^{\mathrm{}}(I,L^2(𝐑))}`$.
We use , Lemma 4.3 (with $`k=l=0`$ , $`b=b_1=\frac{1}{2}+`$ , $`c_1=\frac{1}{4}+`$):
$`w_{L^{\mathrm{}}(I,L^2(𝐑))}cw_{X^{0,\frac{1}{2}+}(I)}`$ (91)
$`c((m_++m_{})v_{X^{0,\frac{1}{4}}(I)}+(\stackrel{~}{n}_++\stackrel{~}{n}_{})v_{X^{0,\frac{1}{4}}(I)}`$
$`+(m_++m_{})\stackrel{~}{u}_{X^{0,\frac{1}{4}}(I)})|I|^{\frac{1}{4}}`$
$`c(m_++m_{}_{X^{0,\frac{1}{2}+}(I)}v_{X^{0,\frac{1}{2}+}(I)}+\stackrel{~}{n}_++\stackrel{~}{n}_{}_{X^{0,\frac{1}{2}+}(I)}v_{X^{0,\frac{1}{2}+}(I)}`$
$`+m_++m_{}_{X^{0,\frac{1}{2}+}(I)}\stackrel{~}{u}_{X^{0,\frac{1}{2}+}(I)})|I|^{\frac{1}{4}}`$
$`c(N^{\frac{1}{2}\frac{1}{4}s}N^s+N^{1s}N^s+N^{\frac{1}{2}\frac{1}{4}s})|I|^{\frac{1}{4}}`$
$`c(N^{12s}+N^{\frac{1}{2}\frac{1}{4}s})N^{(1s)}`$
$`cN^{\frac{3}{2}+\frac{3}{4}s+}`$
by (89),(84),(72),(66), because $`12s<\frac{1}{2}\frac{1}{4}ss>\frac{6}{7}`$
We also repeat the corresponding estimate for $`m_\pm =Sm_\pm =z_\pm `$ from (88):
$$z_\pm _{X^{0,\frac{1}{2}+}(I)}=m_\pm _{X^{0,\frac{1}{2}+}(I)}cN^{\frac{1}{2}\frac{1}{4}s}$$
(92)
## 5 The iteration process
In the preceding sections we have constructed a solution $`(u,n)`$ of our original problem (1),(2) with data (3) $`(u_0,n_0,n_1)`$ in the time interval $`I=[0,|I|]`$ with $`|I|=N^{4(1s)}`$. Namely, if we define $`u:=v+\stackrel{~}{u}`$ , $`n_\pm :=m_\pm +\stackrel{~}{n}_\pm `$ we easily see by (75),(76) that $`(u,n_+,n_{})`$ satisfies the system (47). Moreover the initial conditions $`u(0)=u_0`$ , $`n_\pm (0)=n_{0\pm }`$ are satisfied. This initial value problem is equivalent to the original system (1),(2) by (45). The initial data are transformed via (45),(46) by $`n_{0\pm }=\frac{1}{2}(n_0+n_1)`$ , $`2iA^{1/2}n_1=n_{0+}n_0`$ or conversely by $`n_{0\pm }=n_0\pm iA^{1/2}n_1`$.
In order to continue the solution of (47),(48) we take as new initial data for our system (47) the triple $`(\stackrel{~}{u}(|I|)+w(|I|),\stackrel{~}{n}_+(|I|)+z_+(|I|),\stackrel{~}{n}_{}(|I|)+z_{}(|I|))`$ instead of $`(u_{01},n_{0+},n_0)`$. When we have shown that this problem has a solution $`(\stackrel{~}{\stackrel{~}{u}},\stackrel{~}{\stackrel{~}{n}}_+,\stackrel{~}{\stackrel{~}{n}}_{})`$ in the time interval $`[|I|,2|I|]`$ of equal length $`|I|`$ we insert this solution into the system (75),(76) in place of $`(\stackrel{~}{u},\stackrel{~}{n}_+,\stackrel{~}{n}_{})`$ and solve this problem with data $`(e^{i|I|_x^2}u_{02},0,0)`$ in $`[|I|,2|I|]`$. Adding up the solutions we get a solution of the original problem in $`[|I|,2|I|]`$ as before. This defines an iteration process. At each step we have to ensure the same bounds on the initial data which were used in the first step. The replacement of $`u_{02}`$ by$`e^{i|I|_x^2}u_{02}`$ is harmless, because $`e^{i|I|_x^2}`$ is unitary in $`\dot{H}^s(𝐑)`$. These bounds are controlled by the energy and the $`L^2`$-conservation law (cf. (50),(51)). Thus we have to estimate these quantities independently of the iteration step. This is easy for $`L^2`$-conservation, the increment when replacing $`u_{01}`$ by $`u_{01}(|I|)+w(|I|)`$ using $`L^2`$-conservation is given by
$`\left|\stackrel{~}{u}(|I|)+w(|I|)_{L^2(𝐑)}u_{01}_{L^2(𝐑)}\right|=\left|\stackrel{~}{u}(|I|)+w(|I|)_{L^2(𝐑)}\stackrel{~}{u}(|I|)_{L^2(𝐑)}\right|`$
$`w(|I|)_{L^2(𝐑)}c_2N^{\frac{3}{2}+\frac{3}{4}s+}`$
by (91), where $`c_2=c_2(\overline{c},M)`$.
The number of iteration steps in order to reach the given time $`T`$ is $`\frac{T}{|I|}=TN^{4(1s)+}`$. This means that we have to ensure in order to get uniform control over the $`L^2`$-norm of $`\stackrel{~}{u},\stackrel{~}{\stackrel{~}{u}},\mathrm{}`$:
$$c_2TN^{4(1s)+}+N^{\frac{3}{2}+\frac{3}{4}s+}<M$$
where $`c_2=c_2(2\overline{c},2M)`$ (remark that initially the $`L^2`$-norm of $`\stackrel{~}{u}`$ was also bounded by $`M`$).
This is fulfilled for $`N`$ sufficiently large if $`4(1s)\frac{3}{2}+\frac{3}{4}s<0s>\frac{10}{13}`$ which is fulfilled.
The increment of the energy is given by
$`\left|E(\stackrel{~}{u}(|I|)+w(|I|),n(|I|)+m(|I|),n_t(|I|)+m_t(|I|))E(u_{01},n_0,n_1)\right|`$
$`=\left|E(\stackrel{~}{u}(|I|)+w(|I|),n(|I|)+m(|I|),n_t(|I|)+m_t(|I|))E(\stackrel{~}{u}(|I|),n(|I|),n_t(|I|))\right|`$
$`2(\stackrel{~}{u}_x(|I|)+w_x(|I|))w_x(|I|)+(n(|I|)+m(|I|))w(|I|)`$
$`+(A^{1/2}n_t(|I|)+A^{1/2}m_t(|I|))A^{1/2}m_t(|I|)`$
$`+{\displaystyle _{\mathrm{}}^{\mathrm{}}}|m(|I|)||\stackrel{~}{u}(|I|)+w(|I|)|^2𝑑x+{\displaystyle _{\mathrm{}}^{\mathrm{}}}|n(|I|)|||\stackrel{~}{u}(|I|)+w(|I|)|^2|\stackrel{~}{u}(|I|)|^2|𝑑x`$
Using (54),(90) the first term is bounded by
$$c(N^{1s}+N^{\frac{3}{2}2s+})N^{\frac{3}{2}2s+}cN^{1s}N^{\frac{3}{2}2s+}$$
the second and third one using (54),(92) by
$$c(N^{1s}+N^{\frac{1}{2}\frac{1}{4}s})N^{\frac{1}{2}\frac{1}{4}s}cN^{1s}N^{\frac{1}{2}\frac{1}{4}s}$$
The fourth term is estimated using Gagliardo-Nirenberg and (54),(90),(91),(92):
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}|m(|I|)||\stackrel{~}{u}(|I|)+w(|I|)|^2𝑑x2{\displaystyle _{\mathrm{}}^{\mathrm{}}}|m(|I|)|(|\stackrel{~}{u}(|I|)|^2+|w(|I|)|^2)𝑑x`$
$`2m(|I|)_{L^2(𝐑)}(\stackrel{~}{u}(|I|)_{L^2(𝐑)}\stackrel{~}{u}(|I|)_{L^{\mathrm{}}(𝐑)}+w(|I|)_{L^2(𝐑)}w(|I|)_{L^{\mathrm{}}(𝐑)})`$
$`cm(|I|)_{L^2(𝐑)}(\stackrel{~}{u}(|I|)_{L^2(𝐑)}^{\frac{3}{2}}\stackrel{~}{u}_x(|I|)_{L^2(𝐑)}^{\frac{1}{2}}+w(|I|)_{L^2(𝐑)}^{\frac{3}{2}}w_x(|I|)_{L^2(𝐑)}^{\frac{1}{2}})`$
$`cN^{\frac{1}{2}\frac{1}{4}s}(N^{\frac{1s}{2}}+N^{\frac{3}{2}(\frac{3}{2}+\frac{3}{4}s)+}N^{\frac{1}{2}(\frac{3}{2}2s)+})`$
$`cN^{\frac{1}{2}\frac{1}{4}s}N^{\frac{1s}{2}}=cN^{\frac{3}{4}s}`$
The fifth term is similarly estimated as follows:
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}|n(|I|)|||\stackrel{~}{u}(|I|)+w(|I|)|^2|\stackrel{~}{u}(|I|)|^2|𝑑x`$
$`2n(|I|)_{L^2}(\stackrel{~}{u}(|I|)_{L^2}+w(|I|)_{L^2})w(|I|)_L^{\mathrm{}}`$
$`cn(|I|)_{L^2}(\stackrel{~}{u}(|I|)_{L^2}+w(|I|)_{L^2})w(|I|)_{L^2}^{\frac{1}{2}}w_x(|I|)_{L^2}^{\frac{1}{2}}`$
$`cN^{1s}(1+N^{\frac{3}{2}+\frac{3}{4}s+})N^{\frac{1}{2}(\frac{3}{2}+\frac{3}{4}s)+}N^{\frac{1}{2}(\frac{3}{2}2s)+}`$
$`cN^{1s}N^{\frac{5}{8}s+}`$
It is easy to see that the decisive bound is the one for the first term. Thus the increment of the energy is bounded by
$$c_3N^{1s}N^{\frac{3}{2}2s+}$$
where $`c_3=c_3(\overline{c},M)`$.
Thus the condition which ensures uniform control of the energy of $`(\stackrel{~}{u},\stackrel{~}{n}),(\stackrel{~}{\stackrel{~}{u}},\stackrel{~}{\stackrel{~}{n}}),\mathrm{}`$ is the following:
$$c_3TN^{4(1s)+}N^{1s}N^{\frac{3}{2}2s+}<\overline{c}N^{2(1s)}$$
(93)
where $`c_3=c_3(2\overline{c},2M)`$ (recall that by (53) initially the energy is bounded by $`\overline{c}N^{2(1s)}`$ ).
This is fulfilled for $`N`$ sufficiently large provided
$$4(1s)+(1s)+\frac{3}{2}2s<2(1s)s>\frac{9}{10}$$
So, here is the point where the decisive bound on $`s`$ appears. The uniform control of the energy implies by (50),(51) uniform control of the $`L^2`$-norm of $`(\stackrel{~}{u}_x,\stackrel{~}{n},A^{1/2}\stackrel{~}{n}_t),(\stackrel{~}{\stackrel{~}{u}}_x,\stackrel{~}{\stackrel{~}{n}},A^{1/2}\stackrel{~}{\stackrel{~}{n}}_t),\mathrm{}`$
###### Theorem 5.1
Let $`1>s>9/10`$. The Zakharov system (1),(2),(3) with data $`(u_0,n_0,n_1)H^s(𝐑)\times L^2(𝐑)\times \dot{H}^1(𝐑)`$ is globally well-posed. More precisely for any $`T>0`$ there exists a unique solution
$$(u,n,n_t)X^{s,\frac{1}{2}+ϵ_1}[0,T]\times X^{0,\frac{1}{2}+ϵ_2}[0,T]\times \dot{X}^{1,\frac{1}{2}+ϵ_2}[0,T]$$
(94)
for $`ϵ_1,ϵ_2>0`$ small enough. This solution satisfies
$$(u,n,n_t)C^0([0,T],H^s(𝐑)\times L^2(𝐑)\times \dot{H}^1(𝐑))$$
(95)
Proof: On any of the intervals $`I`$ of the preceding considerations we have by (72),(68),(66) (+ interpolation) $`\stackrel{~}{u}X^{s,\frac{1}{2}+ϵ_1}(I),`$ $`\stackrel{~}{n}_\pm X^{0,\frac{1}{2}+ϵ_2}(I)`$ and by (85),(87) $`vX^{s,\frac{1}{2}+ϵ_1}(I),`$ $`m_\pm X^{0,\frac{1}{2}+ϵ_2}(I)`$. This gives (94) by (45),(46). Uniqueness in this class was proven in already. (95) follows immediately from (94).
It is not difficult to give bounds on the growth of the solutions now. It is elementary to show that the most restrictive bound on $`N`$ comes from condition (93) in the whole range $`9/10<s<1`$, namely
$$N>cT^{\frac{1}{5s\frac{9}{2}}+}$$
(96)
According to the construction of the solution above we have the following structure
$$u(t)=\stackrel{~}{u}(t)+e^{it_x^2}u_{02}+w(t)=e^{it_x^2}u_0+r(t)$$
where
$$r(t)=\stackrel{~}{u}(t)e^{it_x^2}u_{01}+w(t)$$
on $`I`$ first, but then also on $`[0,T]`$.
We have shown that
$$r(t)_{H^{1,2}(𝐑)}cN^{1s}$$
(97)
(remark that $`e^{it_x^2}u_{01}_{H^{1,2}(𝐑)}cN^{1s}`$). Choosing $`N`$ according to (96) gives the following bound for $`0tT`$:
$$r(t)_{H^{1,2}(𝐑)}cT^{\frac{1s}{5s\frac{9}{2}}+}$$
Similarly
$$n_\pm (t)_{L^2(𝐑)}cN^{1s}cT^{\frac{1s}{5s\frac{9}{2}}+}$$
Thus we have shown
###### Theorem 5.2
The solution of the preceding theorem fulfills for $`t0`$:
$$u(t)=e^{it_x^2}u_0+r(t)$$
with
$$r(t)_{H^{1,2}(𝐑)}c\left(1+|t|^{\frac{1s}{5s\frac{9}{2}}+}\right)$$
and
$$n(t)_{L^2(𝐑)}+n_t(t)_{\dot{H}^{1,2}(𝐑)}c\left(1+|t|^{\frac{1s}{5s\frac{9}{2}}+}\right)$$
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# Spectral function of one hole in several one-dimensional spin arrangements
## I Introduction
The understanding of spin-charge separation is a central point in the physics of low-dimensional electronic systems. That issue is most clearly seen in one dimension where the exact solution of the 1D Hubbard model reveals that the low-energy physics is dominated by decoupled, collective charge and spin excitations (also called holon and spinon, respectively). The idea that the spin and charge degrees of freedom separate has also been proposed to explain the properties of 2D cuprate superconductors. However, even in one dimension, the Bethe ansatz solution did not yet give a complete answer for the spectral function $`A(k,\omega )`$. Only for $`U\mathrm{}`$ an exact expression for $`A(k,\omega )`$ was found due to the factorization property of the wave function according to Ogata and Shiba. The results for the insulating half-filled case could also be generalized to other filling factors.
The spin-charge separation was observed in ARPES measurements of the one dimensional, dielectric cuprate SrCuO<sub>2</sub>. The spinon and holon branch of the spectral function were seen, in contrast to the analogous experiment of one hole in the CuO<sub>2</sub> plane where spin and charge are coupled and the spin polaron quasiparticle has a dispersion proportional to $`J`$ as proposed theoretically in Refs. . The ARPES spectra in SrCuO<sub>2</sub> were analyzed using the pure $`t`$-$`J`$ model. On the other hand, many 1D compounds, like for instance CuGeO<sub>3</sub>, are characterized by frustration in the magnetic subsystem which may lead to a gap in the spin excitation spectrum. For the special frustration $`J^{}/J=0.5`$ and in the limit $`J0`$ an analytic expression for $`A(k,\omega )`$ was derived recently, under the assumption that the wave function factorizes. Besides the frustration, also temperature effects are important as it was observed in ARPES measurements on Na<sub>0.96</sub>V<sub>2</sub>O<sub>5</sub>. So, there is a clear need to study the spectral function systematically under the influence of frustration and temperature and to derive analytic expressions.
The present work focuses mainly on the effect of frustration and temperature on the spectral function in the insulating case. For that, we rederive first the exact solution of Sorella and Parola in a straightforward way using Green’s function technique. That is possible due to our finding of a set of eigenoperators of the Liouvillian (Sec. III) in the strong coupling limit $`J0`$. As a consequence, our derivation is applicable for any magnetic state and any temperature in that limit. Especially, one can show that the result which was derived in Ref. does not depend on the assumption that the wave function factorizes. We present analytic expressions for the spectral function of one hole in several magnetic states: (i) the ground-state of the antiferromagnetic Heisenberg model, (ii) the Majumdar-Ghosh wave function at the special frustration $`J^{}/J=0.5`$, and (iii) the ideal paramagnetic state at temperatures much larger than the exchange energy $`k_BTJ`$.
For large finite coupling we compare two methods to account for corrections $`J(t^2/U)`$, namely the projection method and a variational ansatz using the set of eigenoperators of the $`t`$-term. We show that the former method yields a reasonable description of the spinon dispersion in the pure $`t`$-$`J`$ model (Sec. IV) and an approximate result for the spectral function of the Majumdar-Ghosh model. For $`J^{}=J/2`$ it misses the bound state below the continuum which is obtained by the more accurate variational method (Sec. VB). The bound state has a finite spectral weight but a very small separation from the continuum. Both methods show in the Majumdar-Ghosh case that the low energy region for momenta $`k`$ between $`\pi /2`$ and $`\pi `$ (lattice constant $`a=1`$) will be filled with states, that the spinon dispersion (i.e. that collective excitation corresponding to the lower edge of the continuum) becomes symmetric around $`\pi /2`$, and they indicate an overdamped holon branch. The damping of the holon branch is extremely large for very high temperatures (Sec. VI).
Before presenting our results let us shortly discuss the different understandings of the term “spin-charge separation” as it can be met in the literature. The naive picture means that the low energy effective Hamiltonian may be written as
$$\widehat{H}=\widehat{H}_h+\widehat{H}_s,[\widehat{H}_h,\widehat{H}_s]=0,$$
(1)
and the electron operator is the product
$$c_{i\sigma }=s_{i\sigma }h_i^{},$$
(2)
where spinon $`s`$ and holon $`h`$ can be basically regarded as free particles. Then the normalized (i.e. $`A(k,\omega )𝑑\omega =1`$) spectral function is
$$A(k,\omega )=\frac{1}{L}\underset{Q}{}2f(Q)\delta \left[\omega ϵ_h(kQ)ϵ_s(Q)\right],$$
(3)
where $`f(Q)=\theta (\frac{\pi }{2}\left|Q\right|)`$ is the Fermi distribution function of spinons, $`\theta (x)`$ is the Heaviside step function, $`L`$ is the number of sites, $`ϵ_h,ϵ_s`$ being holon and spinon energies, respectively.
However, the naive understanding is not that one which is realized in 1D electron systems. There, it was found that the eigenstates factorize in the limit $`U\mathrm{}`$ in the form
$$\psi (x_1,\mathrm{},x_N,y_1\mathrm{}y_M)=\psi _{SF}(x_1,\mathrm{},x_N)\varphi _H(y_1\mathrm{}y_M),$$
(4)
where $`x_1,\mathrm{},x_N`$ are the spatial coordinates of the $`N`$ electrons on a $`L`$-site ring, and the $`y_1\mathrm{}y_M`$ ’coordinates’ label the position of the spin-up electrons on the squeezed Heisenberg ring, i.e. on the $`N`$ occupied sites. The $`\psi _{SF}`$ is a spinless fermion state, and $`\varphi _H`$ is an eigenstate of an $`N`$-site Heisenberg Hamiltonian with periodic boundary conditions. The product form of equation (4) should not be interpreted as a trivial decoupling between charge and spin. In fact, the momentum of the spin wave function imposes a twisted boundary condition on the spinless fermion wave function. As a result, the Fermi distribution function $`f(Q)`$ in (3) will be replaced by a function $`Z(Q)`$ that is the expectation value of a chain of spin operators that has to be determined from the pure spin system (for the details see Sec. III). The singularity of $`Z(Q)`$ produces additional peaks in $`A(k,\omega )`$. This correct answer for the spectral function may be understood as a manifestation of the phase string effect. It means that spinon and holon interact with each other via a nonlocal phase-string. Instead of (2) we should write
$$c_{i\sigma }=s_{i\sigma }h_i^{}\mathrm{exp}\left[\frac{\pi }{2}\underset{l>i}{}h_l^{}h_l+\frac{\pi }{2}\underset{l>i}{}(s_{l\sigma }^{}s_{l\sigma }1)\right].$$
It should be noted that the phenomenon of spin-charge separation is not restricted to the limit $`U\mathrm{}`$ in the 1D Hubbard model. At any finite $`U`$ the spin and charge fluctuations propagate with different velocities. That means that after some time the spin and charge degrees of freedom will be separated in space. But there is no analytic solution for the spectral function of the Hubbard model at arbitrary values of $`U`$ and also the present calculation treats terms of order $`(t^2/U)`$ as a perturbation. In that sense we will understand here spin-charge separation as a manifestation of the factorization property (4) in the spectral density. Sharp maxima in the continuum correspond to collective excitations whereas possible bound states indicate special eigenfunctions with a strong coupling between spin and charge. Another possible effect of additional terms in the Hamiltonian is the broadening (i.e. the damping) of the collective excitations.
## II Model and spectral density
To describe the low energy physics of compounds with a 1D electronic structure it is sufficient in most cases to take into account only that band which is closest to the Fermi energy (see for instance Ref. ). Treating the on-site Coulomb interaction explicitly, one obtains the well known 1D Hubbard model. In the present calculation we restrict ourselves to the strong coupling limit $`Ut`$ where we may project out the subspace of doubly occupied sites, and for the lower Hubbard band we obtain the effective Hamiltonian
$$\widehat{H}=\widehat{t}+\widehat{J}+\widehat{t}_3,$$
(5)
where
$$\widehat{t}=t\underset{i,g,\alpha }{}X_i^{\alpha 0}X_{i+g}^{0\alpha },$$
(6)
$$\widehat{J}=\frac{J}{2}\underset{i,\alpha ,\beta }{}X_i^{\alpha \beta }X_{i+1}^{\beta \alpha },$$
(7)
$$\widehat{t}_3=t_3\underset{i,g,\alpha ,\beta }{}X_i^{\alpha 0}X_{i+g}^{\beta \alpha }X_{i+2g}^{0\beta },$$
(8)
and $`\alpha ,\beta =`$; $`g`$ are the nearest neighbors $`g=\pm 1`$. The Hamiltonian is valid near half filling ($`X_i^{++}+X_i^{}=1`$). The parameters $`J`$ and $`t_3`$ are connected with the original values of the Hubbard model by
$$J=4t_3=4t^2/U,$$
(9)
but the Hamiltonian (5) is more general, if we relax the condition (9). It may be derived directly from the more realistic three-band Hubbard model. Then, the $`t_3`$-term often becomes negligible and one obtains the $`t`$-$`J`$ model. The Hamiltonian (5) is written in terms of Hubbard projection operators that act in the subspace of on-site states
$$X_i^{\alpha \beta }|\alpha ,i\beta ,i|,\alpha ,\beta =0,,,2.$$
(10)
They are related with bare fermionic and spin operators through
$$X_i^{\sigma 0}=c_{i,\sigma }^{}(1n_{i,\sigma }),X_i^{\sigma 2}=\sigma c_{i,\sigma }n_{i,\sigma },$$
(11)
$$X_i^+=S_i^+=c_{i,}^{}c_{i,},X_i^{\sigma \sigma }=\frac{1}{2}+\frac{\sigma }{2}(c_{i,}^{}c_{i,}c_{i,}^{}c_{i,})=\frac{1}{2}+\sigma S_i^z,$$
(12)
with $`\sigma =\pm 1`$. Other relations are easy to obtain with the use of the main property of Hubbard operator algebra
$$X_i^{\alpha \beta }X_i^{\gamma \lambda }=\delta _{\beta \gamma }X_i^{\alpha \lambda },$$
(13)
that follows immediately from the definition (10). The commutation relations for operators on different sites are fermionic for operators that change the number of particles by odd integers, like (11), and bosonic for others. In the presence of frustration in the magnetic system, which is discussed for instance for CuGeO<sub>3</sub>, the $`t`$-$`J`$ Hamiltonian may be generalized by inclusion of the $`J^{}`$-term
$$\widehat{J}^{}=\frac{J^{}}{2}\underset{i,\alpha ,\beta }{}X_i^{\alpha \beta }X_{i+2}^{\beta \alpha }.$$
(14)
Our aim is to calculate the one-particle two-time retarded Green’s function $`G(k,\omega )`$ and the spectral density of one hole in the magnetic state
$$A(k,\omega )=\frac{1}{\pi }\mathrm{Im}G(k,\omega +i0^+),$$
(15)
that is roughly proportional to the ARPES signal intensity. We define
$$2\pi \delta (kk^{})G(k,\omega )=X_k^{\sigma 0}|X_k^{}^{0\sigma }i_t^{}^{\mathrm{}}𝑑te^{i\omega (tt^{})}\{X_k^{\sigma 0}(t),X_k^{}^{0\sigma }(t^{})\},$$
(16)
where
$$X_k^{\sigma 0}=\sqrt{2}\underset{m=\mathrm{}}{\overset{+\mathrm{}}{}}\mathrm{e}^{ikm}X_m^{\sigma 0},$$
$$\{X_k^{\sigma 0},X_k^{}^{0\sigma }\}=2\pi \delta (kk^{}),$$
and where $`\{\mathrm{},\mathrm{}\}`$ means the anticommutator. The expectation value denotes the thermal average over a grand canonical ensemble:
$$\mathrm{}=Q^1\text{Sp}[\mathrm{e}^{\beta (\widehat{H}\mu \widehat{N})}\mathrm{}],Q=\text{Sp}\mathrm{e}^{\beta (\widehat{H}\mu \widehat{N})}.$$
(17)
Here Sp implies taking the trace of an operator, $`\widehat{N}`$ is the particle number operator, $`\beta =(kT)^1`$ is an inverse temperature, and $`\mu `$ represents the chemical potential. The time dependence of the operator $`B(t)`$ is given by $`B(t)=\mathrm{e}^{it(\widehat{H}\mu \widehat{N})}B\mathrm{e}^{it(\widehat{H}\mu \widehat{N})}`$.
## III Eigenoperator and holon dispersion
Let us consider first the limit $`U\mathrm{}`$ in the Hubbard model or $`J,t_30`$ in (5). Then only the $`t`$-term
$$\widehat{t}=t\underset{i,g,\alpha }{}X_i^{\alpha 0}X_{i+g}^{0\alpha },$$
is nonzero. Note that it is a true many-body Hamiltonian due to the constraint of no double occupancy, as we see from Eq. (11). We introduce the set of operators
$$v_{m,r}=\underset{\alpha _1,\mathrm{},\alpha _r}{}X_m^{\sigma \alpha _1}X_{m+g}^{\alpha _1\alpha _2}\mathrm{}X_{m+rg}^{\alpha _{r1}\alpha _r}X_{m+r}^{\alpha _r0},g=\mathrm{sign}(r),$$
(18)
for which
$$[v_{m,r},\widehat{t}]=t\left(v_{m,r1}+v_{m,r+1}\right)$$
(19)
holds at half filling. Any operator (18) can be considered as a string operator of a certain length consisting of a hole and an attached string of spin flips. Such a string can be produced by creating a hole in the Néel state and applying several times the kinetic energy (6) which creates misaligned spins. Similar string operators were used to describe the spin polaron quasiparticle in the 2D case. We make double Fourier transform
$$v_{k,q}=\sqrt{2}\underset{m,r=\mathrm{}}{\overset{+\mathrm{}}{}}\mathrm{e}^{ikmiqr}v_{m,r},$$
(20)
and see that
$$[v_{k,q},\widehat{t}]=2tv_{k,q}\mathrm{cos}q.$$
(21)
The interpretation in terms of the string operator (18) is quite easy. We see that only the right end of the operator (18) was influenced by the $`t`$-term. Therefore, we may identify the right end of $`v_{m,r}`$ with the holon excitation. In the next Section, it will become clear that the left end of $`v_{m,r}`$ may be connected with the spinon.
The operators $`v_{k,q}`$ (20) are eigenoperators of the Liouvillian $``$ of the problem, where $`\widehat{A}[\widehat{H},\widehat{A}]`$. Note that it is one of the rarest, if not the unique case in many-body physics that the explicit form for a set of eigenoperators can be given. From (21) we see that the equation of motion for the corresponding string operator Green’s function closes and it has a simple pole form
$$v_{k,q}|v_{k^{},q^{}}^{}=\frac{\{v_{k,q},v_{k^{},q^{}}^{}\}}{\omega 2t\mathrm{cos}q},\{v_{k,q},v_{k^{},q^{}}^{}\}=8\pi ^2\delta (kk^{})\delta (qq^{})Z(kq+\pi ),$$
(22)
where the spectral weight is the expectation value
$$Z(q+\pi )=\frac{1}{2}\underset{r=\mathrm{}}{\overset{+\mathrm{}}{}}\mathrm{e}^{iqr}\mathrm{\Omega }_r,$$
(23)
of a chain of $`X`$-operators
$`\mathrm{\Omega }_r`$ $`=`$ $`{\displaystyle \underset{\alpha _1,\mathrm{},\alpha _r,\sigma }{}}X_m^{\sigma \alpha _1}X_{m+g}^{\alpha _1\alpha _2}\mathrm{}X_{m+rg}^{\alpha _{r1}\alpha _r}X_{m+r}^{\alpha _r\sigma }`$ (24)
$`=`$ $`(2𝐒_m𝐒_{m+g}+{\displaystyle \frac{1}{2}})(2𝐒_{m+g}𝐒_{m+2g}+{\displaystyle \frac{1}{2}})\mathrm{}(2𝐒_{m+rg}𝐒_{m+r}+{\displaystyle \frac{1}{2}}),`$ (25)
to be calculated for the pure spin-system without any hole. The expectation value $`\mathrm{\Omega }_r`$ of (25) cannot depend on the starting point $`m`$ due to the translational symmetry of the problem. The operator (25) was introduced in Ref. and explicit values on a 26-site Heisenberg ring were given for $`T=0`$, when $`\mathrm{}`$ becomes the average over the ground-state. Asymptotically, the following behavior was found
$$\mathrm{\Omega }_l\frac{1}{\sqrt{l}}\text{Re}\left[A\mathrm{e}^{i\pi l/2}\right]$$
(26)
which leads to a square root singularity of $`Z(Q)`$. Using additionally the exact values $`\mathrm{\Omega }_0=1`$ and $`\mathrm{\Omega }_1=12\mathrm{ln}2`$ the following formula may be derived
$$Z(Q)=\left(0.393+0.835/\sqrt{\mathrm{cos}Q}\right)\theta (\frac{\pi }{2}\left|Q\right|).$$
(27)
For the hole Green’s function (16) we have
$$X_k^{\sigma 0}|X_k^{}^{0\sigma }=_\pi ^{+\pi }\frac{dq}{2\pi }_\pi ^{+\pi }\frac{dq^{}}{2\pi }v_{k,q}|v_{k^{},q^{}}^{}=2\pi \delta (kk^{})_\pi ^{+\pi }\frac{dq}{\pi }\frac{Z(kq+\pi )}{\omega 2t\mathrm{cos}q},$$
(28)
and the spectral density is obtained in the way
$$A(k,\omega )=_\pi ^{+\pi }\frac{dQ}{\pi }Z(Q)\delta \left[\omega +2t\mathrm{cos}(kQ)\right].$$
(29)
That gives the exact answer in the strong coupling limit $`(J0)`$ where only the holon dispersion $`ϵ_h(q)=2t\mathrm{cos}q`$ is important. But the ratio $`J^{}/J`$ may be arbitrary and (29) is not only exact in the Heisenberg case with $`Z(Q)`$ from (27). Instead, from our derivation follows its validity for arbitrary magnetic states and it is not restricted to zero temperature. Then, however, $`Z(Q)`$ is different. In the following we will give exact results for i) the Majumdar-Ghosh wave function at the special frustration $`J^{}/J=0.5`$, and ii) the ideal paramagnetic case at $`k_BTJ`$. Two cases are quite trivial, namely the saturated ferromagnetic case and the classical Néel state. The former one leads to $`Z(Q+\pi )\delta (Q)`$ and a spectral function like for free fermions, whereas the latter case leads to the Brinkmann-Rice continuum ($`Z(Q)=1/2`$). A magnetic state inbetween the classical Néel and the Heisenberg case could, in principle, also be considered, for which the one-hole spectral function was derived in Ref. treating the spin-fluctuations as perturbation.
## IV Spinon dispersion
In real systems the ratio $`J/t`$ is roughly 0.3. It means that they are in the regime of strong coupling and the above consideration correctly describe the largest energy scale $`t`$. Now, we want to estimate the corrections that arise from other terms of the Hamiltonian (5). First we note that $`v_{k,q}`$ are eigenoperators for the $`t_3`$-term
$$[v_{k,q},\widehat{t}_3]=2t_3v_{k,q}\mathrm{cos}2q,$$
(30)
which leads to the replacement $`ϵ_h(q)ϵ_h(q)+2t_3\mathrm{cos}2q`$ in the denominator of (22). The commutation with $`\widehat{J}`$ gives (see Appendix A for the details)
$$[v_{k,q},\widehat{J}]=J\mathrm{cos}(kq)v_{k,q}+\frac{J}{2}(v_{k,q}^{}+v_{k,q}^{\prime \prime }),$$
(31)
where $`v_{k,q}^{}`$ and $`v_{k,q}^{\prime \prime }`$ are Fourier transforms of the operators
$$v_{m,r}^{}=\underset{\gamma ,\alpha _1,\mathrm{},\alpha _r}{}X_m^{\sigma \gamma }X_{mg}^{\gamma \alpha _1}X_{m+g}^{\alpha _1\alpha _2}\mathrm{}X_{m+r1}^{\alpha _{r1}\alpha _r}X_{m+r}^{\alpha _r0},$$
(32)
$$v_{m,r}^{\prime \prime }=\underset{\gamma ,\alpha _1,\mathrm{},\alpha _r}{}X_m^{\sigma \alpha _1}X_{m+g}^{\alpha _1\alpha _2}\mathrm{}X_{m+rg}^{\alpha _{r1}\alpha _r}X_{m+r+g}^{\alpha _r\gamma }X_{m+r}^{\gamma 0}.$$
(33)
It is impressive that terms, which come from the commutation of “inner” $`X_n`$ operators in $`v_{m,r}`$ with $`n`$ between the points $`m`$ and $`m+r`$ cancel each other and only the terms coming from the ends remain. The term $`v_{m,r}^{\prime \prime }`$ presents a distortion of the right end of $`v_{m,r}`$ by means of the exchange part and may be interpreted as the loss of magnetic energy due to the presence of a holon. On the other hand, the term $`v_{m,r}^{}`$, with a distorted left end will be shown to give rise to the spinon dispersion. We really observe the “separate” motion of the holon that is represented by the right end of $`v_{m,r}`$ and of the spinon that is the left end of $`v_{m,r}`$. The holon motion is governed by the $`t`$-term and the spinon motion by the $`J`$-term. We put the word “separate” in quotes because the motion remains correlated due to the set of “inner” $`X_n`$ operators, connecting the ends of $`v_{m,r}`$.
We need an approximate approach to account for $`v_{k,q}^{}+v_{k,q}^{\prime \prime }`$. For this purpose we use the projection technique
$$v_{k,q}^{}+v_{k,q}^{\prime \prime }\frac{\{v_{k,q}^{}+v_{k,q}^{\prime \prime },v_{k,q}^{}\}}{\{v_{k,q},v_{k,q}^{}\}}v_{k,q}.$$
(34)
Now, the Green’s function for the string operator has the form
$$v_{k,q}|v_{k^{},q^{}}^{}=\frac{8\pi ^2\delta (kk^{})\delta (qq^{})Z(kq+\pi )}{\omega 2t\mathrm{cos}q+2t_3\mathrm{cos}2qϵ_s(kq)},$$
(35)
where $`ϵ_s(kq)`$ is defined by the equation
$$\{[v_{k,q},\widehat{J}],v_{k^{},q^{}}^{}\}8\pi ^2\delta (kk^{})\delta (qq^{})ϵ_s(kq)Z(kq+\pi ).$$
(36)
The contribution of the $`v_{m,r}^{\prime \prime }`$ term to the spinon dispersion is determined by an expression of the form
$$\{v_{m,r}^{\prime \prime },v_{m^{},r^{}}^{}\}=\frac{1}{2}\delta _{m+r,m^{}+r^{}}\mathrm{\Omega }_{r,r^{}}^{\prime \prime },$$
(37)
where the precise order of the $`X`$-operators in $`\mathrm{\Omega }_{r,r^{}}^{\prime \prime }`$ can be easily inferred from (33) and is given in the Appendix A. There, it is also shown that for slowly decaying spin correlation functions (as in the present case, see (26)) the correlation functions $`\mathrm{\Omega }_{r,r^{}}^{\prime \prime }`$ can be approximated to be a function of $`rr^{}`$ only, in the way:
$$\mathrm{\Omega }_{r,r^{}}^{\prime \prime }\mathrm{\Omega }_{rr^{},0}^{\prime \prime }\mathrm{\Omega }_1\mathrm{\Omega }_{rr^{}}.$$
(38)
That leads to a constant shift of the energy $`ϵ_s`$ as the only effect of $`v_{m,r}^{\prime \prime }`$ which will be neglected further on. The contribution of the $`v_{m,r}^{}`$ term can be written analogously to (37), defining the spin correlation functions $`\mathrm{\Omega }_{r,r^{}}^{}=\mathrm{\Omega }_{rr^{},0}^{}`$. The correlation functions $`\mathrm{\Omega }_{l,0}^{}`$ differ from $`\mathrm{\Omega }_{l+1}`$ only by the exchange of two $`X`$-operators. Therefore, for large $`l`$, we may expect that
$$\mathrm{\Omega }_{l,0}^{}\mathrm{\Omega }_{l+1}.$$
(39)
That leads after Fourier transformation to the contribution of the $`v_{m,r}^{}`$ term to the spinon dispersion (Appendix A). Together with the contribution of $`v_{k,q}`$ we obtain for the spinon dispersion
$$ϵ_s(Q\pi )Z(Q)=\frac{J}{2}\left\{\mathrm{cos}Q\left[Z(Q)+\frac{1}{2}\right]+\frac{1}{2}\mathrm{\Omega }_{0,0}^{}\mathrm{\Omega }_1\frac{1}{2}\mathrm{sin}Q_0^{2\pi }\frac{d\kappa }{2\pi }Z(\kappa )\left(\mathrm{cot}\frac{Q\kappa }{2}+\mathrm{cot}\frac{Q+\kappa }{2}\right)\right\},$$
(40)
and the hole spectral function becomes
$$A(k,\omega )=_\pi ^{+\pi }\frac{dQ}{\pi }Z(Q)\delta \left[\omega +2t\mathrm{cos}(kQ)2t_3\mathrm{cos}2(kQ)ϵ_s(Q\pi )\right].$$
(41)
The curve that we obtained for $`ϵ_s`$ with the formula (40) is close to
$$ϵ_s(Q\pi )\alpha J\mathrm{cos}Q,\alpha 2.$$
(42)
as shown in Fig. 1. The functional form (42) is consistent with Bethe-ansatz and field-theoretical considerations. (Sorella and Parola derived a contribution $`J\pi /2\mathrm{cos}Q1.6J\mathrm{cos}Q`$.) We tested it also by comparing the first two terms of Fourier expansion of the product $`\mathrm{cos}QZ(Q)`$ with $`\mathrm{\Omega }_{0,0}^{}`$ and $`\mathrm{\Omega }_{1,0}^{}`$ that give values for $`\alpha `$ in (42) of 2.1 or 1.8, respectively (for the pair correlation functions we took the data of Ref. ). Therefore, we are using the simplified formula (42) instead of (40) in the following analysis of the spectral density. We have checked that the differences are negligible.
The spectral density (for $`t_3=0`$) is shown in Fig. 2. One can clearly distinguish between the spinon and holon features at the lower edge of the spectral density dispersing at an energy scale $`J`$ (from $`k=0`$ to $`k=\pi /2`$) or $`t`$ (from $`k=\pi /2`$ to $`k=\pi `$). At $`k=k^{}`$, which is determined by $`t\mathrm{cos}k^{}=J`$, another holon branch splits off the lower edge of the spectrum and disperses towards $`k=0`$ at an energy scale of $`t`$ (and a corresponding holon branch splits off the upper edge of the spectrum). For $`k`$ values inbetween 0 and $`k^{}`$ one has three peaks in the spectral function (one spinon and lower and upper holon branch). One can easily imagine the situation in the doped case. Then the spinon and holon branches start at the Fermi energy with two different velocities.
In contrast to the 2D case, there is no separate bound state at the lower edge of the spectrum indicating that there are only collective spin and charge excitations. Most of those features were also observed in the ARPES experiment. In the naive picture of spin-charge separation (3) the spectral density would have square root singularities only either at the lower or at the upper edge of the spectrum. In Fig. 2, however, there are additional holon branches due to the square root singularity in $`Z(Q)`$ (see also Ref. ). For $`J0`$ the spinon feature in the spectral density, i.e. the lower edge of the spectrum, becomes completely flat between $`0`$ and $`\pi /2`$. The corresponding pictures were already given in Ref. . Fig. 2 agrees also qualitatively with the finite cluster results.
## V Majumdar-Ghosh model
We have shown that our approach is applicable for any magnetic state for $`J0`$. Now, we are going to present the spectral function of one hole in the $`t`$-$`J`$-$`J^{}`$ model with the special frustration $`J^{}=J/2`$ (called here Majumdar-Ghosh (MG) model for simplicity). In that case rigorous analytic results may be obtained since the ground-state wave function of the MG spin Hamiltonian is exactly known. It is the combination of two simple dimer states
$$\mathrm{\Psi }_{MG}=(\mathrm{\Phi }_1+\mathrm{\Phi }_2)/\sqrt{2},$$
(43)
where
$$\mathrm{\Phi }_1=\underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}[2n,2n+1],\mathrm{\Phi }_2=\underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}[2n1,2n],$$
and the singlet bond is denoted as
$$[l,m]\frac{1}{\sqrt{2}}\underset{\sigma }{}\sigma X_l^{\sigma 0}X_m^{\sigma 0}|vac.$$
We are considering the MG model as a representative example for the case that there is a gap in the spin excitation spectrum (and also in the charge channel). To give the result for the spectral density in the strong coupling limit $`J0`$ one has to find the modified quasiparticle residue $`Z(Q)`$ in (29). It can be simply derived from the correlation functions (see also Ref. )
$`\mathrm{\Omega }_l`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[\mathrm{\Phi }_1|\mathrm{\Omega }_l|\mathrm{\Phi }_1+\mathrm{\Phi }_2|\mathrm{\Omega }_l|\mathrm{\Phi }_2\right]`$ (44)
$`\mathrm{\Omega }_{2n}`$ $`=`$ $`({\displaystyle \frac{1}{2}})^n,\mathrm{\Omega }_{2n+1}={\displaystyle \frac{1}{2}}({\displaystyle \frac{1}{2}})^{n+1},n0,`$ (45)
in the following way
$$Z(Q)=\frac{1}{2}+\underset{n=1}{\overset{\mathrm{}}{}}[(\frac{1}{2}\frac{\mathrm{e}^{iQ}}{4})(\frac{\mathrm{e}^{2iQ}}{2})^n+h.c.]=\frac{3}{2}\frac{1+\mathrm{cos}Q}{5+4\mathrm{cos}2Q}.$$
(46)
The corrections for small $`Jt`$ may only be derived approximatively and we present two methods, projection method and variational procedure having different accuracy.
### A Projection method
First we calculate the spinon dispersion $`ϵ_s`$ in the same way as it was done in the Heisenberg case in Sec. IV. But we should keep in mind that its applicability is less justified for the MG model than for the pure $`t`$-$`J`$ model due to the much faster decay of spin correlation functions (compare (45) with (26)). As before, we approximate $`\mathrm{\Omega }_{r,r^{}}^{\prime \prime }\mathrm{\Omega }_1\mathrm{\Omega }_{rr^{}}`$ which results in a constant energy shift from the $`v_{k,q}^{\prime \prime }`$ term. Therefore, the first contribution to $`ϵ_s`$ coming from $`\widehat{J}`$ is merely determined by $`\mathrm{\Omega }_{r,r^{}}^{}`$ (see Appendix B) resulting in
$$ϵ_{sJ}(Q\pi )=2J\mathrm{cos}Q.$$
(47)
We have a second contribution to $`ϵ_s`$ from $`\widehat{J}^{}`$
$$ϵ_{sJ^{}}(Q\pi )=J^{}\left[4\mathrm{cos}Q+\frac{5}{4}+\mathrm{cos}2Q\right].$$
(48)
We see that for $`J^{}=J/2`$ the terms proportional $`\mathrm{cos}Q`$ cancel and we find
$$ϵ_s(Q\pi )=J\left[\frac{5}{8}+\frac{1}{2}\mathrm{cos}2Q\right],$$
(49)
which is symmetric around $`\pi /2`$.
The spectral density is presented in Fig. 3. We see that in contrast to the $`t`$-$`J`$ model the structures coming from $`Z(Q)`$ (the holon branches) are much less pronounced, whereas square root singularities exist at the lower and upper edges of the spectrum. Their intensities are proportional to $`Z(k)`$ or $`Z(k\pi )`$ at the lower and upper edges, respectively. Therefore, the square root singularity vanishes for $`k=\pi `$ at the lower edge. Furthermore, one can see that the low energy region for $`k`$ between $`\pi /2`$ and $`\pi `$ being empty in Fig. 2 is now filled with states. The spectrum becomes more symmetric around $`\pi /2`$ and the low-energy edge is given by the spinon dispersion $`ϵ_s(Q)`$. The strong damping of the holon branch is due to the suppression of the singularity at the spinon Fermi edge (at $`Q=\pi /2`$ in $`Z(Q)`$). It is a universal feature for any 1D magnetic state having a gap in the spin excitation spectrum. The suppression of holon weight was also found by Voit for the Luther-Emery phase in the Luttinger liquid. The form of the spectral density in Fig. 3 resembles also roughly the exact diagonalization study in Ref. . But a single bound state with a finite spectral weight that was obtained there, is missing in Fig. 3. That deficiency is due to the special projection procedure (34) which can only result in a continuous spectral density. Therefore, one has to go beyond the projection method.
### B Variational ansatz
Here we will use the set of string operators (18) as a set defining a variational wave function for the whole Hamiltonian. Due to the knowledge of the exact ground-state (43) all necessary matrix elements can be calculated without any further approximation. More precisely, we will diagonalize the Hamiltonian $`\widehat{H}=\widehat{t}+\widehat{J}+\widehat{J}^{}`$ in the space spanned by the set of basis operators
$$v_{k,r}=\frac{1}{\sqrt{L}}\underset{\mathrm{}}{\overset{+\mathrm{}}{}}\text{e}^{ik(m+r)}v_{m,r},$$
(50)
where $`L`$ is the number of lattice sites and $`v_{m,r}`$ was defined in (18). For that purpose one has to calculate the overlap matrix resulting in
$$S_{r,r^{}}=\{v_{k,r},v_{k,r^{}}^{}\}=\frac{1}{2}\mathrm{\Omega }_{rr^{}}.$$
(51)
The kinetic energy part of the Hamilton matrix is given by:
$$\frac{t}{2}E_{r,r^{}}^k=\{[v_{k,r},\widehat{t}],v_{k,r^{}}^{}\}=\frac{t}{2}\left(\text{e}^{ik}\mathrm{\Omega }_{rr^{}1}+\text{e}^{ik}\mathrm{\Omega }_{rr^{}+1}\right).$$
(52)
The calculation of the exchange part of the Hamilton matrix is quite lengthy but straightforward. It shall not be given here in detail. To present the results we define a matrix $`\underset{¯}{\underset{¯}{E}}^x`$:
$$\frac{J}{4}E_{r,r^{}}^x=\{[v_{k,r},\widehat{J}+\widehat{J^{}}],v_{k,r^{}}^{}\},$$
(53)
whose matrix elements are listed in Appendix C. The Hamilton matrix is then given by
$$\underset{¯}{\underset{¯}{E}}=\frac{t}{2}\underset{¯}{\underset{¯}{E}}^k+\frac{J}{4}\underset{¯}{\underset{¯}{E}}^x,$$
(54)
and the matrix GF
$$G_{r,r^{}}=v_{k,r}|v_{k,r^{}}^{}$$
(55)
can be found by solving the equation
$$(\omega +i\mathrm{\Gamma }+\underset{¯}{\underset{¯}{E}}\underset{¯}{\underset{¯}{S}}^1)\underset{¯}{\underset{¯}{G}}=\underset{¯}{\underset{¯}{S}},\mathrm{\Gamma }>0.$$
(56)
Finally, the GF (16) can be obtained by $`G(k,\omega )=2G_{0,0}`$.
The numerical results for $`J=0`$ and $`J=0.4`$ at three different momenta are presented in Fig. 4. The curves for $`J=0`$ coincide with the analytic expression (29,46). A number of 400 basis functions and a broadening of $`\mathrm{\Gamma }=0.05`$ are sufficient to reach the thermodynamic limit in contrast to the exact diagonalization method yielding only a sequence of $`\delta `$-peaks. For $`J=0.4`$ we can confirm the features found by the projection method, i.e. the low energy intensity between $`\pi /2`$ and $`\pi `$, the symmetric spinon dispersion and the overdamped holon branch. In addition, the exchange terms produce two new features not present in Sec. VA: a resonance peak near zero energy and a bound state below the continuum. The resonance peak is visible near $`k=\pi `$ and becomes an antiresonance near $`k=0`$. Careful inspection of the exact diagonalization data indicates also a very high peak at the resonance position for $`k=\pi `$ and a small gap at $`k=0`$, but a better understanding of the resonance/antiresonance feature is still required.
The bound state is not visible in Fig. 4 due to the broadening $`\mathrm{\Gamma }`$ which is too large. Instead, we present in Fig. 5 the spectral weight of the lowest eigenstate $`w_1`$ for $`k=\pi /2`$ and $`J=0.4`$ in dependence on the number of basis functions. It is clearly seen that the weight tends to a constant value ($`w_10.1`$) in difference to the weight $`w_3`$ of the third eigenstate. At the same time, the separation $`e_1=E_3E_1`$ between the first and the third eigenvalues $`E_{1/3}`$ stays finite for $`N\mathrm{}`$ but the separation is very small ($`e_10.02`$ in units of $`t`$). For $`J=0`$, both $`w_1`$ and $`e_1`$ tend to zero for $`N\mathrm{}`$. That means that the bound state is connected with the presence of a gap in the spin excitation spectrum.
## VI Ideal paramagnetic state
Such a state is realized for very high temperatures $`T`$, much larger than the exchange energy $`k_BTJ`$. In that case spins at neighboring sites are completely uncorrelated. But the temperature is assumed to be lower than the Hubbard $`U`$ such that the constraint of no double occupancy is preserved. Then the correlation functions become simply
$$\mathrm{\Omega }_l=\left(\frac{1}{2}\right)^l,$$
which results in
$$Z(Q)=\frac{3}{8}\frac{1}{\frac{5}{4}+\mathrm{cos}Q}.$$
(57)
The calculation of the spinon part (without frustration) gives
$$ϵ_s(Q\pi )=\frac{J}{2}\left[2\mathrm{cos}Q+\frac{1}{2}\right].$$
(58)
To calculate it one has to note that (39) is no approximation in the present case. The effect of the $`\mathrm{\Omega }_{r,r^{}}^{\prime \prime }`$ terms can only be treated approximatively (see (38)) but it was checked by the variational method that its influence on the spectral function can be neglected.
The information on $`Z(Q)`$ and $`ϵ_s`$ is sufficient to calculate the spectral function (Fig. 6). It is surprising that the strong singularities at the band edges survive despite the large temperature. The lower edge disperses according to the dispersion of the spinon (58) with a width proportional to $`J`$ and has its minimum at $`k=\pi `$ (in contrast to the frustrated case Fig. 3 with a minimum at $`k=\pi /2`$). But a peak connected with the holon dispersion proportional to $`t`$ is not seen in Fig. 6. Such a peak appears in the finite temperature spectral function of the 2D $`t`$-$`J`$ model and it can be expected since the first moment of the spectral function disperses according to $`t\mathrm{cos}k`$. Its absence in 1D is a nontrivial and unexpected result. It can be understood in the present context since the holon branch is strongly damped due to the suppression of the singularity in $`Z(Q)`$ at $`Q=\pi /2`$. Apparently, that suppression is more strong in (57) than in the frustrated case (46) such that the holon branch is still visible in Fig. 3 but it disappears nearly in Fig. 6. One should note that the above result holds only in the region $`Jk_BTU`$. One may speculate that a further increase of the temperature such that the constraint of no double occupancy is lifted should lead to drastic changes in the spectral function. The strong singularities at the lower or upper band edges should disappear and a free dispersion should become visible.
## VII Conclusion
In conclusion we could derive analytic expressions for the spectral function of one hole in several magnetic states. The expressions are rigorous in the limit $`J0`$, but our approach allows also to calculate the small $`J`$ corrections. We analyzed the frustration and temperature effects. Results were given for the special frustration $`J^{}=J/2`$ with a gap in the spin excitation spectrum and for the ideal paramagnetic case. Both effects, frustration and temperature, lead to low-energy excitations between $`\pi /2`$ and $`\pi `$, and to a strong damping of the holon branches in the spectral function caused by the suppression of the singularity at the Fermi edge of spinons. The exchange terms in the MG model were found to be responsible for the finite weight of the lowest eigenstate and its finite, but small, energy separation from the rest of the spectrum, i.e. the bound state. The proposed scenario of holon branch damping seems to be a universal feature of frustration and temperature. Therefore, our results are of direct importance for photoemission experiments on strongly frustrated 1D compounds like CuGeO<sub>3</sub>, for instance. However, edge-shared cuprate chains have a smaller energy scale and less ideal 1D behavior in comparison with corner-shared compounds, which hinders direct comparison with experiment. But it cannot be excluded that a small frustration is also present in SrCuO<sub>2</sub> such that our study gives one possible reason, why no real, separate holon branch could be observed in the experimental spectra of SrCuO<sub>2</sub> between $`k=0`$ and $`\pi /2`$. In the spin gap case we found a very small energy separation of the bound state from the continuum such that it is nearly impossible to detect it in a photoemission experiment.
Acknowledgements
The authors thank the DFG, the INTAS organization (project No. INTAS-97-11066), and the Max-Planck society for financial support and S.-L. Drechsler, W. Brenig, K. Becker, A. Muramatsu and H. Eschrig for discussions. R.K. acknowledges the hospitality of the University of Technology Dresden, where the main part of this work has been carried out. R.K. also thanks O.A. Starykh for valuable discussions and providing the water-way in the ocean of literature on 1D magnetism.
Appendix A: Spinon dispersion of the Heisenberg case
In this Appendix we outline the main steps to derive the spinon dispersion of the pure $`t`$-$`J`$ model using the projection method. For long chains of $`X`$-operators it is convenient to introduce the notations
$$\underset{\alpha _1,\mathrm{},\alpha _r}{}X_{n_1}^{\sigma \alpha _1}X_{n_2}^{\alpha _1\alpha _2}X_{n_3}^{\alpha _2\alpha _3}\mathrm{}X_{n_{r1}}^{\alpha _{r1}\alpha _r}X_{n_r}^{\alpha _r0}(n_1|n_2|n_3|\mathrm{}|n_{r1}|n_r]$$
and
$$\underset{\sigma ,\alpha _1,\mathrm{},\alpha _r}{}X_{n_1}^{\sigma \alpha _1}X_{n_2}^{\alpha _1\alpha _2}X_{n_3}^{\alpha _2\alpha _3}\mathrm{}X_{n_{r1}}^{\alpha _{r1}\alpha _r}X_{n_r}^{\alpha _r\sigma }(n_1|n_2|n_3|\mathrm{}|n_{r1}|n_r).$$
which means especially that (18) may be rewritten as $`v_{m,r}(m|m+g|m+2g|\mathrm{}|m+r]`$. In such a notation we obtain for the commutation with the Heisenberg Hamiltonian
$$[X_m^{\sigma 0},\widehat{J}]=\frac{J}{2}\underset{g,\gamma }{}X_{m+g}^{\sigma \gamma }X_m^{\gamma 0}=\frac{J}{2}\underset{g}{}(m+g|m],$$
(A.1)
$$[X_m^{\alpha \beta },\widehat{J}]=\frac{J}{2}\underset{g,\gamma }{}(X_m^{\alpha \gamma }X_{m+g}^{\gamma \beta }X_{m+g}^{\alpha \gamma }X_m^{\gamma \beta })=\frac{J}{2}\underset{g}{}\{(m|m+g|(m+g|m|\},$$
(A.2)
and then
$`[v_{m,r},\widehat{J}]`$ $`=`$ $`{\displaystyle \frac{J}{2}}\{(m|mg|m+g|m+2g|\mathrm{}|m+r](mg|m|m+g|m+2g|\mathrm{}|m+r]`$ (A.4)
$`(m+g|m+2g|\mathrm{}|m+r](m|\mathrm{}|m+rg|m+r+g|m+r]\},`$
where $`g=\mathrm{sign}(r)`$. In deriving (LABEL:comJapp) it is important that the commutation of the ”inner” operators $`X_{m+l}^{\alpha \beta }`$ with $`l<r`$ do not give rise to additional terms since the corresponding sums cancel each other. That is a direct consequence of one-dimensionality.
Now, we consider the holon contribution to $`ϵ_s`$ coming from $`v_{k,q}^{\prime \prime }`$
$$\{v_{k,q}^{\prime \prime },v_{k^{},q^{}}^{}\}=2\pi \delta (kk^{})\underset{r,r^{}}{}\text{e}^{iq^{}r^{}iqr+ik(rr^{})}\mathrm{\Omega }_{r,r^{}}^{\prime \prime }$$
(A.6)
with
$$\mathrm{\Omega }_{r,r^{}}^{\prime \prime }=(m|\mathrm{}|m+rg|m+r+g|m+r|m+rg^{}|\mathrm{}|m+rr^{}),$$
(A.7)
and $`g=\text{sign}(r)`$, $`g^{}=\text{sign}(r^{})`$. We see that in general $`\mathrm{\Omega }_{r,r^{}}^{\prime \prime }`$ depends both on $`rr^{}`$ and on $`r`$. Eqn. (A.6) can also be written as
$$\{v_{k,q}^{\prime \prime },v_{k^{},q^{}}^{}\}=2\pi \delta (kk^{})\underset{l}{}\text{e}^{i(kq)l}S_l,$$
(A.8)
with
$$S_l=\underset{r=\mathrm{}}{\overset{+\mathrm{}}{}}\text{e}^{i(qq^{})r}\mathrm{\Omega }_{r,rl}^{\prime \prime }.$$
Due to the slow decay of spin correlation functions in the 1D Heisenberg state, one can expect that the main contribution to $`S_l`$ comes from regions where $`|r||l|`$. There holds $`g=g^{}`$ and we may rewrite and approximate (A.7) by
$$\mathrm{\Omega }_{r,rl}^{\prime \prime }=(0|\mathrm{}|l)(r+g|r)\mathrm{\Omega }_l\mathrm{\Omega }_1.$$
(A.9)
Then, the explicit dependence on $`r`$ drops out and we obtain
$$\{v_{k,q}^{\prime \prime },v_{k^{},q^{}}^{}\}=8\pi ^2\delta (kk^{})\delta (qq^{})Z(kq+\pi )\mathrm{\Omega }_1,$$
(A.10)
i.e. a simple constant shift of the energy $`ϵ_s`$.
The contribution of the $`v_{k,q}^{}`$ term to the spinon dispersion is determined by the sequence of spin operators
$$\mathrm{\Omega }_{r,r^{}}^{}=(m|mg|m+g|\mathrm{}|m+r|m+rg^{}|\mathrm{}|m+rr^{})$$
(A.11)
instead of (A.7). The expectation value of that term has to be calculated for the magnetic system without holes. In difference to $`\mathrm{\Omega }_{r,r^{}}^{\prime \prime }`$, it depends only on $`rr^{}`$ without further approximation
$$\mathrm{\Omega }_{r+l,r}^{}=\mathrm{\Omega }_{l,0}^{}=(0|1|1|\mathrm{}|l),(l>0),$$
(A.12)
and $`\mathrm{\Omega }_{l,0}^{}=\mathrm{\Omega }_{l,0}^{}`$. For $`l=0,1`$ it can be expressed through pair correlation functions
$`\mathrm{\Omega }_{0,0}^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}+2𝐒_0𝐒_2,`$ (A.13)
$`\mathrm{\Omega }_{1,0}^{}`$ $`=`$ $`{\displaystyle \frac{1}{4}}+2𝐒_0𝐒_1+𝐒_0𝐒_2.`$ (A.14)
For large $`l>0`$ we may expect
$$\mathrm{\Omega }_{l,0}^{}\mathrm{\Omega }_{l+1}.$$
(A.15)
Using this approximation we obtain the following contribution to the spinon dispersion $`ϵ_s`$ which stems from the $`v_{k,q}^{}`$ term
$$\frac{J}{2}\{v_{k,q}^{},v_{k^{},q^{}}^{}\}=8\pi ^2\delta (qq^{})\delta (kk^{})ϵ_s^{}(kq)Z(kq+\pi )$$
(A.16)
with
$$ϵ_s^{}(k)Z(k+\pi )=\frac{J}{4}\underset{l=\mathrm{}}{\overset{+\mathrm{}}{}}\text{e}^{ikl}\mathrm{\Omega }_{l,0}^{}.$$
(A.17)
After some algebra we find
$$4ϵ_s^{}(Q\pi )Z(Q)/J=\mathrm{\Omega }_{0,0}^{}2\mathrm{\Omega }_12\mathrm{cos}Q\left[Z(Q)\frac{1}{2}\right]\mathrm{sin}QY(Q),$$
(A.18)
where
$$Y(Q)=\underset{l=1}{\overset{+\mathrm{}}{}}2(1)^l\mathrm{sin}(Ql)\mathrm{\Omega }_l=\frac{1}{2\pi }_0^{2\pi }𝑑\kappa Z(\kappa )\left[\mathrm{cot}\frac{Q\kappa }{2}+\mathrm{cot}\frac{Q+\kappa }{2}\right].$$
(A.19)
The complete expression for the spinon dispersion follows from
$$ϵ_s(Q\pi )=J\mathrm{cos}Q+ϵ_s^{}(Q\pi )+\mathrm{const},$$
and is given in Eqn. (40) neglecting the constant energy shift.
Appendix B: Majumdar-Ghosh model with projection method
The commutation with the frustration Hamiltonian (14) is very similar to (LABEL:comJapp). It gives
$$[v_{m,r},\widehat{J}^{}]=\frac{J^{}}{2}\left\{v_{m,r}^{(3)}+v_{m,r}^{(4)}\right\}$$
(B.1)
where
$`v_{m,r}^{(3)}`$ $`=`$ $`\{(m|m2g|m+g|\mathrm{}|m+r](m2g|m|m+g|\mathrm{}|m+r]`$ (B.4)
$`(m|m+g)(m+2g|\mathrm{}|m+r]+(m|m+g|mg|m+2g|\mathrm{}|m+r]`$
$`(m|mg|m+g|\mathrm{}|m+r]\},`$
$`v_{m,r}^{(4)}`$ $`=`$ $`\{(m|\mathrm{}|m+rg|m+r+2g|m+r]`$ (B.7)
$`(m|\mathrm{}|m+r2g|m+r+g|m+rg|m+r]`$
$`+(m|\mathrm{}|m+rg|m+r+g|m+r]\}.`$
Again we see that all terms coming from commutations at “inner” operators cancel. But now the motion of spinons becomes more complicated. The same considerations that show the absence of dispersion from the $`v_{m,r}^{\prime \prime }`$ term are applicable to $`v_{m,r}^{(4)}`$.
Taking into account that
$$\mathrm{\Omega }_{0,0}^{}=\frac{1}{2},\mathrm{\Omega }_{2n,0}^{}=\frac{1}{4}\left(\frac{1}{2}\right)^n,\mathrm{\Omega }_{2n1,0}^{}=\left(\frac{1}{2}\right)^n,n>0,$$
we obtain the following contributions to $`ϵ_s=ϵ_{sJ}+ϵ_{sJ^{}}`$ (we drop dispersionless terms). From $`\widehat{J}`$ comes
$`ϵ_{sJ}(Q\pi )Z(Q)`$ $`=`$ $`{\displaystyle \frac{J}{2}}\left[2\mathrm{cos}QZ(Q)+Z^{}(Q)\right],`$ (B.8)
$`Z^{}(Q)`$ $`=`$ $`{\displaystyle \frac{3(\frac{1}{8}+\mathrm{cos}Q)}{5+4\mathrm{cos}2Q}}+{\displaystyle \frac{1}{4}}+{\displaystyle \frac{1}{8}}=Z(Q)2\mathrm{cos}Q.`$ (B.9)
Thus
$$ϵ_{sJ}(Q\pi )=2J\mathrm{cos}Q$$
(B.10)
has the same form (42) that we have assumed for the $`t`$-$`J`$ model.
For the term that comes from the left distorted end of $`v_{m,r}`$ due to $`\widehat{J}^{}`$ we have
$$\{v_{m,r_1}^{(3)},v_{m+r_1r_2,r_2}^{}\}\frac{1}{2}\mathrm{\Omega }_{r_1r_2}^{(3)},$$
$$\mathrm{\Omega }_{2n}^{(3)}=\left(\frac{1}{2}\right)^n,\mathrm{\Omega }_{2n1}^{(3)}=4\left(\frac{1}{2}\right)^n,n>1,$$
$$\mathrm{\Omega }_0^{(3)}=\left(\frac{1}{2}\right),\mathrm{\Omega }_1^{(3)}=\frac{5}{4}.$$
The contribution from $`\widehat{J}^{}`$ is
$`ϵ_{sJ^{}}(Q\pi )Z(Q)`$ $`=`$ $`{\displaystyle \frac{J^{}}{2}}\left[4Z^{}(Q)+{\displaystyle \frac{3}{4}}+{\displaystyle \frac{6}{8}}\mathrm{cos}Q\right]`$ (B.11)
$`=`$ $`Z(Q){\displaystyle \frac{J^{}}{2}}\left[8\mathrm{cos}Q+2({\displaystyle \frac{5}{4}}+\mathrm{cos}2Q)\right]`$ (B.12)
and we obtain (48).
Appendix C: Matrix elements of the variational basis set
The matrixelements of $`E_{r,r^{}}^x`$ in the neighborhood of $`r,r^{}=0`$ are given by:
| | $`r^{}`$ | -3 | -2 | -1 | 0 | 1 | 2 | 3 |
| --- | --- | --- | --- | --- | --- | --- | --- | --- |
| $`r`$ | | | | | | | | |
| -3 | | 3/8 | 0 | 3/16 | -3/32 | -9/32 | 3/32 | 3/16 |
| -2 | | 0 | 3/8 | -9/16 | 3/8 | 0 | -3/8 | 3/32 |
| -1 | | 3/16 | -9/16 | 3/8 | -3/16 | 3/8 | 0 | -9/32 |
| 0 | | -3/32 | 3/8 | -3/16 | 0 | -3/16 | 3/8 | -3/32 |
| 1 | | -9/32 | 0 | 3/8 | -3/16 | 3/8 | -9/16 | 3/16 |
| 2 | | 3/32 | -3/8 | 0 | 3/8 | -9/16 | 3/8 | 0 |
| 3 | | 3/16 | 3/32 | -9/32 | -3/32 | 3/16 | 0 | 3/8 |
One has two different regions in the matrix. The first one is defined for $`r>0,r^{}>0`$ and $`rr^{}+2`$ where we have the matrix elements:
$$\begin{array}{cc}r=2n,r^{}=2m\text{or}r=2n+1,r^{}=2m+1:\hfill & E_{r,r^{}}^x=\frac{3}{8}\left(\frac{1}{2}\right)^{nm},\hfill \\ r=2n+1,r^{}=2m:\hfill & E_{r,r^{}}^x=\frac{3}{8}\left(\frac{1}{2}\right)^{nm},\hfill \\ r=2n,r^{}=2m1:\hfill & E_{r,r^{}}^x=\frac{3}{16}\left(\frac{1}{2}\right)^{nm+1},\hfill \end{array}$$
(C.1)
and the second one for $`r^{}2,r2`$ with
$$E_{r,r^{}}^x=\frac{3}{2}\mathrm{\Omega }_{rr^{}}.$$
(C.2)
There are special matrix elements along the diagonal (3/8) and along the side diagonal (alternatively -9/16 or 0) and also for the two lines $`(n1)`$:
$`E_{2n,0}^x={\displaystyle \frac{3}{4}}\left({\displaystyle \frac{1}{2}}\right)^n`$ $`E_{2n+1,0}^x={\displaystyle \frac{3}{16}}\left({\displaystyle \frac{1}{2}}\right)^n`$
$`E_{2n,1}^x=0`$ $`E_{2n+1,1}^x={\displaystyle \frac{9}{16}}\left({\displaystyle \frac{1}{2}}\right)^n`$
The matrix is filled by
$$E_{r,r^{}}^x=E_{r,r^{}}^x=E_{r^{},r}^x.$$
(C.3)
Figures
Fig. 1: Comparison of the spinon dispersion $`ϵ_s(Q\pi )`$ as calculated from the projection method (broken line, $`J=0.4`$) with $`2J\mathrm{cos}Q`$ (full line).
Fig. 2: Spectral density of the $`t`$-$`J`$ model for $`J=0.4`$ and $`t=1`$.
Fig. 3: Spectral density of the frustrated $`t`$-$`J`$ model ($`J=0.4`$ and $`t=1`$) at the special frustration $`J^{}=0.5J`$ (using the Majumdar-Ghosh wave function) within the projection method.
Fig. 4: Spectral density of the Majumdar-Ghosh model $`A(k,\omega )`$ for three different momenta $`k/\pi `$ and $`t=1`$, $`J=0.4`$ (full lines) or $`J=0`$ (dashed lines) with a variational set of 400 basis functions and a broadening of $`\mathrm{\Gamma }=0.05`$.
Fig. 5: Weight $`w_1`$ and energy separation $`e_1=E_3E_1`$ of the lowest eigenvalue $`E_1`$ at $`k/\pi =0.5`$, $`t=1`$, $`J=0.4`$ (full lines) as a function of the inverse number of basis functions $`1/N`$. The dashed lines are the weights $`w_3`$ and the energy separation $`e_3=E_5E_3`$ of the third eigenvalue $`E_3`$.
Fig. 6: Spectral density of the $`t`$-$`J`$ model ($`J=0.4`$ and $`t=1`$) in the ideal paramagnetic state.
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# Numerical Simulations in Cosmology III: Dark Matter Halos
## 1 Introduction
During the last decade there was an increasingly growing interest in testing predictions of variants of the cold dark matter (CDM) models at subgalactic ($`100\mathrm{kpc}`$) scales. This interest was first induced by indications that observed rotation curves in the central regions of dark matter dominated dwarf galaxies are at odds with predictions of hierarchical models. Specifically, it was argued (Moore 1994; Flores & Primack 1994) that circular velocities, $`v_c(r)[GM(<r)/r]^{1/2}`$, at small galactocentric radii predicted by the models are too high and increase too rapidly with increasing radius compared to the observed rotation curves. The steeper than expected rise of $`v_c(r)`$ implies that the shape of the predicted halo density distribution is incorrect and/or that the DM halos formed in CDM models are too concentrated (i.e., have too much of their mass concentrated in the inner regions).
In addition to the density profiles, there is an alarming mismatch in predicted abundance of small-mass ($`10^810^9h^1\text{M}\text{}`$) galactic satellites and the observed number of satellites in the Local Group (Kauffmann, White & Guiderdoni 1993; Klypin et al. 1999; Moore et al. 1999). Although this discrepancy may well be due to feedback processes such as photoionization that prevent gas collapse and star formation in the majority of the small-mass satellites (e.g., Bullock, Kravtsov & Weinberg 2000), the mass scale at which the problem sets in is similar to the scale in the spectrum of primordial fluctuations that may be responsible for the problems with density profiles. In the age of precision cosmology that forthcoming MAP and Planck cosmic microwave background anisotropy satellite missions are expected to bring, tests of the cosmological models at small scales may prove to be the final frontier and the ultimate challenge to our understanding of the cosmology and structure formation in the Universe. However, this obviously requires detailed predictions and checks from the theoretical side and higher resolution/quality observations and thorough understanding of their implications and associated caveats from the observational side. In this paper we focus on the theoretical predictions of the density distribution of DM halos and some problems with comparing these predictions to observations.
A systematic study of halo density profiles for a wide range of halo masses and cosmologies was done by Navarro, Frenk & White (1996, 1997; hereafter NFW), who argued that analytical profile of the form $`\rho (r)=\rho _s(r/r_s)^1(1+r/r_s)^2`$ provides a good description of halo profiles in their simulations for all halo masses and in all cosmologies. Here, $`r_s`$ is the scale radius which, for this profile corresponds to the scale at which $`d\mathrm{log}\rho (r)/d\mathrm{log}r|_{r=r_s}=2`$. The parameters of the profile are determined by the halo’s virial mass $`M_{\mathrm{vir}}`$ and concentration defined as $`cr_{\mathrm{vir}}/r_s`$. NFW argued that there is a tight correlation between $`c`$ and $`M_{\mathrm{vir}}`$, which implies that the density distributions of halos of different masses can in fact be described by a one-parameter family of analytical profiles. Further studies by Kravtsov, Klypin & Khokhlov (1997), Kravtsov et al. (1999), Jing (1999), Bullock et al. (2000), although confirming the $`c(M_{\mathrm{vir}})`$ correlation, indicated that there is a significant scatter in the density profiles and concentrations for DM halos of a given mass.
Following the initial studies by Moore (1994) and Flores & Primack (1994), Kravtsov et al. (1999) presented a systematic comparison of the results of numerical simulations with rotation curves of a sample of seventeen dark matter dominated dwarf and low surface brightness (LSB) galaxies. Based on these comparisons, we argued that there does not seem to be a significant discrepancy in the shape of the density profiles at the scales probed by the numerical simulations ($`0.020.03r_{\mathrm{vir}}`$, where $`r_{\mathrm{vir}}`$ is halo’s virial radius). However, these conclusions were subject to several caveats and had to be tested. First, observed galactic rotation curves had to be re-examined more carefully and with higher resolution. The fact that all of the observed rotation curves used in earlier analyses were obtained using relatively low-resolution HI observations, required checks of the possible beam smearing effects. Also, a possibility of non-circular random motions in the central regions that could modify the rotation velocity of the gas (e.g., Binney & Tremain 1987, p. 198) had to be considered. Second, the theoretical predictions had to be tested for convergence and extended to scales $`0.01r_{\mathrm{vir}}`$.
Moore et al. (1998; see also a more recent convergence study by Ghigna et al. 1999) presented a convergence study and argued that mass resolution has a significant impact on the central density distribution of halos. They argued that at least several million particles per halo are required to model reliably the density profiles at scales $`0.01r_{\mathrm{vir}}`$. Based on these results, Moore et al. (1999) advocated a density profile of the form $`\rho (r)(r/r_0)^{1.5}[1+(r/r_0)^{1.5}]^1`$, that behaves similarly ($`\rho r^3`$) to the NFW profile at large radii, but is steeper at small $`r`$: $`\rho r^{1.5}`$. Most recently, Jing & Suto (2000) presented a systematic study of density profiles for halo masses range $`2\times 10^{12}h^1\text{M}\text{}5\times 10^{14}h^1\text{M}\text{}`$. The study was uniform in mass and force resolution featuring $`510\times 10^5`$ particles per halo and force resolution of $`0.004r_{\mathrm{vir}}`$. They found that galaxy-mass halos in their simulations are well fitted by profile<sup>1</sup><sup>1</sup>1Note that his profile is somewhat different than profile advocated by Moore et al., but behaves similarly to the latter at small radii. $`\rho (r)(r/r_0)^{1.5}[1+r/r_0]^{1.5}`$, but that cluster-mass halos are well described by the NFW profile, with logarithmic slope of the density profiles at $`r=0.01r_{\mathrm{vir}}`$ changing from $`1.5`$ for $`M_{\mathrm{vir}}10^{12}h^1\text{M}\text{}`$ to $`1.1`$ for $`M_{\mathrm{vir}}5\times 10^{14}h^1\text{M}\text{}`$. Jing & Suto interpreted these results as evidence that profiles of DM halos are not universal.
Rotation curves of a number of dwarf and LSB galaxies have recently been reconsidered using H$`\alpha `$ observations (e.g., Swaters, Madore & Trewhella 2000; van den Bosch et al. 2000). The results show that for majority of galaxies H$`\alpha `$ rotation curves are significantly different in their central regions than the rotation curves derived from HI observations. This indicates that the HI rotation curves are affected by beam smearing (Swaters et al. 2000). It is also possible that some of the difference may be due to real differences in kinematics of the two tracer gas components (ionized and neutral hydrogen). Preliminary comparisons of the new H$`\alpha `$ rotation curves with model predictions show that NFW density profiles are consistent with the observed shapes of rotation curves (van den Bosch 2000; Navarro & Swaters 2000). Moreover, cuspy density profiles with inner logarithmic slopes as steep as $`1.5`$ also seem to be consistent with the data (van den Bosch 2000). Nevertheless, CDM halos appear to be too concentrated (Navarro & Swaters 2000; McGaugh et al. 2000; Navarro & Steinmetz 2000) as compared to galactic halos, and therefore the problem remains.
New observational and theoretical developments show that comparison of model predictions to the data is not straightforward. Decisive comparisons require reaching convergence of theoretical predictions and understanding the kinematics of the gas in the central regions observed galaxies. In this paper we present convergence tests designed to test effects of mass resolution on the density profiles of halos formed in the currently popular CDM model with cosmological constant ($`\mathrm{\Lambda }`$CDM) and simulated using the multiple mass resolution version of the Adaptive Refinement Tree code (ART). We also discuss some caveats in drawing conclusions about the density profiles from the fits of analytical functions to numerical results and their comparisons to the data.
## 2 Dark Matter Halos: the NFW and the Moore et al. profiles
Before we fit the analytical profiles to real dark matter halos or compare them with observed rotational curves, it is instructive to compare different analytical approximations. Although the NFW and Moore et al. profiles predict different behavior of $`\rho (r)`$ in the central regions of a halo, the scale where this difference becomes significant depends on the specific values of halo’s characteristic density and radius. Table 2 presents different parameters and statistics associated with the two analytical profiles. For the NFW profile more information can be found in Klypin et al. (1999), Lokas & Mamon (2000), and in Widrow (2000).
Each profile is set by two independent parameters. We choose these to be the characteristic density $`\rho _0`$ and radius $`r_s`$. In this case all expressions describing different properties of the profiles have simple form and do not depend on concentration. The concentration or the virial mass appear only in the normalization of the expressions. The choice of the virial radius (e.g., Lokas & Mamon 2000) as a scale unit results in more complicated expressions with explicit dependence on the concentration. In this case, one also has to be careful about definition of the virial radius, as there are several different definitions in the literature. For example, it is often defined as the radius, $`r_{200}`$, within which the average density is 200 times the critical density. In this paper the virial radius is defined as radius within which the average density is equal to the density predicted by the top-hat model: it is $`\delta _{\mathrm{TH}}`$ times the average matter density in the Universe. For the $`\mathrm{\Omega }_0=1`$ case the two existing definitions are equivalent. In the case of $`\mathrm{\Omega }_0=0.3`$ models, however, the virial radius is about 30% larger than $`r_{200}`$.
There is no unique way of defining a consistent concentration for the different analytical profiles. Again, it is natural to use the characteristic radius $`r_s`$ to define the concentration: $`cr_{\mathrm{vir}}/r_s`$. This simplifies expressions. At the same time, if we fit the same dark matter halo with the two profiles, we will get different concentrations because the values of corresponding $`r_s`$ will be different. Alternatively, if we choose to match the outer parts of the profiles (say, $`r>r_s`$) as closely as possible, we may choose to change the ratio of the characteristic radii $`r_{s,\mathrm{NFW}}/r_{s,\mathrm{Moore}}`$ in such a way that both profiles reach the maximum circular velocity $`v_{\mathrm{circ}}`$ at the same physical radius $`r_{\mathrm{max}}`$. In this case, the formal concentration of the Moore et al. profile is 1.72 times smaller than that of the NFW profile. Indeed, with this normalization profiles look very similar in the outer parts as one finds in Figure 1. Table 2 also gives two other “concentrations”. The concentration $`C_{1/5}`$ is defined as the ratio of virial radius to the radius, which encompasses 1/5 of the virial mass (Avila-Reese et al. 1999). For halos with $`C_{\mathrm{NFW}}5.5`$ this 1/5 mass concentration is equal to $`C_{\mathrm{NFW}}`$. One can also define the concentration as the ratio of the virial radius to the radius at which the logarithmic slope of the density profile is equal -2. This scale corresponds to $`r_s`$ for the NFW profile and $`0.35r_s`$ for the Moore et al. profile.
Figure 1 presents comparison of the analytic profiles normalized to have the same virial mass and the same radius $`r_{\mathrm{max}}`$. We show results for halos of low and high values of concentration representative of cluster- and low-mass galaxy halos, respectively. The bottom panels show the profiles, while the top panels show corresponding logarithmic slope as a function of radius. The figure shows that the two profiles are very similar throughout the main body of the halos. Only in the very central region the differences become significant. The difference is more apparent in the logarithmic slope than in the actual density profiles. Moreover, for galaxy-mass halos the difference sets in at a rather small radius $`0.01r_{\mathrm{vir}}`$, which would correspond to scales $`<1\mathrm{kpc}`$ for the typical dark matter dominated dwarf and LSB galaxies. In most analyses involving galaxy-size halos, the differences between NFW and Moore et al. profiles are irrelevant, and NFW profile should provide an accurate description of the density distribution.
Note also that for galaxy-size (e.g., high-concentration) halos the logarithmic slope of the NFW profile does not reach its asymptotic inner value of $`1`$ at scales as small as $`0.01r_{\mathrm{vir}}`$. For $`10^{12}h^1\text{M}\text{}`$ halos the logarithmic slope of the NFW profile is $`1.41.5`$, while for cluster-size halos this slope is $`1.2`$. This dependence of slope at a given fraction of the virial radius on the virial mass of the halo is very similar to the results plotted in Figure 3 of Jing & Suto (2000). These authors interpreted it as evidence that halo profiles are not universal. It is obvious, however, that their results are consistent with NFW profiles and the dependence of the slope on mass can be simply a manifestation of the well-studied $`c_{\mathrm{vir}}(M)`$ relation.
To summarize, we find that the differences between the NFW and the Moore et al. profiles are very small ($`\mathrm{\Delta }\rho /\rho <10\%`$) for radii above 1% of the virial radius. The differences are larger for halos with smaller concentrations. In case of the NFW profile, the asymptotic value of the central slope $`\gamma =1`$ is not achieved even at radii as small as 1%-2% of the virial radius.
## 3 Properties of Dark Matter Halos
Some properties of halos depend on large-scale environment in which the halos are found. We will call a halo distinct if it is not inside a virial radius of another (larger) halo. Halo is called subhalo, if it is inside of another halo. The number of subhalos depends on the mass resolution – the deeper we go, the more subhalos we will find. Most of the results below are based on a simulation, which was complete to masses down to $`10^{11}h^1\text{M}\text{}`$ or, equivalently, to the maximum circular velocity of 100 km/s.
Mass and Velocity Distribution functions. Extensive analysis of halo mass and velocity function was done by Sigad et al. (2000) for halos in the $`\mathrm{\Lambda }`$CDM model. Additional results can also be found in Ghigna et al. (1999); Moore et al. (1999b); Klypin et al. (1999); Gottlöber, Klypin & Kravtsov (1998). Figure 2 compares the mass function of subhalos and distinct halos. The Press-Schechter approximation overestimates the the mass function by a factor of 2 for $`M<5\times 10^{12}h^1\text{M}\text{}`$ and it somewhat underestimates it at larger masses. More advanced approximation given by Sheth & Tormen is more accurate. On scales below $`10^{14}h^1\text{M}\text{}`$ the mass function is close to a power law with the slope $`\alpha 1.8`$. There is no visible difference in the slope for subhalos and for the distinct halos.
For each halo one can measure the maximum circular velocity $`V_{\mathrm{max}}`$. In many cases (especially for subhalos) $`V_{\mathrm{max}}`$ is a better measure of how large is the halo. It is also closer related with observed properties of galaxies hosted by halos. Figure 3 presents the velocity distribution function of different types of halos. In addition to distinct halos and subhalos, we show also isolated halos and halos in groups and clusters. Here isolated halos are defined as halos with mass less than $`10^{13}h^1\text{M}\text{}`$, which are not inside a larger halo and which do not have subhalos more massive than $`10^{11}h^1\text{M}\text{}`$. The velocity function is approximated by a power law $`dn=\mathrm{\Phi }_{}V_{\mathrm{max}}^\beta dV_{\mathrm{max}}`$ with the slope $`\beta 3.8`$ for distinct halos. The slope depends on environment: $`\beta 3.1`$ for halos in groups and $`\beta 4`$ for isolated halos. Klypin et al. (1999) and Ghigna et al. (1999) found that the slope $`\beta 3.84`$ of the velocity function extends to much smaller halos with velocities down to 10 km/s.
Correlation between characteristic density and radius. The halo density profiles are approximated by the Navarro-Frenk-White profile:
$$\rho =\frac{\rho _0}{(r/r_0)[1+r/r_0]^2}$$
(1)
Kravtsov et al. (1999) find the correlation between the two parameters of halos $`\rho _0`$ and $`r_s`$. Fifure 4 compares results for the dark matter halos with those for the the dark matter dominated Low Surface Brightness (LSB) galaxies and dwarf galaxies. Halos are consistent with observational data: smaller halos are denser.
Correlations between mass, concentration, and redshift. Navarro et al. (1997) argued that the halo profiles have a universal shape in the sense that profile is uniquely defined by virial mass of the halo. Bullock et al. (2000) analyzed concentrations of thousands of halos at different redshifts. To some degree they confirm conclusions of Navarro et al. (1997): halo concentration correlates with its mass. But some significant deviations were also found. There is no one-to-one relation between concentration and mass. It appears the the universal profile should only be treated as a trend: halo concentration does increase as the halo mass decreases, but there are large deviations for individual halos from that “universal” shape. Halos have intrinsic scatter of concentration: at $`1\sigma `$ level halos with the same mass have $`\mathrm{\Delta }(\mathrm{log}c_{\mathrm{vir}})=0.18`$ or, equivalently, $`\mathrm{\Delta }V_{\mathrm{max}}/V_{\mathrm{max}}=0.12`$.
Velocity anisotropy. Inside a large halo subhalos or dark matter particles do not move on either circular or radial orbits. Velocity ellipsoid can be measured at each position inside halo. It can be characterized by anisotropy parameter defined as $`\beta (r)=1V_{}^2/2V_r^2`$. Here $`V_{}^2`$ is the velocity dispersion perpendicular to the radial direction and $`V_r^2`$ is the radial velocity dispersion. For pure radial motions $`\beta =1`$. For isotropic velocities $`\beta =0`$. Function $`\beta (r)`$ was estimated for halos in different cosmological models ( see Colín et al. (1999) for references). By studying 12 rich clusters with many subhalos inside each of them, Colín et al. (1999) found that both the subhalos and dark matter particles can be described by the same anisotropy parameter
$$\beta (r)=0.15+\frac{2x}{x^2+4},x=r/r_{\mathrm{vir}}$$
(2)
## 4 Halo profiles: Convergence study
The following results are based on Klypin et al. (2000)
### 4.1 Numerical simulations
Using the ART code (Kravtsov, Klypin & Khokhlov, 1997; Kravtsov, 1999), we simulate a flat low-density cosmological model ($`\mathrm{\Lambda }`$CDM) with $`\mathrm{\Omega }_0=1\mathrm{\Omega }_\mathrm{\Lambda }=0.3`$, the Hubble parameter (in units of $`100\mathrm{kms}^1\mathrm{Mpc}^1`$) $`h=0.7`$, and the spectrum normalization $`\sigma _8=0.9`$. We have run two sets of simulations with $`30h^1\text{Mpc}`$ and $`25h^1\text{Mpc}`$ computational box. The first simulations were run to the present moment $`z=0`$. The second set of simulations had higher mass resolution and therefore produced more halos but were run only to $`z=1`$.
In all of our simulations step in the expansion parameter was chosen to be $`\mathrm{\Delta }a_0=2\times 10^3`$ on the zero level of resolution. This gives about 500 steps for an entire run to $`z=0`$. A test run was done with twice smaller time-step for a halos of comparable mass (but with smaller number of particles) as studied in this paper. We did not find any visible deviations in the halo profile. In the first set of simulations, the highest level of refinement was ten, which corresponds to $`500\times 2^{10}500,000`$ time steps at the tenth level. For the second set of simulation, nine levels of refinement were reached which corresponds to $`128,000`$ steps at the ninth level.
In the following sections we present results for four halos. The first halo ($`A`$) was the only halo selected for resimulation in the first set of simulations. In this case the selected halo was relatively quiescent at $`z=0`$ and had no massive neighbors. The halo was located in a long filament bordering a large void. It was about 10 Mpc away from nearest cluster-size halo. After the high-resolution simulation was completed we found that the nearest galaxy-size halo was about 5 Mpc away. The halo had a fairly typical merging history with $`M(t)`$ track slightly lower than the average mass growth predicted using extended Press-Schechter model. The last major merger event occured at $`z2.5`$; at lower redshifts the mass growth (the mass in this time interval has grown by a factor of three) was due to slow and steady mass accretion.
The second set of simulations was done in a different way. In the low resolution run we selected three halos in a well pronounced filament. Two of the halos are neighbors located at about 0.5 Mpc from each other. The third halo was 2 Mpc away from this pair. Thus, the halos were not selected to be too isolated as was the case in the first set of runs. Moreover, the simulation was analyzed at an earlier moment ($`z=1`$) where halos are more likely to be unrelaxed. Therefore, we consider the halo $`A`$ from the first set as an example of a rather isolated well relaxed halo. In many respects, this halo is similar to halos simulated by other research groups that used multiple mass resolution techniques. The three halos from the second set of simulations can be viewed as representative of more typical halos, not necessarily well relaxed and located in more crowded environments.
Parameters of the simulated dark matter halos are listed in Table 2. Columns in the table present (1) Halo “name” (halos A<sub>1</sub>, A<sub>2</sub>, A<sub>3</sub> are the halo A re-simulated times with different resolutions); (2) redshift at which halo was analyzed; (3-5) virial mass, comoving virial radius, and maximum circular velocity. At $`z=0`$ ($`z=1`$) the virial radius was estimated as the radius within which the average overdensity of matter is 340 (180) times larger than the mean cosmological density of matter at that redshift; (6) the number of particles within the virial radius; (7) the smallest particle mass in the simulation; (8) formal force resolution achieved in the simulation. As we will show below, convergent results are expected at scales larger than four times the formal resolution; (9) halo concentration as estimated from NFW profile fits to halo density profiles; (10) maximum relative error of the NFW fit: $`\rho _{\mathrm{NFW}}/\rho _{\mathrm{halo}}1`$ (the error was estimated inside $`50h^1\text{kpc}`$ radius); (11) the same as in the previous column, but for the fits of profile advocated by Moore et al.
Halo $`A`$ in the first set of simulations was re-simulated three times with increasing mass resolution. For each simulation, we considered outputs at four time moments in the interval to $`z=00.03`$. Parameters of the halos in these simulations averaged over the four moments are presented in the first three rows of the Table 2. We do not find any systematic change with resolution in the values of halo parameters both on the virial radius scale and around the maximum of the circular velocity ($`r=(3040)h^1\text{kpc}`$).
Left panel in Figure 6 shows the central region of the halo A<sub>1</sub> (see Table 2). This plot is similar to the Fig.1a in Moore et al. (1998) in that all profiles are drawn to the formal force resolution. The straight lines indicate slopes of two power-laws: $`\gamma =1`$ and $`\gamma =1.4`$. The figure indeed shows that at around 1% of the virial radius the slope is steeper than -1 and the central slope increases as we increase the mass resolution. Moore et al. (1998) interpreted this behavior as evidence that profiles are steeper than predicted by the NFW profile. We also note that the results of our highest resolution run $`A_1`$ are qualitatively consistent with results of Kravtsov et al. (1998). Indeed, if the profiles are considered down to the scale of two formal resolutions, the density profile slope in the very central part of the profile $`r0.01r_{\mathrm{vir}}`$ is close to $`\gamma =0.5`$.
The profiles in Figure 6 reflect the density distribution in the cores of simulated halos. However, the interpretation of these profiles is not straightforward because it requires assessment of numerical effects. The formal resolution usually does not even correspond to the scale where the numerical force is fully Newtonian (usually it is still considerably “softer” than the Newtonian value). In the ART code, the interparticle force reaches (on average) the Newtonian value at about two formal resolutions(see Kravtsov et al. 1997). The effects of force resolution can be studied by resimulating the same objects with higher force resolution and comparing the density profiles. Such convergence study was done in Kravtsov et al. (1998) where it was found that for a fixed mass resolution halo density profiles converge at scales above two formal resolutions. Second, local dynamical time for particles moving in the core of a halo is very short. For example, particles on the circular orbit of the radius $`1h^1\text{kpc}`$ from the center of halo $`A`$ makes about 200 revolutions over the Hubble time. Therefore, if the time step is insufficiently small, numerical errors in these regions will tend to grow especially fast. The third possible source of numerical errors is mass resolution. Poor mass resolution in simulations with good force resolution may, for example, lead to two-body effects (e.g., Knebe et al. 2000). Insufficient number of particles may also result in “grainy” potential in halo cores and thereby affect accuracy of orbit integration. In these effects, the mass resolution may be closely inter-related with force resolution.
It is clear thus that in order to make conclusions not affected by numerical errors, one has to determine the range of trustworthy scales using convergence analysis. Right panel in Figure 6 shows that for the halo A simulations the convergence for vastly different mass and force resolution is reached for scales $`4`$ formal force resolutions (all profiles in this figure are plotted down to the radius of 4 formal force resolutions). For all resolutions, there are more than 200 particles within the radius of four resolutions from the halo center. For the highest resolution simulation (halo A<sub>1</sub>) the convergence is reached at scales $`0.005r_{\mathrm{vir}}`$.
In order to judge which profile provides a better description of the simulated profiles we fitted the NFW and Moore et al. analytic profiles. Figure 7 presents results of the fits and shows that both profiles fit the numerical profile equally well: fractional deviations of the fitted profiles from the numerical one are smaller than 20% over almost three decades in radius. It is clear thus that the fact that numerical profile has slope steeper than $`1`$ at the scale of $`0.01r_{\mathrm{vir}}`$ does not mean that good fit of the NFW profile (or even analytic profiles with shallower asymptotic slopes) cannot be obtained.
There is certainly a certain degree of degeneracy in fitting various analytic profile to numerical results. Figure 8 illustrates this further by showing results of fitting profiles (solid lines) of the form $`\rho (r)(r/r_0)^\gamma [1+(r/r_0)^\alpha ]^{(\beta \alpha )/\gamma }`$ to the same (halo $`A_1`$) simulated halo profile shown as solid circles. The legend in each panel indicates the corresponding values of $`\alpha `$, $`\beta `$, and $`\gamma `$ of the fit; the digit in parenthesis indicates whether the parameter was kept fixed ($`0`$) or not ($`1`$) during the fit. The right two panels show fits of the NFW and Moore et al. profiles; the bottom left panel shows fit of the profiles used by Jing & Suto (2000). The top left panel shows a fit in which the inner slope was fixed but $`\alpha `$ and $`\beta `$ were fit. The figure shows that all four analytic profiles can provide a nice fit to the numerical profile in the whole range $`0.0051r_{\mathrm{vir}}`$.
### 4.2 Halo profiles at $`z=1`$
As we mentioned in § 4.1, the halo A analyzed in the previous section is somewhat special because it was selected as an isolated relaxed halo. In order to reach unbiased conclusions, in this section we will present analysis of halos from the second set of simulations (halos B, C, and D in Table 2) which were not selected to be relaxed or isolated. Based on the results of convergence study presented in the previous section, we will consider profiles of these halos only at scales above four formal resolutions use results starting only from 4 formal resolutions and not less than 200 particles. Note that these conditions are probably more stringent than necessary because these halos were simulated with $`57`$ times more particles per halo. There is an advantage in analyzing halos at a relatively high redshift. Halos of a given mass will have lower concentration (see Bullock et al. 2000). Lower concentration implies a large scale at which the asymptotic inner slope is reached. Profiles of the high redshift halos should therefore be more useful in discriminating between the analytic models with different inner slopes.
We found that substantial substructure is present inside the virial radius in all three halos. Figure 9 shows profiles of these halos at $`z=1`$. There profiles are not as smooth as that of halo A<sub>1</sub> due to the substructure. Note that bumps and depressions visible in the profiles cannot are significantly larger amplitude than the shot noise. Halo $`C`$ appeared to be the most relaxed of the three halos. This halo had the last major merger somewhat earlier than the other two. Halo $`D`$ had a major merger event at $`z2`$. Remnant of the merger are still visible as a hump at radii around $`100h^1\text{kpc}`$. The non-uniformities of profiles cause by substructure may substantially bias analytic fits to the entire range of scales below the virial radius. Therefore, we used only the central, presumably more relaxed, regions in the analytic fits: $`r<50h^1\text{kpc}`$ for halo D and $`r<100h^1\text{kpc}`$ for halos B and C (fits using only central $`50h^1\text{kpc}`$ did not change results).
The best fit parameters were obtained by minimizing the maximum fractional deviation of the fit: $`\mathrm{max}(\mathrm{abs}(\mathrm{log}\rho _{\mathrm{fit}})\mathrm{log}\rho _{\mathrm{halo}})`$. Minimizing the sum of squares of deviations ($`\chi ^2`$), as is often done, can result in larger errors at small radii with the false impression that the fit fails because it has a wrong central slope. The fit that minimizes maximum deviations, improves the NFW fit for points in the range of radii $`(520)h^1\text{kpc}`$, where the NFW fit would appear to be below the data points if the fit was done by the $`\chi ^2`$ minimization. This improvement comes at the expense of few points around $`1h^1\text{kpc}`$. For example, if we fit halo B by using $`\chi ^2`$ minimization, the concentration decreases from 12.3 (see Table 2) to 11.8 We also made a fit for halo B assuming even more stringent limits on the effects of numerical resolution. By minimizing the maximum deviation we fitted the halo starting at six times the formal resolution. Inside this radius there were about 900 particles. Resulting parameters of the fit were close to those in Table 2: $`C_{\mathrm{NFW}}=11.8`$, and maximum error of the NFW fit was 17%.
We found that the errors in the Moore et al. fits were systematically smaller than those of the NFW fits, though the differences were not dramatic. Moore et al. fit failed for halo $`D`$. It formally gave very small errors, but this was done for a fit with unreasonably small concentration $`C=2`$. When we constrained the approximation to have about twice larger concentration as compared with the best NFW fit, we were able to obtain a reasonable fit (this fit is shown in Figure 9). Nevertheless, the central part was fit poorly in this case.
Our analysis therefore failed to determine which analytic profile provides a better description of the density distribution in simulated halos. Despite the larger number of particles per halo and lower concentrations of halos, results are still inconclusive. Moore at al. profile is a better fit to the profile of halo C; the NFW profile is a better fit to the central part of halo D. Halo B represents an intermediate case where both profiles provide equally good fits (similar to the analysis of halo A).
Note that there seem to be real deviations in parameters of halos of the same mass. Halos B and D have the same virial radii and nearly the same circular velocities, yet their concentrations are different by 30%. We find the same differences in estimates of $`C_{1/5}`$ concentrations, which do not depend on specifics of an analytic fit. The central slope at around $`1\text{kpc}`$ also changes from halo to halo.
## 5 Summary
We run a series of simulations with vastly different mass and force resolution with the goal of studying the shape of the dark matter halo profile in central parts of galaxy-size halos. We use a modified version of the ART code, which is capable of handling particles with different masses, variable force, and time resolution. In runs with the highest resolution, the code achieved (formal) dynamical range of $`2^{17}=131,072`$ with 500,000 steps for particles at highest level of resolution.
Our conclusions regarding the convergence of the profiles differ from those of Moore et al. (1997). If we take into account only radii, at which we believe numerical effects (the force resolution, the resolution of initial perturbations, and the two-body scattering) are small, then we find that the slope and the amplitude of the density do not change when we change the force and mass resolution. This result is consistent with what was found in simulations of the “Santa Barbara” cluster (Frenk et al., 1999): at a fixed resolved scale results do not change as the resolution increases. For the ART code the results converged at 4 times the formal force resolution and more than 200 particles. These limits of convergence very likely depend on particular code used and on the duration of integration.
We reproduce results of Moore et al. regarding the convergence and results of Kravtsov et al. (1998) regarding shallow central profiles, but only when we consider points inside unresolved scales. We conclude that those results followed from overly optimistic interpretation of numerical accuracy of simulations.
For the galaxy-size halos considered in this paper with masses $`M_{\mathrm{vir}}=7\times 10^{11}h^1\text{M}\text{}2\times 10^{12}h^1\text{M}\text{}`$ and concentrations $`C=917`$ both the NFW profile $`\rho r^1(1+r)^2`$ and the Moore et al. profile $`\rho r^{1.5}(1+r^{1.5})^1`$ give good fits with accuracy about 10% for radii not smaller than 1% of the virial radius. None of the profiles is significantly better then the other.
Halos with the same mass may have different profiles. No matter what profile is used – NFW or Moore et al. – there is no universal profile: just halo mass does not yet define the density profile. Nevertheless, the universal profile is extremely useful notion which should be interpreted as the general trend $`C(M)`$ of halos with larger mass to have lower concentration. Deviations from the general $`C(M)`$ are real and significant (Bullock et al., 2000). It is not yet clear, but seems very likely that the central slopes of halos also have real fluctuations. The fluctuations in the concentration and the central slopes are important for interpretation of the central parts of rotation curves.
We acknowledge the support of the grants NAG- 5- 3842 and NST- 9802787. A.V.K. acknowledges support by NASA through Hubble Fellowship grant HF-01121.01-99A from the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5-26555. Computer simulations presented in this paper were done at the National Center for Supercomputing Applications (NCSA), Urbana-Champaign, Illinois.
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# Variational wave functions of a vortex in cyclotron motion
## Acknowledgments
It is a pleasure to thank David Thouless for initiating the idea of this work and stimulating discussions. The author is also grateful to Ping Ao for valuable comments and suggestions. This work was partially supported by NSF DMR-9815932.
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# Mid-Infrared Observations of Normal Star-Forming Galaxies: The Infrared Space Observatory Key Project Sample1footnote 11footnote 1This paper is based on observations with the Infrared Space Observatory (ISO). ISO is an ESA project with instruments funded by ESA member states (especially the PI countries: France, Germany, The Netherlands and the United Kingdon) and with the participation of ISAS and NASA.
## 1 Introduction
Normal galaxies, defined as those with on-going star formation, account for most of the luminous mass in the local Universe. Their luminosity is derived from stars and they span a broad range of observed morphologies, luminosities, and infrared-to-blue ratios. Results from IRAS have shown that the infrared colors from the four IRAS bands are sensitive indices of the radiation intensity in the interstellar medium. The mid-infrared (5-20 $`\mu `$m) emission is dominated by very small grains fluctuating to high temperatures and polycyclic aromatic hydrocarbons (PAHs) (Helou et al. 2000). These grains are not in thermal equilibrium but they still convert heating photons and thereby trace star formation. Larger grains in thermal equilibrium dominate the emission from normal galaxies at longer wavelengths. By comparing the mid-infrared emission to other components of the galaxy such as H I, H<sub>2</sub>, ionized gas, and starlight, one can derive the physical properties of the interstellar gas, dust and radiation field in galaxies (e.g. Vigroux et al. 1999).
The Infrared Space Observatory (ISO) U.S. Key Project on Normal Galaxies (PI: G. Helou, proposal id: SF\_GLX\_\*, hereafter the “Key Project”) was proposed to study the interstellar medium of a broad range of normal galaxies using three of the four instruments aboard ISO: ISOCAM, ISOLWS, and ISOPHOT. Under NASA guaranteed time, the U.S. Key Project obtained ISO observations of 69 galaxies. Nine relatively nearby and extended galaxies were chosen to provide spatially resolved cases so that various phases of the interstellar medium could be studied independently. The remaining 60 galaxies in the Key Project sample cover the full range of observed morphologies, luminosities, and IRAS colors seen in normal galaxies. Using this diverse sample, we hope to gain new insight into the star formation process on the scale of galaxies, especially its drivers and inhibitors (Helou et al. 1996).
## 2 The U.S. Key Project ISOCAM Sample
The ISOCAM sample of galaxies for the Key Project is presented in Table 1 where we list each galaxy’s position, optical size, morphology, recession velocity, distance estimate, and far-infrared-to-blue ratio. The positions, optical sizes, and velocities come from the NASA Extragalactic Database<sup>10</sup><sup>10</sup>10The NASA/IPAC Extragalactic Database is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration. and were last updated February 2000. The distances were computed in the Local Group reference frame assuming a Hubble constant of 75 km s<sup>-1</sup> Mpc<sup>-1</sup>, except for IC 10 for which we rely on results from recent Cepheid variable observations (Wilson et al. 1996; Saha et al. 1996), and for NGC 1569 in which case a distance has been computed from Hubble Space Telescope observations of resolved supergiants (O’Connell, Gallagher & Hunter 1994). The mean distance of the sample is 34 Mpc. The RC3 (de Vaucouleurs et al. 1991) optical morphological classification for each galaxy was kindly re-examined by H. Corwin using the POSS plates and more recent CCD observations by the Key Project team and those found in the literature; we have updated a total of 17 RC3 morphologies. Figure 1 shows the histogram of optical morphological types. The morphologies are lenticular or later except for one elliptical, NGC 6958. There is an approximately even distribution of spirals and a large number of irregulars in the final sample.
In all, five of the nine relatively extended Key Project galaxies and 56 of the 60 smaller Key Project galaxies were mapped with ISOCAM. These galaxies span a broad range in physical properties such as total luminosity, dust temperature, and star formation activity. Figure 2 is a color-color diagram of ratios of IRAS flux densities for the sample (see Table 4). Quiescent galaxies such as NGC 7418 lie along the upper left of the distribution while more active star forming galaxies such as NGC 1569 are found to the lower right.
## 3 The ISOCAM Observations
Observations were made using the 32$`\times `$32 pixel Si:Ga long wavelength (LW) array of ISOCAM (Césarsky et al. 1996) aboard ISO (Kessler et al. 1996). We utilized the narrow-band LW1 filter (4.5 $`\mu `$m, $`\delta \lambda =1.0`$ $`\mu `$m), the broad-band LW2 filter (6.75 $`\mu `$m, $`\delta \lambda =3.5`$ $`\mu `$m), and broad-band LW3 filter (15.0 $`\mu `$m, $`\delta \lambda =6.0`$ $`\mu `$m) to sample the mid-infrared emission from the galaxies. The mid-infrared regime represents a transition from stellar to interstellar emission, though the former should be negligible at 6.75 $`\mu `$m for the majority of our sample. From ISOPHOT observations (Helou et al. 2000) we know there are strong aromatic emission features at 6.2, 7.7 and 8.6 $`\mu `$m, all of which fall within the 6.75 $`\mu `$m filter. The 15 $`\mu `$m filter contains some emission from aromatics at 11.3 and 12.7 $`\mu `$m, and continuum emission from very small grains. We note that strong \[Ne II\] (12.8 $`\mu `$m) and \[Ne III\] (15.6 $`\mu `$m) fine structure lines have been observed in the spectra of H II regions (e.g. M17; Césarsky et al. 1996). Rotational H<sub>2</sub> lines are also seen in the LW3 bandpass (Kunze et al. 1996, Valentijn et al. 1996, Timmerman et al. 1996), and while these lines are important physically, their contribution to the total flux is very small (Helou et al. 2000).
All of our ISOCAM observations used a gain setting of 2 and a 6″$`\times `$6″ pixel scale. Tables 2 and 3 list the date, revolution (orbit number during the ISO mission), and exposure times for each galaxy. For six galaxies (NGC 1313, NGC 3620, NGC 5713, IC 4595, NGC 6753, NGC 7771), the 4.5 $`\mu `$m observations were not done at the same time as the 6.75 and 15 $`\mu `$m observations, so the 4.5 $`\mu `$m maps for these galaxies are at a different roll angle compared to the 6.75 and 15 $`\mu `$m maps.
The bulk of our ISOCAM observations are in the form of 2$`\times `$2 overlapping raster maps (Table 2) with a raster step of 75″, or 12.5 pixel widths, executed along the axes of the array. Thus the final maps are 4$`\stackrel{}{\mathrm{.}}`$45$`\times `$4$`\stackrel{}{\mathrm{.}}`$45, and the half-integer pixel raster spacing yields an inner 1$`\stackrel{}{\mathrm{.}}`$95$`\times `$1$`\stackrel{}{\mathrm{.}}`$95 portion that is sampled every 3″; the final maps are rebinned onto a 3″ pixel<sup>-1</sup> grid. Five galaxies have larger raster maps (Table 3), with up to 8$`\times `$8 pointings and a step size of 81″, or 13.5 pixel widths. These larger maps are also rebinned onto a 3″ pixel<sup>-1</sup> grid.
## 4 Data Reduction
The ISOCAM observations were reduced using the CAM Interactive Analysis (CIA) software package (Delaney 1998) using Pipeline 7 datasets. The reduction steps include dark subtraction, deglitching, transient removal, and flat-field correction. We experimented with various methods for each step, with the goal of finding a reduction method that worked well for the entire galaxy sample.
The arrival of the first wave of ISOCAM observations in mid-1996 allowed us to explore in some detail various reduction and analysis techniques. This early work indicated that the library dark distributed with each dataset did an acceptable job in removing the effects of the dark current. The default deglitcher, which uses the multi-resolution median transform method (Starck et al. 1999), did the best job in removing the transitory cosmic ray and charged particle hits without simultaneously flagging the source data as bad data. Library sky flats from some revolutions have cosmic-ray/charged particle residuals, so we investigated how to use the observations themselves to create a target field sky flat. We found that as long as the galaxy is not too extended or too bright, we could create a good sky flat from the target field data. The only remaining task was to remove the effects of the finite response time, which can artificially increase or decrease flux detections, referred to below as “transient removal.”
All transient removal techniques that existed within CIA circa 1997 were tested to identify the one best suited for our type of observations. In particular, we studied in detail the relative merits of the ‘fit3,’ ‘inverse,’ and ‘fouks-schubert’ methods. All other options—dark, deglitcher, target field sky flat—were unchanged in the tests. Individual pixels from both on and off-source observations were examined before and after transient removal, and the final maps were inspected. The ‘fit3’ (Siebenmorgen et al. 1997) method appeared to do the best job but it leaves the first raster position, the one most affected by the transient, at a slightly different flux level than the remaining raster positions. This is because ‘fit3’ removes the transient effect independently for each raster position. Thus, the average map from each raster position will have a slightly different sky value, with the first raster position being more discrepant (usually by more than 5%) compared to the remaining raster positions. To solve this problem, we experimented with masking out the first raster position. Several of the transient removal methods were re-tested, and again the ‘fit3’ method appeared to work best in terms of removing the transient from individual pixels, both on and off-source, and producing a “good” final map.
In summary, the data reduction steps we use were as follows. First, starting with the Standard Processed Data (SPD) of the pipeline processing, a library dark is subtracted from each frame. Next, each frame is deglitched using the multi-resolution median transform method. Then all of the frames in the first raster position are masked out (not used to produce the final map). Next, transients are removed using the Saclay ‘fit3’ model. The frames from each raster position are averaged together, and then these average images are used to make a target field sky flat. For the 4.5 $`\mu `$m maps the library sky flat was used, as the sky flux in these maps is too low to produce good target field sky flats from the data. Next, the flat-fielded average images are mosaicked to create the final map. The final step converts the pixel units from ADU/gain/second to mJy.
## 5 The ISOCAM Images
### 5.1 The Small Maps
Figure 3 presents optical Digitized Sky Survey (DSS) images and ISOCAM maps for the 2$`\times `$2 raster sample of galaxies, arranged one row per galaxy. Figures 4-8 present the larger raster maps, also including the optical DSS images for comparison. As a consequence of the ‘dead’ column #24 in the ISOCAM LW array, all of the maps contain two zero value columns towards their right edge. There is also a notch in the lower left corner of the ISOCAM maps, due to our masking of the first raster position. In many cases the maps still show low level residuals due to transients from the sources. In particular, residuals just below the source are common and result from the source being in that position during the third raster observation. There are also weak residuals in some maps from the second raster position (the upper left in the ISOCAM maps). These residuals are due to our scheme of masking the first raster position, as they are very weak or non-existent in maps made with all four raster positions. Such low level residuals are masked out before source fluxes are measured.
NGC 4490 and UGC 2855 are relatively large galaxies compared to the 4$`\stackrel{}{\mathrm{.}}`$45$`\times `$4$`\stackrel{}{\mathrm{.}}`$45 raster maps; for both galaxies the amount of sky area in the final maps is insufficient. We tried using sky flats from the observations themselves and the library sky flats. The results using the library sky flats are superior to the target field sky flats, leaving very weak transient residuals and a much smoother sky. The library sky flats do not appear to add any artifacts to the final maps (as was the case for a few of our galaxies during the testing of reduction methods), and so for these two galaxies fluxes are determined from ISOCAM maps reduced using library sky flats. We note that the ISOCAM map of NGC 4490 does not contain the entire galaxy; our 6.75 and 15 $`\mu `$m fluxes should be considered lower limits for this galaxy.
We successfully detected all galaxies at all mid-infrared wavelengths attempted in our ISOCAM program, (though just barely for two nearby galaxies; see §5.2.3). Overall the ISOCAM maps show many of the same features present in the optical images; prominent spiral arms or rings, and distinct clumps and regions of relatively high mid-infrared emission are quite evident in some of the later type galaxies such as IC 4595, NGC 1385, NGC 3705, NGC 4713, NGC 7418, and UGC 2855. However, in most cases the bulges seen in the optical are not apparent in the mid-infrared maps. More importantly, the mid-infrared morphologies at 4.5, 6.75, and 15 $`\mu `$m are qualitatively similar. There are no obvious differences in how the light is distributed, though upon closer inspection of the central regions of some galaxies there are hints of variations with wavelength. A quantitative look at the relative distribution of mid-infrared emission at 6.75 and 15 $`\mu `$m will be presented in Section 7.1. The typical sky background levels are 1, 5, and 20 MJy sr<sup>-1</sup>, with rms fluctuations at the 0.1, 0.1, and 0.2 MJy sr<sup>-1</sup> levels, at 4.5, 6.75, and 15 $`\mu `$m, respectively.
### 5.2 The Larger Maps
#### 5.2.1 IC 10
IC 10 was observed twice. In the first set of observations the center of the galaxy was chosen to be the center of the map, as with the maps for all our other galaxies. However, we noticed that the mid-infrared emission from IC 10 extended more to the west of center, leaving some interesting structure on the edges of the maps. A second set of observations was taken a year later to more properly map the mid-infrared emission. The final set of 6.75 and 15 $`\mu `$m maps results from a co-add of the maps from the two epochs. The 11.4 $`\mu `$m (LW8 filter) and 4.5 $`\mu `$m maps were obtained only in the second epoch.
In the mid-infrared IC 10 appears to consist of a collection of bright knots, except at 4.5 $`\mu `$m where the galaxy appears to be quite faint. IC 10 is classified as a barred irregular galaxy but the bar is not apparent in the mid-infrared maps. There are two wispy arms extending to the north which contain distinct knots of stronger mid-infrared emission. An overall faint and diffuse mid-infrared emission is detected towards the west and south of the galaxy center, extending almost to the edge of our maps. A surface brightness analysis of the 6.75 and 15 $`\mu `$m ISOCAM observations of IC 10 can be found in Dale et al. (1999).
#### 5.2.2 NGC 1313
Observations for this galaxy were taken in two parts, a western piece and an eastern piece, with the two pieces mosaicked together to make a final map. The 4.5 $`\mu `$m observation was taken 5 months later at a different roll angle. In the mid-infrared, the bar of the galaxy appears as discrete emission sources. The spiral arms are prominent as both diffuse emission and knots of strong emission. Similar to what is seen optically, there is also a faint wisp of emission directly to the south of galaxy center. The bright knots of optical emission in the eastern arm also show strongly in the mid-infrared. Our ISOCAM map is unfortunately too small to cover the large star-forming complex just to the southwest of the galaxy (Ryder et al. 1995). The galaxy is barely detected at 4.5 $`\mu `$m. A surface brightness analysis of the 6.75 and 15 $`\mu `$m observations of NGC 1313 can be found in Dale et al. (1999).
#### 5.2.3 NGC 2366 and NGC 6822
NGC 2366 and NGC 6822 are low luminosity irregulars that contain several moderate to supergiant-sized H II regions (Hunter et al. 2000). The mid-infrared maps for these galaxies are of insufficient depth: the 1$`\sigma `$ rms fluctuations in the sky level for NGC 2366 are 0.05 and 0.14 MJy sr<sup>-1</sup> at 6.75 and 15 $`\mu `$m, respectively, while the numbers for NGC 6822 are 0.09 and 0.28 MJy sr<sup>-1</sup>. Very little 6.75 or 15 $`\mu `$m emission is detected from these two galaxies, suggesting they contain only a small amount of dust. In NGC 2366 most of the mid-infrared emission coincides with the supergiant H II region NGC 2363 in the southwestern part of the galaxy. There is a much fainter region of emission to the northeast of that which does not obviously correspond to anything in particular in the optical. In NGC 6822 the mid-infrared emissison is found in four H II regions. There is a bit of diffuse emission towards the center of the galaxy, but it is of too low signal-to-noise to measure with any degree of confidence. The fluxes given in Table 4 are the combined fluxes of the H II regions; we caution that these are only estimates of the total flux, as the true sizes of the galaxies in the mid-infrared are difficult to determine. There are a few foreground stars seen in the 6.75 $`\mu `$m maps.
#### 5.2.4 NGC 6946
At optical wavelengths, NGC 6946 is a striking example of a nearly face-on ($`i`$$``$30) spiral galaxy. The galaxy is equally rich in detail in the mid-infrared. The spiral arms can easily be traced in the mid-infrared maps, and each contains several knots of brighter emission. The ‘sharp’ edge of the northeastern-most arm is real and not an artifact of the mosaicking. It is also seen in IRAS HiReS maps of NGC 6946. Preliminary work on NGC 6946 from these ISOCAM observations (using earlier data reduction techniques) by Helou et al. (1996) indicated that there is very little mid-infrared color variation throughout the disk of the galaxy. This result is echoed by the work of Tuffs et al. (1996), who used ISOPHOT data to suggest that the bulk of the far-infrared luminosity is due to a uniformly colored and rather cold diffuse emission from the disk. Work by Malhotra et al. (1996), again using these same ISOCAM observations reduced with prior techniques, showed that the mid-infrared emission is consistent with an exponential disk of scale length of 75″, based on a flux profile computed from the median flux within annuli between 70 and 200″. They also indicate that the mid-infrared arm-interarm contrast, being close to that observed in $`R`$-band but lower than observed in H$`\alpha `$, suggests that non-ionizing radiation plays an important role in the heating of the dust. A surface brightness analysis of the ISOCAM observations for NGC 6946, using our new reduction methods, can be found in Dale et al. (1999). We stress that the main results do not change: the color variations analyzed in Dale et al. are quite close to those presented earlier. In fact, the small discrepancies ($``$10%) between the two analyses can be taken as a first-order estimate of the uncertainties in the data. Moreover, using the newly-processed data we find a similar mid-infrared disk scale length (72″). Finally, we point out that the starbursting center is quite bright in the mid-infrared; our integration times nearly saturated two pixels at 15 $`\mu `$m and saturated as many as four pixels at 6.75 $`\mu `$m. The impact of these saturated and nearly saturated pixels on the overall flux density is measurably small, approximately a 15% (25%) effect at 15 (6.75) $`\mu `$m for affected pixels (the method used to approximate this effect is described in Section 6.1). Our values for the total flux density should be formally taken as lower limits, though the underestimation is no more than a few percent.
## 6 Flux Density Estimates
Before beginning any quantitative analysis on the ISOCAM images presented in Figures 3-8, we cleaned the images by masking any foreground stars and cosmic ray/transient residuals. To allow direct comparison between the various filters we smoothed the maps to a similar resolution. To do this we needed to first estimate the approximate resolution for the 6.75 and 15 $`\mu `$m maps. There are a few maps which have strong emission from a foreground star (e.g. NGC 1569 and NGC 6946), and several galaxies appear to be unresolved at these mid-infrared wavelengths (IRAS F10565, NGC 4418, IC 860, IC 883, CGCG1510.8+0725, IRAS F23365+3604, Markarian 331). Analysis of these images indicate a (FWHM) resolution of $``$7″ and $``$8″ at 6.75 at 15 $`\mu `$m, respectively. This information was used to convolve the 6.75 and 15 $`\mu `$m maps, using slightly different Gaussian smoothing profiles for each filter, to generate maps with approximately uniform resolution (9″ FWHM). As a final step the images were registered to the coordinate system of the 15 $`\mu `$m map by assuming the peak emission at 6.75 $`\mu `$m is spatially coincident with the peak emission at 15 $`\mu `$m.
### 6.1 Mid-Infrared Flux Densities
Multi-aperture photometry is our preferred method for extracting mid-infrared flux densities. For each galaxy care is taken to define the apertures and background annulus using the flux density distribution from both the 6.75 and 15 $`\mu `$m (and when available the 4.5 $`\mu `$m) maps. In each case, a series of evenly spaced, concentric apertures are centered on the central emission peak (or as in the case of NGC 1569 which has two central emission peaks, the brightest of the central emission peaks), with the outermost aperture sized according to the greatest extent of the emission in either map. The background annulus was then defined starting from the outermost aperture, unless a localized contamination (e.g. star) warrants starting the background annulus at a slightly larger radius; the typical background annulus has a width of 15″. The mean value within the background annulus is taken to be the sky background level. As described in Section 5.1, the typical rms fluctuation in the background level is 0.1–0.2 MJy sr<sup>-1</sup>. The background-subtracted flux density within the outermost aperture is the total flux density.
The mid-infrared flux densities and their uncertainties for the galaxies are given in Table 4 along with the IRAS flux densities. Column 1 in Table 4 lists the galaxy. Columns 2–5 list the IRAS flux densities in Janskys and columns 6–7 list the ISO 6.75 and 15 $`\mu `$m flux densities and their uncertainties in Jy (see §6.2 for a comparison of ISOCAM and IRAS 12 $`\mu `$m fluxes). We assume a 20% uncertainity in the flux density calibration of ISOCAM data (Blommaert & Cesarsky 1998; Biviano 1998). By far, this uncertainty level dominates any other uncertainties that are introduced by our method of flux density extraction (e.g. choice of aperture and background annulus, sky rms fluctuation, etc.). However, our finite aperture measurements systematically underestimate the total flux density of the galaxies by a small amount due to the limiting $``$0.1 MJy sr<sup>-1</sup> sensitivity of our observations. We can estimate the amount of low surface brightness emission that is missed by extrapolating the observed isophotal surface brightness trends. Assuming such trends hold out to a surface brightness of 0.001 MJy sr<sup>-1</sup> (about two disk scale-lengths), no more than a few percent of the flux is missing. Overall, for most galaxies we thus find a total flux density uncertainty of 20%; this number is slightly higher ($``$25%) for about one fourth of the sample due to more significant sky background fluctuations.
The uncertainty is slightly larger in cases where i) the mid-infrared emission extends beyond the observed target field, or ii) there is saturation, two additional effects that systematically lower the flux density measure. Only the fluxes for UGC 2855 and NGC 4490 suffer from the former effect; extrapolations of the observed profiles, as described above, show we are missing $``$ 5% of such flux for these two galaxies. A second concern is that observations of a few galaxies saturated the analog to digital converter for a small number of pixels, the effect being more serious for the longer elementary integrations for the 6.75 $`\mu `$m bandpass. We have gauged the amplitude of this effect using the known response function of the ISOCAM detectors: the actual flux density should be larger than the initial frame’s measure by a factor 1/0.6 (ISOCAM Observer’s Manual). For most of the 11 galaxies exhibiting saturated pixels the net effect is an underestimate of the total flux density by 6% or less; the 6.75 and 15 $`\mu `$m flux densities for NGC 4418 are underestimated by about 15% and 25%, respectively, comparable to the calibration uncertainty level. The flux densities computed for these galaxies are lower limits and are indicated by a colon in Table 4; they may be assumed to be uncertain at the 25% level.
Table 5 lists the total 4.5 $`\mu `$m flux densities for 13 galaxies in our sample. For NGC 1569 we obtained 2$`\times `$2 ISOCAM maps in eight filters; NGC 1569 flux densities and uncertainties are given in Table 6. The typical LW1 flux density uncertainty is 25-30% due to a lower overall signal to noise compared to the LW2 and LW3 maps.
### 6.2 Calibration Check: The Consistency of ISO and IRAS Fluxes
Lu et al. (2000) show that the agreement between ISOPHOT and ISOCAM LW2 fluxes is better than 20%. An ISO-independent check of the calibration of our ISOCAM maps can be provided by IRAS data. Figure 9 presents the ratio of ISOCAM fluxes to IRAS 12 $`\mu `$m fluxes versus $`\frac{I_\nu (60\mu \mathrm{m})}{I_\nu (100\mu \mathrm{m})}`$ color. The dashed lines indicate the expected ratios for cirrus (log $`\frac{I_\nu (6.75\mu \mathrm{m})}{I_\nu (12\mu \mathrm{m})}`$ = $`0.29`$ and log $`\frac{I_\nu (15\mu \mathrm{m})}{I_\nu (12\mu \mathrm{m})}`$ = $`0.01`$ for $`\rho `$ Ophiuchi; W. Reach, private communication). The expected cirrus color agrees well with our observed $`\frac{I_\nu (6.75\mu \mathrm{m})}{I_\nu (12\mu \mathrm{m})}`$ color ratio: log $`\frac{I_\nu (6.75\mu \mathrm{m})}{I_\nu (12\mu \mathrm{m})}`$\[$`0.5<`$ log $`\frac{I_\nu (60\mu \mathrm{m})}{I_\nu (100\mu \mathrm{m})}`$$`<0.2`$\]=$`0.26\pm 0.11`$ (the uncertainty represents the population dispersion excluding NGC 1155 and NGC 6958). The agreement is even better for the more quiescent galaxies, even though real galaxies have other emission (star-formation regions) in them in addition to cirrus. Moreover, we expect the $`\frac{I_\nu (6.75\mu \mathrm{m})}{I_\nu (12\mu \mathrm{m})}`$ color to decrease for warmer galaxies (relatively high $`\frac{I_\nu (60\mu \mathrm{m})}{I_\nu (100\mu \mathrm{m})}`$ color) since the PAH features in the 6.75 $`\mu `$m filter become relatively weak compared to the rising continuum in the 12 $`\mu `$m band, and this is indeed observed. Photospheres can contribute to the flux at 6.75 $`\mu `$m and throw off the ratio, especially for measurements of elliptical galaxies. Since this effect does not seem significant in this sample, we conclude that in most of our sample the 6.75 $`\mu `$m and longer wavelength bands are dominated by interstellar emission.
For the $`\frac{I_\nu (15\mu \mathrm{m})}{I_\nu (12\mu \mathrm{m})}`$ color ratio comparison with cirrus, the agreement is not as good, but the ratio observed for galaxies is lower than that of cirrus only by 35%: log $`\frac{I_\nu (15\mu \mathrm{m})}{I_\nu (12\mu \mathrm{m})}`$\[$`0.5<`$ log $`\frac{I_\nu (60\mu \mathrm{m})}{I_\nu (100\mu \mathrm{m})}`$$`<0.2`$\]=$`0.14\pm 0.10`$ (excluding NGC 6958). Surprisingly, the offset goes in the wrong direction for astrophysical effects to resolve the discrepancy: the 15 $`\mu `$m band should pick up more emission in galaxies than the cirrus spectrum represents, and we reach the “cirrus ratio” only for mildly active galaxies. There are a few possible explanations for this discrepancy, most importantly the uncertainty in the cirrus spectrum (calibration uncertainty and spatial variations in the cirrus behavior), the uncertainties in the various filter band passes, and incompleteness in the ISO flux integral. The ISO 15 $`\mu `$m maps are integrated only down to the $``$0.2 MJy sr<sup>-1</sup> isophote, whereas the large angular size of the IRAS detectors tended to pick up contributions from extended diffuse emission. However, as pointed out in Section 6.1, extrapolations of the observed surface brightness profiles down to 0.001 MJy sr<sup>-1</sup> show we are missing at most 5% of the total ISO flux. Moreover, if the problem is ISO flux incompleteness, we might expect a similar discrepancy in the $`\frac{I_\nu (6.75\mu \mathrm{m})}{I_\nu (12\mu \mathrm{m})}`$ comparison; the ISO 6.75 $`\mu `$m maps go only slightly deeper, to $``$0.1 MJy sr<sup>-1</sup>. We conclude that it is unlikely that ISO flux incompleteness is responsible.
## 7 Analysis
### 7.1 Mid-Infrared Flux Curve of Growth Profiles
Results of the aperture photometry described in Section 6.1 are presented in Figures 10 and 11. From the smoothed maps, the curve of growth profiles for the 6.75 and 15 $`\mu `$m emission are represented by the solid and dashed lines, respectively. Each profile is flux-normalized to the total flux observed in the respective filters. The bold, solid vertical line indicates the resolution scale in physical units and corresponds to the 4$`\stackrel{}{\mathrm{.}}`$5 HWHM Gaussian smoothing profile. The emission from most galaxies is very well resolved in the maps: for about half (45%) of the galaxies in our sample, 10% or less of the mid-infrared light arises from within the central ISOCAM beam; one-sixth of the galaxies emit $`>`$ 25% in the central beam. While the physical extent and concentration of the mid-infrared emission varies from galaxy to galaxy, the 6.75 and 15 $`\mu `$m profiles are remarkably similar for a given source.
Against this general similarity, some galaxies display strong central $`\frac{I_\nu (15\mu \mathrm{m})}{I_\nu (6.75\mu \mathrm{m})}`$ ratios, which can be identified in Figures 10 and 11 by a narrower 15 $`\mu `$m than 7 $`\mu `$m curve. The most notable cases are NGC 1156, NGC 1569, NGC 4519, IC 4662, NGC 6958, and NGC 2366. As discussed below (Section 7.2), larger values of $`\frac{I_\nu (15\mu \mathrm{m})}{I_\nu (6.75\mu \mathrm{m})}`$ point to more intense heating of the dust, suggesting that the galaxies with this central color enhancement harbor high densities of H II regions in their centers, or nuclear starbursts. This is indeed consistent with the properties of NGC 1156, NGC 1569, IC 4662, and NGC 2366, which are all Magellanic Irregulars with bright star-forming regions near their centers (Hunter et al. 2000; Ho et al. 1995). It is interesting to note that the reverse situation, namely a decrease in the $`\frac{I_\nu (15\mu \mathrm{m})}{I_\nu (6.75\mu \mathrm{m})}`$ color, is not observed in this sample.
### 7.2 The ISO-IRAS Color Diagram
The $`\frac{I_\nu (6.75\mu \mathrm{m})}{I_\nu (15\mu \mathrm{m})}`$ ratio has emerged as an interesting diagnostic of the radiation environment (Helou et al. 1997). As evident in the ISO-IRAS color-color diagram (Figure 12), this ratio remains relatively constant and near unity as the interstellar medium of galaxies proceeds from quiescent to mildly active, where the level of activity is indicated by a rising $`\frac{I_\nu (60\mu \mathrm{m})}{I_\nu (100\mu \mathrm{m})}`$ (Helou 1986). As dust heating increases further, the flux at 15 $`\mu `$m increases steeply compared to 7 $`\mu `$m. The data plotted in Figure 12 are consistent with an inflection in the mean trend occurring near log $`\frac{I_\nu (60\mu \mathrm{m})}{I_\nu (100\mu \mathrm{m})}`$=$`0.2`$, which we interpret as due to excess emission in the 15 $`\mu `$m band rather than a drop in the 7 $`\mu `$m band. The main argument for this interpretation is that the $`I_\nu (6.75\mu `$m)/FIR ratio does not drop as precipitously as $`\frac{I_\nu (6.75\mu \mathrm{m})}{I_\nu (15\mu \mathrm{m})}`$ for these objects. We assign to this emission a characteristic temperature 100 K $`<T_{\mathrm{MIR}}<`$ 200 K, since that is the range that would allow a blackbody to contribute considerably to the 15 $`\mu `$m band but not to the 7 $`\mu `$m band; the estimates hold for modified blackbodies as well. Such values of $`T_{\mathrm{MIR}}`$ are typical of heating intensities about $`10^4`$ times greater than the diffuse interstellar radiation field in the Solar Neighboorhood (Helou et al. 1997). While such a temperature could result from classical dust heated within or just outside H II regions, there is no decisive evidence as to the size or location of grains involved or whether in fact they are in equilibrium with the radiation field. It is simpler at this time to associate this component empirically with the observed emission spectrum of H II regions and their immediate surroundings (Tran 1998; Contursi et al. 2000). This emission has severely depressed aromatic feature or PAH emission, and is dominated by a steeply rising though not quite a blackbody continuum near 15 $`\mu `$m, consistent with mild fluctuations in grain temperatures, $`\mathrm{\Delta }T/T0.5`$. This H II region hot dust component characterized by color temperatures 100 K $`<T_{\mathrm{MIR}}<`$ 200 K becomes detectable in systems where the color temperature from the $`\frac{I_\nu (60\mu \mathrm{m})}{I_\nu (100\mu \mathrm{m})}`$ ratio is only $`T_{\mathrm{FIR}}50`$ K (Helou et al. 1988). This disparity between color temperatures derived from different wavelengths demonstrates the broad distribution of dust temperatures within any galaxy.<sup>11</sup><sup>11</sup>11The quoted far-infrared color temperatures are merely approximations arising from graybody profiles that exhibit similar flux ratios for the two wavelength bands. See Helou et al. (1988) for more details on $`\frac{I_\nu (60\mu \mathrm{m})}{I_\nu (100\mu \mathrm{m})}`$ color temperatures. The combined data from ISO and IRAS on these systems are consistent with an extension of the “two-component model” of infrared emission (Helou 1986). The low $`\frac{I_\nu (6.75\mu \mathrm{m})}{I_\nu (15\mu \mathrm{m})}`$ ratio is associated with the active component, and combines in a variable proportion with a component with a $`\frac{I_\nu (6.75\mu \mathrm{m})}{I_\nu (15\mu \mathrm{m})}`$ near unity (Dale et al. 1999).
### 7.3 Mid-Infrared vs Optical Sizes
It is interesting to measure the size of a galaxy in the mid-infrared in comparison to an equivalent measure at optical wavelengths, e.g. $`R(0.44\mu \mathrm{m})_{25}`$, the length of the semi-major axis out to 25 mag arcsec<sup>-2</sup> in $`B`$ (obtained from the RC3 catalog; de Vaucouleurs et al. 1991). Using fits of isophotal ellipses to the mid-infrared distribution, we measure semi-major axis lengths from both the 6.75 and 15 $`\mu `$m maps at several surface brightness levels between 0.2 and 1.0 MJy sr<sup>-1</sup>. We fit elliptical isophotes to the mid-infrared images using both standard and customized IRAF packages (see Haynes et al. 1999 for further details on the GALPHOT surface photometry fitting routines). We do not measure sizes at lower surface brightness levels due to S/N constraints, since as discussed in Section 5.1, the typical rms fluctuation in the sky background is $`0.10.2`$ MJy sr<sup>-1</sup>. The mid-infrared semi-major axes are corrected for ISOCAM resolution using a quadratic formula $`R=\sqrt{R_{obs}^24\stackrel{}{\mathrm{.}}5^2}`$. The open circles in Figure 13 show the mean $`R(6.75\mu \mathrm{m})`$/$`R(0.44\mu \mathrm{m})_{25}`$ and mean $`R(15\mu \mathrm{m})`$/$`R(0.44\mu \mathrm{m})_{25}`$ for the sample as a function of the mid-infrared surface brightness level at which mid-infrared sizes are measured. Error bars reflect the one sigma population dispersions in the ratios at each surface brightness, divided by the square root of the number of galaxies; histograms of the ratios at each surface brightness level are also displayed at the top of the figure. The data are reasonably well approximated by an exponential dependence of mid-infrared surface brightness on radius, and the trend holds for the entire range of measured surface brightness levels. Invoking a small extrapolation of the trend to lower surface brightness levels shows that, on average, mid-infrared sizes at 6.75 and 15 $`\mu `$m match $`R(0.44\mu \mathrm{m})_{25}`$ at a surface brightness level of $`I_\nu (6.75\mu \mathrm{m})0.04`$ MJy sr<sup>-1</sup> and $`I_\nu (15\mu \mathrm{m})0.09`$ MJy sr<sup>-1</sup> respectively. The uncertainties in these surface brightness estimates are 25%.
Since $`R(0.44\mu \mathrm{m})_{25}`$ is comparable to the location of the Sun in the Milky Way, one would expect similar values of the mid-infrared surface brightness of the Milky Way in the Local Neighborhood. For a typical high Galactic latitude H I column density of $`23\times 10^{20}`$ cm<sup>-2</sup> (Kulkarni & Heiles 1988), and for an emissivity of $`4\pi \nu I_\nu (12\mu \mathrm{m})=1.1\times 10^{31}`$ W per H atom (Boulanger & Pérault 1988), the expected mid-infrared brightness would be $`I_\nu (12\mu \mathrm{m})0.08`$ MJy sr<sup>-1</sup>, indeed comparable to the values above. In addition, one can translate the 12 $`\mu `$m band flux density into expected 6.75 and 15 $`\mu `$m flux densities using the cirrus spectrum of Reach and Boulanger (1998). One finds $`I_\nu (6.75\mu \mathrm{m})0.04`$ MJy sr<sup>-1</sup> and $`I_\nu (15\mu \mathrm{m})0.09`$ MJy sr<sup>-1</sup>, the same numbers we find for the Key Project sample. In short, the typical mid-infrared emission from the outskirts of normal galaxies is reassuringly similar to what we see in the quiescent regions of the Milky Way.
We have compared galaxy mid-infrared sizes to their optical sizes, in particular those defined at the 25 $`B`$-mag arcsec<sup>-2</sup> level. Since this surface brightness corresponds to 0.018 MJy sr<sup>-1</sup> (assuming $`f_B=4260\times 10^{0.4m_B}`$ Jy), the observed mid-infrared to $`B`$-band flux ratios at $`R(0.44\mu \mathrm{m})_{25}`$ are
$$\frac{\nu I_\nu (6.75\mu \mathrm{m})}{\nu I_\nu (0.44\mu \mathrm{m})_{25}}0.12\pm 0.03\frac{\nu I_\nu (15\mu \mathrm{m})}{\nu I_\nu (0.44\mu \mathrm{m})_{25}}0.15\pm 0.04$$
(1)
In other words, approximately eight (seven) times the flux is emitted at the center of the $`B`$-band as compared to the emission at 6.75 $`\mu `$m (15 $`\mu `$m), in the outskirts of normal galaxies. Once again we can contrast this with observations of Galactic cirrus. De Vaucouleurs & Pence (1978) estimated the $`B`$-band surface brightness of the Galactic disk towards the Galactic poles to be $`\mu _B=24.92`$ mag arcsec<sup>-2</sup>. Thus we infer a ratio for the Galactic cirrus to be approximately
$$\frac{\nu I_\nu (12\mu \mathrm{m})}{\nu I_\nu (0.44\mu \mathrm{m})}_{\mathrm{Gal}.\mathrm{cirrus}}0.15,$$
(2)
comparable to our results at 6.75 and 15 $`\mu `$m.
Alternatively, it may be more illuminating to compare for quiescent regions the integrated total-infrared flux to the integrated flux in the $`B`$-band. From this comparison we can infer the ratio of heating output to heating input (i.e. the total-infrared flux to the UV and optical energy responsible for the dust heating). As a start, we contrast the integrated flux within the LW2 and LW3 bands to that in the $`B`$-band ($`\mathrm{\Delta }\nu /\nu 0.22`$):
$$\frac{\mathrm{\Delta }\nu I_\nu (6.75)+\mathrm{\Delta }\nu I_\nu (15)}{\mathrm{\Delta }\nu I_\nu (0.44\mu \mathrm{m})_{25}}0.53\pm 0.2.$$
(3)
Note that this is only valid for the outer regions of normal star-forming galaxies. To infer the total-infrared emission for diffuse cirrus regions using only the flux in the LW2 and LW3 bands, we have to rely on a model for the infrared spectral energy distribution of galaxies. Dale et al. (2000) present such a model for normal star-forming galaxies. Constrained by IRAS, ISOCAM, and ISOPHOT observations for our sample of 69 normal galaxies, the model reproduces well the empirical spectra and infrared color trends. It also allows us to determine the infrared energy budget for normal galaxies, and of particular interest here, to translate mid-infrared fluxes into total-infrared fluxes. The model shows that, for cirrus-dominated regions, the total-infrared emission (from 3–1100 $`\mu `$m) is approximately 10.3 times the emission appearing within just the LW2 and LW3 mid-infrared bands (and approximately 2.8 times the far-infrared emission from 42–122 $`\mu `$m). Coupling this result with that expressed by Equation 3, we see that
$$\frac{\mathrm{TIR}(31100\mu \mathrm{m})}{\mathrm{\Delta }\nu \mathrm{I}_\nu (0.44\mu \mathrm{m})_{25}}5.4$$
(4)
in the outskirts of normal galaxies. Observations of cirrus-dominated galaxies show that this finding is representative of quiescent disks: in our sample of normal late-type galaxies, the lowest such integrated total-infrared to blue ratios, assuming TIR/FIR $`2.8`$ for the most quiescent regions, are of order 3.
From the above ratio (Equation 4), one can derive the ratio of the heating output to the heating input as an indication of the optical depth in the interstellar medium. This is expected to be significantly less than unity for the most quiescent regions. All photon energies from the UV to the near-infrared contribute to the heating of dust grains. In fact, photons responsible for the $`B`$-band flux only contribute about 10-15% towards the total far-infrared flux in the outer parts of the disk of M31, a galaxy well-known for its overall quiescent behavior (Xu & Helou 1996 and C. Xu personal communication). Thus if we assume the results from the outer disk regions of M31 can reasonably apply for the outermost portions of our normal galaxy sample, the ratio of the total heating output to the total heating input for quiescent regions is typically of order 0.7 (to within a factor of two).
### 7.4 The Nature of the Flux at 4.5 $`\mu `$m
The transition from stellar to interstellar emission is well illustrated by the spectra of Virgo Cluster galaxies collected by Boselli et al. (1998). Its precise location and therefore the interpretation to attach to mid-infrared fluxes can be parametrized by the ratio of far-infrared to visible light fluxes. Interstellar dust emission takes over by 5 $`\mu `$m when the FIR/B ratio exceeds 0.5, and at shorter wavelengths for higher ratios (see Table 1 for the definition of FIR/B). As might be expected, elliptical galaxies are dominated by stellar emission, both photospheric and from circumstellar dust shells, and therefore provide the templates that one subtracts to isolate the interstellar emission component in spiral galaxies (Boselli et al. 1998; Madden, Vigroux & Sauvage 1999).
We do not expect large contributions from photospheric emission at 4.5 $`\mu `$m for our sample in light of the results from mid-infrared ISOPHOT spectrophotometry for 45 Key Project galaxies. The shape and flux level of the 3–5 $`\mu `$m continuum shows no variation with the strength of the aromatic features in emission for galaxies with global far infrared to blue ratios greater than unity (and very little variation otherwise), direct evidence that this portion of the continuum arises primarily from fluctuating dust grain emission (Helou et al. 2000; Lu et al. 2000). We can similarly explore the nature of the 4.5 $`\mu `$m emission by analyzing the ratios of broad band fluxes. Figure 14 displays the $`\frac{I_\nu (4.5\mu \mathrm{m})}{I_\nu (6.75\mu \mathrm{m})}`$ color ratio as a function of FIR/B. Gauging by the range of observed $`\frac{I_\nu (4.5\mu \mathrm{m})}{I_\nu (6.75\mu \mathrm{m})}`$ ratios, we can immediately rule out the possibility that the emission at both 4.5 and 6.75 $`\mu `$m is purely photospheric: the typical logarithmic value for the early type galaxy sample of Madden, Vigroux & Sauvage (1999) is much higher (log $`\frac{I_\nu (4.5\mu \mathrm{m})}{I_\nu (6.75\mu \mathrm{m})}`$ $`=`$ 0.1–0.4), consistent with the ratio for a $`T=3500`$ K blackbody profile that has been suggested for the stellar component of normal galaxies (Boselli et al. 1998). Thus for our sample of normal galaxies, the average $`\frac{I_\nu (4.5\mu \mathrm{m})}{I_\nu (6.75\mu \mathrm{m})}`$ ratio tells us that the 6.75 $`\mu `$m emission is definitely dominated by emission from interstellar dust grains, but does not constrain the origin of the 4.5 $`\mu `$m emission.
The slight trend in the data can tell us more. First, the trend is in the sense of a “photospheric excess” at 4.5 $`\mu `$m for the galaxies with lower FIR/B. Second, the slope of this trend is consistent with the extreme scenario whereby the 4.5 $`\mu `$m emission is proportional to that in the $`B`$-band, while the 6.75 $`\mu `$m emission is directly correlated with the far-infrared emission; the dotted line indicates the (extinction-corrected) inverse one-to-one trend. Thus, at first glance it would appear that these data suggest there are photospheric contributions to the 4.5 $`\mu `$m emission in normal galaxies, at least for those with FIR/B $``$ 1. The small number of data points, though, limit the robustness of any such claim. Finally, if we assume that for higher FIR/B ratios ($`5`$) the trend levels off, as the results of Helou et al. (2000) and Lu et al. (2000) suggest, then the $`\frac{I_\nu (4.5\mu \mathrm{m})}{I_\nu (6.75\mu \mathrm{m})}`$ ratio for the “pure ISM” is about 0.15.
## 8 Conclusion
We present mid-infrared maps at 6.75 and 15 $`\mu `$m for 61 normal star-forming galaxies; for a subset of 13 of these galaxies we also show maps at 4.5 $`\mu `$m. All galaxies for which observations were attempted at these wavelengths were successfully detected. Qualitatively, the mid-infrared morphology is not a strong function of wavelength, and many of the optical features of the galaxies are also observed in the mid-infrared, except for the bulges of spiral galaxies, consistent with the findings of Helou et al. (1996) in NGC 6946. Moreover, the data support non-negligible photospheric contributions at 4.5 $`\mu `$m for galaxies exhibiting low FIR/B ratios, consistent with the conclusion drawn from mid-infrared ISOPHOT spectroscopy for normal galaxies (Helou et al. 2000; Lu et al. 2000). However, the evidence presented here is tenuous, as we only have broad band data at 4.5 $`\mu `$m for a small number of galaxies.
Mid-infrared curve of growth profiles indicate that the mid-infrared emission is very well resolved by the ISOCAM maps for most of these galaxies. Moreover, the profiles are generally exponential in nature, and the distribution of 6.75 and 15 $`\mu `$m emission is quite similar for most galaxies. However, four of the six galaxies that show an enhanced $`\frac{I_\nu (15\mu \mathrm{m})}{I_\nu (6.75\mu \mathrm{m})}`$ color also show signs of active central star formation, with yet greater enhancement of $`\frac{I_\nu (15\mu \mathrm{m})}{I_\nu (6.75\mu \mathrm{m})}`$ in the central regions.
Quiescent galaxies, those showing low global interstellar heating intensities (i.e. $`\frac{I_\nu (60\mu \mathrm{m})}{I_\nu (100\mu \mathrm{m})}`$ $`0.6`$), have an almost constant $`\frac{I_\nu (6.75\mu \mathrm{m})}{I_\nu (15\mu \mathrm{m})}`$ color near unity. For galaxies with higher global heating intensities, the mid-infrared color drops rapidly with increasing far-infrared color. We interpret this as evidence for global mid-infrared spectral energy distributions that are increasingly dominated by H II region emission, characterized by a relatively steep slope in the mid-infrared continuum and a depressed contribution from PAHs. It is interesting to note that galaxies of all morphological types appear to follow the same color-color trend.
We have estimated the average mid-infrared surface brightness at which the mid-infrared semi-major axis matches that in the optical (at the $`B`$-band 25 mag arcsec<sup>-2</sup> level). We find $`R(0.44\mu \mathrm{m})_{25}`$ $``$ $`R(6.75\mu \mathrm{m})`$ at 0.04 MJy sr<sup>-1</sup> and $`R(0.44\mu \mathrm{m})_{25}`$ $``$ $`R(15\mu \mathrm{m})`$ at 0.09 MJy sr<sup>-1</sup> on average for this sample. These mid-infrared surface brightness levels are consistent with observations of Galactic cirrus and the Solar Neighborhood, implying a reasonable similarity in interstellar heating intensity for the outskirts of normal galaxies and our Galaxy. A final interesting finding from the mid-infrared size analysis centers on the ratio of total heating output to the total heating input for quiescent regions. We find that ratio to be of order 0.7.
This work was supported by ISO data analysis funding from the US National Aeronautics and Space Administration, and carried out at the Infrared Processing and Analysis Center (IPAC) and the Jet Propulsion Laboratory of the California Institute of Technology. The ISOCAM data presented in this paper was analysed using CIA, a joint development by the ESA Astrophysics Division and the ISOCAM Consortium. The Digitized Sky Surveys were produced at the Space Telescope Science Institute under U.S. Government grant NAG W-2166. The images of these surveys are based on photographic data obtained using the Oschin Schmidt Telescope on Palomar Mountain and the UK Schmidt Telescope. The plates were processed into the present compressed digital form with the permission of these institutions.
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# Isotropic phase squeezing and the arrow of time
PACS numbers: 03.65.-w; 03.67.-a
## Abstract
We prove that isotropic squeezing of the phase is equivalent to reversing the arrow of time.
The concept of “squeezing” appeared in the literature in the early 70’s and was extensively studied in order to improve the capacity of quantum information channels and the sensitivity in interferometric measurements . Since then squeezing has been a very popular word in quantum optics. By “squeezing” one refers to a physical process where the uncertainty of an observable is reduced at the expense of increasing the uncertainty of the conjugated observable, according to the Heisenberg inequalities (for extensive reviews see for example Ref. ). In quantum optics quadrature squeezing, namely the squeezing of the probability distribution of the observable $`a_\varphi =(a^{}e^{i\varphi }+ae^{i\varphi })/2`$$`a`$ and $`a^{}`$ denoting the annihilation and creation operators of a given radiation mode—has been achieved experimentally, giving rise to a number of interesting properties, such as phase-sensitive amplification and antibunching . More recently, the density operator of squeezed states has been measured by optical homodyne tomography .
Conjugated quadratures, i.e. quadratures relative to phases $`\varphi `$ and $`\varphi +\pi /2`$, are generalizations of the couple of observables position-momentum. Thus, we can view quadrature squeezing of radiation states as a narrowing process of the probability distribution in the phase space which occurs in a definite direction, corresponding to the phase of the squeezed quadrature \[see Fig. 1(a)\]. Classically, one can imagine a similar process in polar coordinates \[Fig. 1(b)\], where the radial probability distribution is squeezed, while the phase is spread, or viceversa. In the phase space the squared radius corresponds to the total energy of the harmonic oscillator, which is proportional to the photon number operator $`N=a^{}a`$ for a single-mode radiation field. Number squeezing narrows the photon number distribution, with the possibility of achieving sub-Poissonian statistics $`\mathrm{\Delta }N^2<N`$ in photon counting . This process has been investigated extensively and can be experimentally achieved by means of self-phase modulation in Kerr media .
The inverse process, namely phase squeezing, is the subject of the present Letter. We will consider isotropic phase squeezing, namely squeezing of the phase probability distribution independently of the mean value of the phase. Such a process corresponds to noise reduction in the measurement of phase, and it would lead to important results for communications and measurements, such as improved sensitivity of interferometric schemes and the achievement of the capacity of quantum communications based on phase coding.
In the following we will prove that isotropic phase squeezing cannot be realized because it would correspond to reversing the arrow of time. The arrow of time is statistically defined by the direction of the irreversible dynamics of open systems . In quantum-mechanical terms, it is associated to a loss of coherence of the quantum state, e.g. dephasing mechanism of the laser light, which corresponds to a random walk on the phase space . We will prove that any dynamical process that isotropically reduces the phase uncertainty can be described only in terms of a “time-reversed dissipative equation”.
In the literature the Heisenberg-like heuristic inequality $`\mathrm{\Delta }N\mathrm{\Delta }\varphi 1`$ for the couple number-phase is often reported. However, its meaning is only semiclassical, since the quantum phase does not correspond to any self-adjoint operator . Therefore, in order to investigate isotropic phase squeezing, we have first to introduce the concepts of phase measurement and phase probability distribution in a rigorous way.
The quantum-mechanical definition of the phase is well assessed in the framework of quantum estimation theory . In this context the phase of a quantum state is defined by the shift $`\varphi `$ generated by any operator $`F`$ with discrete spectrum. For example $`F=a^{}a`$ for the harmonic oscillator, and $`F=\sigma _z/2`$ for a two-level system, $`\sigma _z`$ being the customary Pauli operator.
Quantum estimation theory provides a general description of quantum statistics in terms of POVM’s (positive operator-valued measures) and seeks the optimal POVM to estimate one or more parameters of a quantum system on the basis of a cost function which assesses the cost of errors in the estimates. For phase estimation, the optimal POVM for pure states $`|\psi `$ with coefficients $`\psi _n=|\psi _n|e^{i\chi _n}0`$ on the basis $`|n`$ of $`F`$ eigenvectors is given by
$`d\mu (\varphi )={\displaystyle \frac{d\varphi }{2\pi }}|e(\varphi )e(\varphi )|`$ (1)
for the class of Holevo’s cost functions—a large class including the maximum likelihood criterion, the $`2\pi `$-periodicized variance, and the fidelity optimization. In Eq. (1) $`|e(\varphi )`$ denotes the (Dirac) normalizable vector
$`|e(\varphi )={\displaystyle \underset{nS}{}}e^{i(n\varphi \chi _n)}|n,`$ (2)
where $`S`$ is the spectrum of $`F`$. In Ref. the solution given in Eqs. (1-2) has also been proved for phase-pure states, namely for states described by a density operator $`\rho `$ satisfying the condition
$`\rho _{nm}n|\widehat{\rho }|m=|\rho _{nm}|e^{i(\chi _n\chi _m)},`$ (3)
and for a nondegenerate phase-shift generator $`F`$ . For states that are not of this kind, there is no available method in the literature to obtain the optimal POVM, and thus the concept itself of phase does not have a well defined meaning.
The phase probability distribution $`dp(\varphi )`$ of a quantum state is evaluated by means of the optimal POVM in Eq. (1) through the Born’s rule $`dp(\varphi )=\mathrm{Tr}[\rho d\mu (\varphi )]`$. The phase uncertainty $`\mathrm{\Delta }\varphi ^2`$ of the state can then be calculated. However, notice that for periodic distributions the customary r.m.s. deviation depends on the chosen window of integration. Definitions of phase uncertainty that do not depend on the interval of integration given in the literature are monotonic increasing functions $`f`$ of some average cost of the Holevo’s class, namely
$`\delta \varphi =f(C),`$ (4)
where $`C`$ represents the cost operator
$`C=c_0{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}c_n(e_+^n+e_{}^n),c_n0,n1,`$ (5)
and $`\mathrm{}`$ represents the quantum ensemble average. In Eq. (5) we introduced the following notation
$`e_+={\displaystyle \underset{nS}{}}e^{i(\chi _{n+1}\chi _n)}|n+1n|\text{and }e_{}=(e_+)^{}.`$ (6)
Typical examples of functions of this kind are the reciprocal peak likelihood and the phase deviation $`2(1|e_+|^2)`$. The former corresponds to $`f(x)=1/x`$ for the cost operator $`C`$ with all $`c_n=1`$; the latter corresponds to $`f=2\left[1(1/4)x^2\right]`$ with $`c_1=1`$ and $`c_n=0`$ $`n1`$ (for phase-pure states $`e_+`$ is a real positive quantity, so $`e_{}=e_+`$ and $`C^2/4=|e_+|^2`$, with $`C0`$).
Now we introduce the concept of isotropic phase squeezing. For the e.m. field, we remind that ordinary quadrature squeezing is effective in reducing the phase uncertainty of a quantum state provided that the average value of the phase is known a priori. As mentioned above, isotropic phase squeezing should reduce the phase uncertainty of the state $`\rho `$ independently of the initial mean phase. In mathematical terms, this condition corresponds to a linear map $`\mathrm{\Gamma }`$ that is covariant for the rotation group generated by the operator $`F`$, namely
$`\mathrm{\Gamma }(e^{iF\varphi }\rho e^{iF\varphi })=e^{iF\varphi }\mathrm{\Gamma }(\rho )e^{iF\varphi }.`$ (7)
A physically realizable linear map $`\mathrm{\Gamma }`$ corresponds to a completely positive (CP) map for density operators that can be written in the Lindblad form
$`{\displaystyle \frac{\rho }{t}}={\displaystyle \underset{n}{}}L[V_n]\rho ,`$ (8)
where $`L[O]\rho O\rho O^{}\frac{1}{2}(O^{}O\rho +\rho O^{}O)`$ denotes the Lindblad superoperator . We do not take into account the customary Hamiltonian term $`i[H,\rho ]`$ in the master equation (8), because for the covariance condition one has $`[H,F]=0`$, and hence the optimal POVM is simply rotated, with the result that the phase uncertainty is not affected by the corresponding unitary evolution (such Hamiltonian term can be equivalently applied in one step before the evolution (8), and preserves the phase purity of the state).
The covariance condition restricts the general form of Eq. (8) to the expression
$`{\displaystyle \frac{\rho }{t}}={\displaystyle \underset{m=\mathrm{}}{\overset{+\mathrm{}}{}}}{\displaystyle \underset{j}{}}L[B_{m,j}]\rho ,`$ (9)
where
$`B_{mj}`$ $`=g_{m,j}(F)e_+^m,m0`$ (10)
$`B_{mj}`$ $`=h_{|m|,j}(F)e_{}^{|m|},m<0.`$
In the following we will focus our attention on the case of a single-mode radiation field, hence we take $`F=a^{}a`$ and the spectrum $`S=`$. We postpone the discussion of the generality of our result at the end of the paper.
We consider the class of phase-pure states as initial states for the master equation (9), since for other kinds of states the phase measurement is not well defined, as mentioned before.
We are now in position to prove the main result of this paper, namely that isotropic phase squeezing is equivalent to reversing the arrow of time. According to Eq. (4) the time derivative of the phase uncertainty $`\delta \varphi `$ has the same sign as the derivative of the average cost $`C`$, which is obtained from Eq. (9) as follows
$`{\displaystyle \frac{C}{t}}=2\text{Re}{\displaystyle \frac{}{t}}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}c_ke_+^k.`$ (11)
A straightforward calculation gives the following contribution for the $`k`$-th term in the sum of Eq. (11)
$`2\text{Re}{\displaystyle \frac{}{t}}e_+^k=`$ (12)
$`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{j}{}}{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}|\rho _{l,l+k}||g_{m,j}(m+l)g_{m,j}(m+l+k)|^2`$
$`+{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{j}{}}{\displaystyle \underset{l=m}{\overset{\mathrm{}}{}}}|\rho _{l,l+k}||h_{m,j}(lm)h_{m,j}(lm+k)|^2`$
$`+{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{j}{}}{\displaystyle \underset{l=\mathrm{max}(0,mk)}{\overset{m1}{}}}|\rho _{l,l+k}||h_{m,j}(lm+k)|^20,`$
which is manifestly non negative for all values of $`k`$. Hence, according to Eqs. (11) and (12), the average cost as well as the phase uncertainty increase versus time. The only possibility to achieve isotropic phase squeezing is then to have a minus sign in front of the master equation (8), which means to reverse the arrow of time.
Our proof rules out also the possibility of isotropic squeezing through a phase measurement followed by a feedback quadrature squeezing. In fact, such a kind of process is described by a CP map as well, and one can explicitly show that the phase uncertainty in the measurement eventually leads to an overall phase diffusion.
The same conclusions regarding the derivatives of $`e_+^k`$ hold also for the cases of unbounded spectrum $`S=`$ and bounded spectrum $`S=_q`$ for a nondegenerate phase shift operator: also in these cases the phase uncertainty can decrease for any input state only if we reverse the arrow of time. For $`S=_q`$ all series in Eq. (12) are bounded and boundary terms appear in addition to the third one. For $`S=`$ Eq. (12) rewrites:
$`2\text{Re}{\displaystyle \frac{}{t}}e_+^k=`$ (13)
$`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{j}{}}{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}|\rho _{l,l+k}||g_{m,j}(m+l)g_{m,j}(m+l+k)|^2`$
$`+{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{j}{}}{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}|\rho _{l,l+k}||h_{m,j}(lm)h_{m,j}(lm+k)|^20.`$
In the general case a phase-covariant master equation does not evolve a phase-pure state into a phase-pure state, hence it may happen in principle that, after a finite time interval in which the phase-purity is lost, phase-purity is then recovered at the end with an overall decrease of $`\delta \varphi `$. However, one cannot follow the evolution of $`\delta \varphi `$ for finite time intervals if the definition itself of the phase is lost during the time evolution. A phase-covariant master equation (9) preserves phase-purity if and only if $`\mathrm{arg}(g_{m,j}(F))=\phi _{m,j}`$ and $`\mathrm{arg}(h_{m,j}(F))=\theta _{m,j}`$ independent on $`F`$, and one can always choose $`\phi _{m,j}=\theta _{m,j}=0`$ identically, due to the bilinear form of the Lindblad superoperator.
In the case of degenerate $`F`$, we can find the optimal POVM for pure states $`|\psi \psi |`$ as follows : we select a vector $`|n_{}`$ for each degenerate eigenspace $`_n`$ corresponding to the eigenvalue $`n`$, such that $`|n_{}`$ is parallel to the projection of $`|\psi `$ on $`_n`$. So the Hilbert space $``$ can be represented as $`_{}_{}`$, where $`_{}`$ is the Hilbert space spanned by the vectors $`|n_{}`$ and $`_{}`$ its orthogonal completion. Since $`|\psi `$ has null component in $`_{}`$ the estimation problem reduces to a nondegenerate one in the Hilbert space $`_{}`$, and the optimal POVM is given by $`d\mu (\varphi )=d\mu _{}(\varphi )d\mu _{}(\varphi )`$, where $`d\mu _{}`$ is the optimal POVM for the nondegenerate estimation problem in $`_{}`$, while $`d\mu _{}(\varphi )`$ is an arbitrary POVM in $`_{}`$. It is clear that the POVM obtained in this way is optimal also for phase-pure states that are mixtures of pure states all with the same $`_{}`$. The most general phase-covariant master equation is again of the form in Eqs. (8) and (10), however, now there are infinitely many possible $`B_{m,j}`$ for fixed $`m,j`$ corresponding to different operators $`e_+`$ which shift the eigenvalue of $`F`$ while spreading the state in the whole $``$ from $`_{}`$ in all possible ways. Therefore, a phase-covariant master equation does not keep the original state in $`_{}`$, apart from the case where one considers operators $`B_{m,j}`$ defined in terms of $`e_+`$ only of the form
$`e_+={\displaystyle \underset{n}{}}e^{i(\chi _{n+1}\chi _n)}|n+1_{}n|.`$ (14)
In the general case, however, a reduction of the phase uncertainty is possible in principle, due to the arbitrariness introduced by $`d\mu _{}(\varphi )`$ in the definition of phase in the degenerate case, the time derivative of $`dp(\varphi )=\mathrm{Tr}[\rho d\mu _{}(\varphi )]`$ generally depending on $`d\mu _{}(\varphi )`$.
We now focus attention back to the case of nondegenerate $`F`$. Looking at Eq. (12), one can see that it is possible to make all terms vanishing, getting a null derivative for the average cost. This is actually possible only for unbounded spectra as $`S=`$ and $`S=`$, where one can find a master equation that preserves the phase uncertainty for any quantum state. For $`S=`$, one has the following conditions on the coefficients of Eq. (12): $`g_{m,j}(F)=g_{m,j}`$ constant and $`h_{m,j}(F)=0`$. Upon defining $`u_m=_j|g_{m,j}|^2`$ one has
$`{\displaystyle \frac{\rho }{t}}={\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}u_m(e_+^m\rho e_{}^m\rho ).`$ (15)
In the case $`S=`$ one has more generally $`h_{m,j}(F)=h_{m,j}`$ constant, and introducing $`v_m=_j|h_{m,j}|^2`$, the phase-uncertainty preserving master equation takes the form
$`{\displaystyle \frac{\rho }{t}}={\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}u_m(e_+^m\rho e_{}^m\rho )+v_m(e_{}^m\rho e_+^m\rho ).`$ (16)
The master equations (15) and (16) are very interesting, since they represent a counterexample to the customary identification of “decoherence” and “dephasing”. The study of physical realizations of Eqs. (15) and (16) could provide insight in the understanding of decoherence and relaxation phenomena.
In conclusion, we have shown that isotropic squeezing of the phase is equivalent to reversing the arrow of time. This result is very general, as it holds for any definition of phase with nondegenerate shift operator, for any definition of phase-uncertainty in the Holevo’s class, and for any initial phase-pure state. In this way we have related the concept of phase to the arrow of time statistically defined by the evolution of open quantum systems, thus enforcing the link between phase and time .
This work is supported by the Italian Ministero dell’Università e della Ricerca Scientifica e Tecnologica under the program Amplificazione e rivelazione di radiazione quantistica, and by the INFM PAIS 1999.
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# Introduction
## Introduction
Building on the work of Kac & Vainerman , Enock & Schwartz , Baaj & Skandalis , Woronowicz and Van Daele , a precise definition of a locally compact quantum group was recently introduced by the authors in , see and for an overview. For an overview of the historical development of the theory we refer to and the introduction of . This theory provides a topological framework to study quantum groups and it unifies locally compact groups, compact quantum groups and Kac algebras.
Because classical groups are usually defined to act on a space it is very natural to make a quantum group act on a quantum space, which will be an algebra. In an algebraic framework the study of Hopf algebras acting on algebras has been very useful.
On the other hand, actions of locally compact groups on von Neumann algebras have always been an important topic in operator algebra theory, see e.g. and . In these works the importance of Haagerup’s results on the canonical implementation of locally compact group actions and his results on the dual weight construction, cannot be overestimated. It is simply used all the time, without noticing it. See and .
Hence it seems natural to study more generally actions of locally compact quantum groups on von Neumann algebras and to try to develop the same machinery of canonical implementation and dual weight construction. This is what is done in the first half of this paper. We strongly believe that this will serve as an important tool in several applications of locally compact quantum groups. We already give some applications in the second half of this paper. Other applications are given by Kustermans in . Both will be explained below.
The special case of Kac algebra actions has been studied by Enock and Schwartz in and . They obtained important results on crossed products, with the biduality theorem as a major achievement. But they never obtained a unitary implementation for an arbitrary action and also Haagerup’s theory of dual weights on the crossed product could not be completely generalized. For instance, it remained an open problem whether the crossed product with a Kac algebra action on a von Neumann algebra is in standard position on its natural Hilbert space. It should be mentioned that in also Baaj & Skandalis obtain a biduality theorem for crossed products with multiplicative unitaries.
A first attempt to obtain the unitary implementation of a Kac algebra action was made by J.-L. Sauvageot in . Unfortunately his proof is wrong, and for this we refer to the discussion in the beginning of section 3.
So in this paper we will define actions of a locally compact quantum group on a von Neumann algebra and we will construct its unitary implementation. We will also give a construction for the dual weight on the crossed product and prove analogous results as those about group actions obtained by Haagerup in and . In particular we prove that the crossed product is in standard position on its natural Hilbert space and we identify the associated modular objects. Hence we do not only give a right proof for the results of Sauvageot, but also we prove new results on the dual weights which make them a workable and applicable tool, and we work in the more general setting of locally compact quantum groups.
In the second half of the paper we will give some applications of these results in the theory of subfactors and inclusions of von Neumann algebras. It has been proved by Enock and Nest in their beautiful papers and that every irreducible, depth 2 inclusion of factors satisfying the regularity condition, can be obtained as $`N^\alpha N`$ where $`\alpha `$ is an outer action of a locally compact quantum group on the factor $`N`$ and $`N^\alpha `$ is the fixed point algebra. We show in this paper that the action $`\alpha `$ is always integrable and that, conversely, for every outer and integrable action $`\alpha `$ on a factor $`N`$ the inclusion $`N^\alpha N`$ is irreducible, of depth 2 and regular. So we obtain a converse to the results of Enock and Nest. The same result is stated for the special case of a dual Kac algebra action in \[7, 11.14\], but not proved. While doing this, we study more generally the problem when the inclusion $`N^\alpha NM\text{α}N`$ is a basic construction, and here $`M\text{α}N`$ denotes the crossed product.
As a final application of our results we prove the equivalence of outerness and minimality of a locally compact quantum group action, under the integrability condition. We also generalize the main theorem of Yamanouchi to actions of arbitrary locally compact quantum groups: when working on separable Hilbert spaces, we prove that every integrable outer action with infinite fixed point algebra is a dual action.
It should also be mentioned that our results on the unitary implementation of a locally compact quantum group action are already applied in a recent paper by Kustermans (see ) in which he constructs induced corepresentations of locally compact quantum groups. Taking into account the importance of induced representations of locally compact groups, it is clear that the results of Kustermans serve as a major motivation for our work.
Acknowledgment : I would like to thank both J. Kustermans and L. Vainerman for their valuable advice, interesting and stimulating discussions and useful suggestions on this topic of locally compact quantum groups in action.
## Definitions and notations
The whole of this paper will rely heavily on the modular theory of von Neumann algebras. Throughout the text we will not make efforts to give references to the original papers, but we will use as a general reference.
When $`\theta `$ is a normal, semifinite and faithful (we say n.s.f.) weight on a von Neumann algebra $`N`$, one can make the so-called GNS-construction $`(K_\theta ,\pi _\theta ,\mathrm{\Lambda }_\theta )`$, where $`K_\theta `$ is a Hilbert space, $`\pi _\theta `$ is a normal representation of $`N`$ on $`K_\theta `$ and $`\mathrm{\Lambda }_\theta :𝒩_\theta K_\theta `$ is a linear map satisfying $`\pi _\theta (x)\mathrm{\Lambda }_\theta (y)=\mathrm{\Lambda }_\theta (xy)`$ for all $`xN`$ and $`y𝒩_\theta `$. Further $`\mathrm{\Lambda }_\theta (𝒩_\theta )`$ is dense in $`K_\theta `$. Here $`𝒩_\theta `$ is the left ideal in $`N`$ defined by $`\{xN\theta (x^{}x)<\mathrm{}\}`$. The representation $`\pi _\theta `$ is faithful and often we will identify $`N`$ and $`\pi _\theta (N)`$. Then we will use the expression: let us represent $`N`$ on the GNS-space of $`\theta `$ such that $`(K_\theta ,\iota ,\mathrm{\Lambda }_\theta )`$ is a GNS-construction.
We will use several standard notations and results of modular theory. We write
$$_\theta ^+=\{xN^+\theta (x)<\mathrm{}\}$$
and we denote with $`(\sigma _t^\theta )_t`$ the modular automorphism group of $`\theta `$. Further we denote with $`𝒯_\theta `$ the Tomita algebra defined by
$$𝒯_\theta =\{xNx\text{is analytic with respect to}(\sigma ^\theta )\text{and}\sigma _z^\theta (x)𝒩_\theta 𝒩_\theta ^{}\text{for all}z\}.$$
Given a GNS-construction $`(K_\theta ,\pi _\theta ,\mathrm{\Lambda }_\theta )`$ we define as usual the modular conjugation $`J_\theta `$ and the modular operator $`_\theta `$. Recall that
$$J_\theta _\theta ^{1/2}\mathrm{\Lambda }_\theta (x)=\mathrm{\Lambda }_\theta (x^{})$$
for all $`x𝒩_\theta 𝒩_\theta ^{}`$ and $`\mathrm{\Lambda }_\theta (𝒩_\theta 𝒩_\theta ^{})`$ is a core for $`_\theta ^{1/2}`$.
When $`\theta _1`$ and $`\theta _2`$ are n.s.f. weights on $`N`$ we denote with $`([D\theta _1:D\theta _2]_t)_t`$ the Connes cocycle as defined in e.g. \[26, 3.1\].
Often we will make use of operator valued weights. When $`N`$ is a von Neumann algebra we denote with $`N^+\text{ext}`$ the extended positive part of $`N`$ as introduced by Haagerup in , see e.g. \[26, 11.1\]. For the notion of operator valued weights we refer to or \[26, 11.5\]. We will denote with $`,`$ the composition of elements of $`N^+\text{ext}`$ and $`N_{}^+`$. When $`T`$ is an operator valued weight we denote with $`𝒩_T`$ the left ideal of elements $`x`$ such that $`T(x^{}x)`$ is bounded.
All tensor products in this paper are either von Neumann algebra tensor products or tensor products of Hilbert spaces. This will always be clear from the context. We will use $`\sigma `$ to denote the flip map on a tensor product $`AB`$ of von Neumann algebras and $`\mathrm{\Sigma }`$ to denote the flip map on a tensor product $`HK`$ of Hilbert spaces. When $`K`$ is a Hilbert space and $`\xi K`$ we denote with $`\theta _\xi `$ the operator in $`B(,K)`$ given by $`\theta (\lambda )=\lambda \xi `$ for all $`\lambda `$. When $`H`$ is a Hilbert space and $`\xi ,\eta H`$ we denote with $`\omega _{\xi ,\eta }`$ the usual vector functional in $`B(H)_{}`$ given by $`\omega _{\xi ,\eta }(x)=x\xi ,\eta `$. We use $`\omega _\xi `$ as a shorter notation for $`\omega _{\xi ,\xi }`$. We will denote with $`𝒟(T)`$ the domain of a (usually densily defined) map $`T`$.
Throughout this paper the pair $`(M,\mathrm{\Delta })`$ will denote a (von Neumann algebraic) locally compact quantum group. This means that
* $`M`$ is a von Neumann algebra and $`\mathrm{\Delta }:MMM`$ is a normal and unital $``$-homomorphism satisfying coassociativity: $`(\mathrm{\Delta }\iota )\mathrm{\Delta }=(\iota \mathrm{\Delta })\mathrm{\Delta }`$.
* There exist n.s.f. weights $`\phi `$ and $`\psi `$ on $`M`$ such that
+ $`\phi `$ is left invariant in the sense that $`\phi \left((\omega \iota )\mathrm{\Delta }(x)\right)=\phi (x)\omega (1)`$ for all $`x_\phi ^+`$ and $`\omega M_{}^+`$.
+ $`\psi `$ is right invariant in the sense that $`\psi \left((\iota \omega )\mathrm{\Delta }(x)\right)=\psi (x)\omega (1)`$ for all $`x_\psi ^+`$ and $`\omega M_{}^+`$.
We refer to , and for the theory of locally compact quantum groups in either the von Neumann algebra or C-algebra language. It should be stressed that in there is given a definition of a locally compact quantum group in the C-algebra framework, and it is proven that one can associate with this a locally compact quantum group in the von Neumann algebra framework. In the above definition of a von Neumann algebraic locally compact quantum group is given and it is shown how to associate with it a C-algebraic locally compact quantum group.
One can then prove the existence of a $`\sigma `$-strong closed map $`S`$ on $`M`$, called the antipode, such that for all $`a,b𝒩_\phi `$ we have
$$(\iota \phi )\left(\mathrm{\Delta }(a^{})(1b)\right)𝒟(S)\text{and}S\left((\iota \phi )\left(\mathrm{\Delta }(a^{})(1b)\right)\right)=(\iota \phi )\left((1a^{})\mathrm{\Delta }(b)\right).$$
Moreover, such elements $`(\iota \phi )\left(\mathrm{\Delta }(a^{})(1b)\right)`$ span a $`\sigma `$-strong core for $`S`$. Then there exists a unique one-parameter group $`(\tau _t)_t`$ of automorphisms of $`M`$ and a unique $``$-anti-automorphism $`R`$ of $`M`$ such that
$$S=R\tau _{i/2}R^2=\iota R\tau _t=\tau _tR\text{for all}t.$$
We call $`\tau `$ the scaling group of $`(M,\mathrm{\Delta })`$ and $`R`$ the unitary antipode. One refers to the expression $`S=R\tau _{i/2}`$ as the polar decomposition of the antipode.
Next one can prove that $`\mathrm{\Delta }R=\sigma (RR)\mathrm{\Delta }`$, where $`\sigma `$ denotes the flip map on $`MM`$. One can also prove that the left and right invariant weights on $`(M,\mathrm{\Delta })`$ are unique up to a positive scalar. So $`\psi `$ and $`\phi R`$ are proportional and we can suppose from the beginning that $`\psi =\phi R`$. We denote with $`(\sigma _t)_t`$ the modular group of $`\phi `$. Then there exists a unique self-adjoint, strictly positive operator $`\delta `$ affiliated with $`M`$ and satisfying $`\sigma _t(\delta )=\nu ^t\delta `$ and $`\psi =\phi _\delta `$, where $`\nu >0`$ is a real number. Formally this means that $`\psi (x)=\phi (\delta ^{1/2}x\delta ^{1/2})`$ and for an exact definition of $`\phi _\delta `$ we refer to \[29, 1.5\]. We call $`\delta `$ the modular element of $`(M,\mathrm{\Delta })`$. The number $`\nu >0`$ is called the scaling constant of $`(M,\mathrm{\Delta })`$.
Let us represent $`M`$ on the GNS-space $`H`$ of $`\phi `$ such that $`(H,\iota ,\mathrm{\Lambda })`$ is a GNS-construction for $`\phi `$. Let $`(H,\iota ,\mathrm{\Gamma })`$ be the canonical GNS-construction for $`\psi =\phi _\delta `$ as defined in \[19, 7.2\]. Then one can define unitaries $`W`$ and $`V`$ in $`B(HH)`$ such that
$`W^{}(\mathrm{\Lambda }(a)\mathrm{\Lambda }(b))`$ $`=(\mathrm{\Lambda }\mathrm{\Lambda })\left(\mathrm{\Delta }(b)(a1)\right)\text{for all}a,b𝒩_\phi `$
$`V(\mathrm{\Gamma }(a)\mathrm{\Gamma }(b))`$ $`=(\mathrm{\Gamma }\mathrm{\Gamma })\left(\mathrm{\Delta }(a)(1b)\right)\text{for all}a,b𝒩_\psi .`$
Here $`\mathrm{\Lambda }\mathrm{\Lambda }`$ and $`\mathrm{\Gamma }\mathrm{\Gamma }`$ denote the canonical GNS-maps for the tensor product weights $`\phi \phi `$ and $`\psi \psi `$. Then $`W`$ and $`V`$ are multiplicative unitaries, which means that they satisfy the pentagonal equation
$$W_{12}W_{13}W_{23}=W_{23}W_{12}.$$
Denoting with <sup>-</sup> the $`\sigma `$-strong closure we have that
$$M=\{(\iota \omega )(W)\omega B(H)_{}\}^{}\text{and}\mathrm{\Delta }(x)=W^{}(1x)W\text{for all}xM.$$
We will denote with $`J`$ and $``$ the modular conjugation and modular operator of $`\phi `$ in the GNS-construction $`(H,\iota ,\mathrm{\Lambda })`$.
Finally we describe how to define the dual locally compact quantum group $`(\widehat{M},\widehat{\mathrm{\Delta }})`$. Define the von Neumann algebra $`\widehat{M}`$ as follows, where again <sup>-</sup> denotes the $`\sigma `$-strong closure.
$$\widehat{M}=\{(\omega \iota )(W)\omega M_{}\}^{}.$$
Then one can define a comultiplication $`\widehat{\mathrm{\Delta }}`$ on $`\widehat{M}`$ by
$$\widehat{\mathrm{\Delta }}(y)=\mathrm{\Sigma }W(y1)W^{}\mathrm{\Sigma }\text{for all}y\widehat{M}$$
where $`\mathrm{\Sigma }`$ denotes the flip map on $`HH`$. When $`\omega M_{}`$ we define $`\lambda (\omega )=(\omega \iota )(W)`$. Of course $`M_{}`$ should be thought of as the L<sup>1</sup>-functions on the quantum group $`(M,\mathrm{\Delta })`$, and then $`\lambda `$ is the left regular representation. We also define
$$=\{\omega M_{}\text{there exists}\eta H\text{such that}\omega (x^{})=\eta ,\mathrm{\Lambda }(x)\text{for all}x𝒩_\phi \}.$$
Such a $`\eta H`$ is necessarily uniquely determined and will be denoted with $`\xi (\omega )`$. Then there exists a unique n.s.f. weight $`\widehat{\phi }`$ on $`\widehat{M}`$ with GNS-construction $`(H,\iota ,\widehat{\mathrm{\Lambda }})`$ such that $`\lambda ()𝒩_{\widehat{\phi }}`$, $`\lambda ()`$ is a $`\sigma `$-strong–norm core for $`\widehat{\mathrm{\Lambda }}`$ and $`\widehat{\mathrm{\Lambda }}(\lambda (\omega ))=\xi (\omega )`$ for all $`\omega `$.
Then $`(\widehat{M},\widehat{\mathrm{\Delta }})`$ will be a locally compact quantum group, and having fixed the GNS-construction $`(H,\iota ,\widehat{\mathrm{\Lambda }})`$ for $`\widehat{\phi }`$ we can now repeat the story about $`(M,\mathrm{\Delta })`$ and obtain $`(\widehat{\sigma }_t)_t`$,$`\widehat{\delta }`$,$`\widehat{W}`$,$`\widehat{V}`$,$`\widehat{J}`$ and $`\widehat{}`$. For all kinds of formulas relating these objects we refer to . We only mention that
$`\widehat{J}J`$ $`=\nu ^{i/4}J\widehat{J}`$
$`\widehat{}^{it}x\widehat{}^{it}=\tau _t(x)`$ $`\text{and}\widehat{J}x^{}\widehat{J}=R(x)\text{for all}xM,t.`$
Finally we denote with $`(M,\mathrm{\Delta })\text{op}`$ the opposite locally compact quantum group $`(M,\mathrm{\Delta }\text{op})`$ where $`\mathrm{\Delta }\text{op}=\sigma \mathrm{\Delta }`$. Further we define $`(M,\mathrm{\Delta })^{}=(M^{},\mathrm{\Delta }^{})`$ where
$$\mathrm{\Delta }^{}(x)=(JJ)\mathrm{\Delta }(JxJ)(JJ)$$
for all $`xM^{}`$ and we call $`(M,\mathrm{\Delta })^{}`$ the commutant locally compact quantum group. Then one can prove that
$$(M,\mathrm{\Delta })\text{op}\widehat{\text{}}=(M,\mathrm{\Delta })\widehat{\text{}}^{}$$
and for this we refer to .
## 1 Actions of locally compact quantum groups
In this section we define actions of locally compact quantum groups on von Neumann algebras and we construct an important tool: the canonical operator valued weight from the von Neumann algebra on which we act to the fixed point algebra, obtained by integrating out the action.
###### Definition 1.1.
Let $`N`$ be a von Neumann algebra. A normal, injective and unital $``$-homomorphism $`\alpha :NMN`$ will be called a left action of $`(M,\mathrm{\Delta })`$ on $`N`$ when $`(\iota \alpha )\alpha =(\mathrm{\Delta }\iota )\alpha `$.
A normal, injective and unital $``$-homomorphism $`\alpha :NNM`$ will be called a right action of $`(M,\mathrm{\Delta })`$ on $`N`$ when $`(\alpha \iota )\alpha =(\iota \mathrm{\Delta })\alpha `$.
In this paper we will only work with left actions and so we drop the predicate *left*. When $`\alpha `$ is a right action, $`\sigma \alpha `$ will be a left action of $`(M,\mathrm{\Delta }\text{op})`$ on $`N`$, where $`\sigma `$ denotes the flip map from $`NM`$ to $`MN`$ and $`\mathrm{\Delta }\text{op}`$ denotes the opposite comultiplication. So it is not a real restriction to work only with left actions. It should be observed that in and they work with right actions.
###### Definition 1.2.
When $`\alpha :NMN`$ is an action of $`(M,\mathrm{\Delta })`$ on $`N`$ we define the fixed point algebra $`N^\alpha `$ as
$$N^\alpha =\{xN\alpha (x)=1x\}.$$
It is clear that $`N^\alpha `$ is a von Neumann subalgebra of $`N`$.
Recall that $`N^+\text{ext}`$ denotes the extended positive part of $`N`$.
###### Proposition 1.3.
Let $`N`$ be a von Neumann algebra and $`\alpha `$ an action of $`(M,\mathrm{\Delta })`$ on $`N`$. For every $`xN^+`$ the element $`T_\alpha (x)=(\psi \iota )\alpha (x)`$ of $`N^+\text{ext}`$ belongs to $`(N^\alpha )^+\text{ext}`$. Further $`T_\alpha `$ is a normal, faithful operator valued weight from $`N`$ to $`N^\alpha `$.
###### Proof..
Let $`xN^+`$ and $`\omega (MN)_{}^+`$. Denote with $`,`$ the composition of an element in $`N^+\text{ext}`$ and an element in $`N_{}^+`$. Then by definition of the operator valued weight $`\psi \iota `$ we get
$`T_\alpha (x),\omega \alpha =(\psi \iota )\alpha (x),\omega \alpha `$ $`=\psi \left((\iota \omega \alpha )\alpha (x)\right)`$
$`=\psi \left((\iota \omega )(\mathrm{\Delta }\iota )\alpha (x)\right)`$
$`=(\psi \iota \iota )\left((\mathrm{\Delta }\iota )\alpha (x)\right),\omega .`$
By the right invariant version of \[21, 3.1\] we get that
$$T_\alpha (x),\omega \alpha =1(\psi \iota )\alpha (x),\omega =1T_\alpha (x),\omega .$$
From this it follows that $`\alpha (T_\alpha (x))=1T_\alpha (x)`$. So we get that $`T_\alpha (x)(N^\alpha )^+\text{ext}`$.
If $`xN^+`$ and $`aN^\alpha `$ we have
$$T_\alpha (axa^{}),\omega =(\psi \iota )\left((1a)\alpha (x)(1a^{})\right),\omega =\psi \left((\iota a^{}\omega a)\alpha (x)\right)=T_\alpha (x),a^{}\omega a.$$
So we get that $`T_\alpha `$ is indeed an operator valued weight. Because both $`\alpha `$ and $`\psi \iota `$ are faithful and normal, also $`T_\alpha `$ is faithful and normal. ∎
One should observe that the same result is stated and used in for Kac algebra actions. Their proof contains a gap because they do not have an invariance formula like \[21, 3.1\], which is indispensable. For Kac algebra actions this was repaired by Zsidó (see , see also \[26, 18.18 and 18.23\]). Also in the case of a group action this was a non-trivial problem (see ). The simpler proof of Zsidó for this invariance formula only works in the Kac algebra setting, where the scaling group is trivial. I would like to thank prof. Enock for the discussion on this topic.
###### Definition 1.4.
An action $`\alpha `$ of $`(M,\mathrm{\Delta })`$ on a von Neumann algebra $`N`$ is called integrable when the operator valued weight $`T_\alpha `$ defined in proposition 1.3 is semifinite.
We will now introduce the well known concept of cocycle equivalent actions (cfr. \[4, I.6\]).
###### Definition 1.5.
Let $`\alpha `$ be an action of $`(M,\mathrm{\Delta })`$ on the von Neumann algebra $`N`$. A unitary $`UMN`$ is called an $`\alpha `$-cocycle if
$$(\mathrm{\Delta }\iota )(U)=U_{23}(\iota \alpha )(U).$$
It is clear that in this case the formula
$$\beta (x)=U\alpha (x)U^{}\text{for all}xN$$
defines an action $`\beta `$ of $`(M,\mathrm{\Delta })`$ on $`N`$.
Two actions $`\alpha `$ and $`\beta `$ of $`(M,\mathrm{\Delta })`$ on $`N`$ are called cocycle equivalent if there exists an $`\alpha `$-cocycle $`U`$ such that $`\beta `$ is given by the formula above.
It is easy to see that $`U^{}`$ is a $`\beta `$-cocycle when $`U`$ is an $`\alpha `$-cocycle and when $`\beta (x)=U\alpha (x)U^{}`$ for all $`xN`$.
## 2 Crossed products, the dual action and the duality theorem
In this section we fix an action $`\alpha `$ of a locally compact quantum group $`(M,\mathrm{\Delta })`$ on a von Neumann algebra $`N`$. We will define the crossed product $`M\text{α}N`$ in a similar way as in . We will also state some classical theorems concerning crossed products, the biduality theorem being the major one, but we will omit proofs because they are completely analogous to the proofs of and . See also , where the results of and are generalized to actions of Woronowicz algebras.
###### Definition 2.1.
We define the crossed product of $`N`$ with respect to the action $`\alpha `$ of $`(M,\mathrm{\Delta })`$ on $`N`$ as the von Neumann subalgebra of $`B(H)N`$ generated by $`\alpha (N)`$ and $`\widehat{M}`$. We denote this crossed product with $`M\text{α}N`$. So we have
$$M\text{α}N=(\alpha (N)\widehat{M})^{\prime \prime }.$$
We will now define in the usual way the dual action, which will be an action of $`(\widehat{M},\widehat{\mathrm{\Delta }}\text{op})`$ on $`M\text{α}N`$.
###### Proposition 2.2.
There exists a unique action $`\widehat{\alpha }`$ of $`(\widehat{M},\widehat{\mathrm{\Delta }}\text{op})`$ on $`M\text{α}N`$ such that
$`\widehat{\alpha }(\alpha (x))`$ $`=1\alpha (x)\text{for all}xN`$
$`\widehat{\alpha }(a1)`$ $`=\widehat{\mathrm{\Delta }}\text{op}(a)1\text{for all}a\widehat{M}.`$
Moreover when we denote with $`\stackrel{~}{W}`$ the unitary $`(JJ)\mathrm{\Sigma }W\mathrm{\Sigma }(JJ)`$ we have
$$\widehat{\alpha }(z)=(\stackrel{~}{W}1)(1z)(\stackrel{~}{W}^{}1)\text{for all}zM\text{α}N.$$
As we already mentioned, Enock and Schwartz deal with right actions in and . Hence they also give a slightly different definition for the crossed product, but in fact our definition agrees with theirs. When $`\alpha `$ is a right action of $`(M,\mathrm{\Delta })`$ on $`N`$ they define $`N_\alpha M`$ to be $`(\alpha (N)\widehat{M}^{})^{\prime \prime }`$. This is in accordance with our definition, because $`\sigma \alpha `$ is a left action of $`(M,\mathrm{\Delta }\text{op})`$ on $`N`$. The dual of $`(M,\mathrm{\Delta }\text{op})`$ is $`(\widehat{M}^{},\widehat{\mathrm{\Delta }}^{})`$, which gives
$$M\text{σα}N=(\sigma \alpha (N)\widehat{M}^{})^{\prime \prime }.$$
So we have $`N_\alpha M=\sigma (M\text{σα}N)`$, which shows that both definitions in fact agree.
Let us introduce the following concept, which will be needed later on. See also \[5, III.1\].
###### Definition 2.3.
Let $`\rho `$ be a self-adjoint, strictly positive operator affiliated with $`M`$. Then a n.s.f. weight $`\theta `$ on $`N`$ is called $`\rho `$-invariant if
$$\theta \left((\omega _{\xi ,\xi }\iota )\alpha (x)\right)=\rho ^{1/2}\xi ^2\theta (x)$$
for all $`x_\theta ^+`$ and $`\xi 𝒟(\rho ^{1/2})`$.
We will always work with $`\delta ^1`$-invariant weights, where $`\delta `$ is the modular element of the locally compact quantum group in action.
Then the following result can be proved as in \[9, 2.9\]. For the last statement of the next proposition observe that $`\tau _t(\delta )=\delta `$ and so the self-adjoint operators $`\delta `$ and $`\widehat{}`$ commute strongly. Hence their product $`\delta \widehat{}`$ is closable.
###### Proposition 2.4.
When $`\theta `$ is a n.s.f. $`\delta ^1`$-invariant weight on $`N`$ with GNS-construction $`(H_\theta ,\pi _\theta ,\mathrm{\Lambda }_\theta )`$, then there exists a unique unitary $`V_\theta MB(H_\theta )`$ such that for all $`\xi 𝒟(\delta ^{1/2})`$, $`\eta H`$ and $`x𝒩_\theta `$
$$(\omega _{\xi ,\eta }\iota )(V_\theta )\mathrm{\Lambda }_\theta (x)=\mathrm{\Lambda }_\theta \left((\omega _{\delta ^{1/2}\xi ,\eta }\iota )\alpha (x)\right).$$
Denote with $`J_\theta `$ and $`_\theta `$ the modular conjugation and modular operator of $`\theta `$. Then $`V_\theta `$ satisfies
$`(\mathrm{\Delta }\iota )(V_\theta )`$ $`=V_{\theta \mathrm{\hspace{0.17em}23}}V_{\theta \mathrm{\hspace{0.17em}13}}`$
$`(\iota \pi _\theta )\alpha (x)`$ $`=V_\theta (1\pi _\theta (x))V_\theta ^{}\text{for all}xN`$
$`V_\theta (\widehat{J}J_\theta )`$ $`=(\widehat{J}J_\theta )V_\theta ^{}`$
$`V_\theta (Q_\theta )`$ $`=(Q_\theta )V_\theta \text{where}Q\text{ is the closure of }\delta \widehat{}.`$
The following result is crucial (see \[9, 2.8\]).
###### Proposition 2.5.
* Let $`\alpha `$ be an integrable action of $`(M,\mathrm{\Delta })`$ on $`N`$ and denote with $`T_\alpha `$ the operator valued weight defined in proposition 1.3. Let $`\mu `$ be a n.s.f. weight on $`N^\alpha `$. Then $`\mu {}_{}{}^{}T_{\alpha }^{}`$ is a $`\delta ^1`$-invariant weight on $`N`$.
* Every dual action is integrable.
With these results at hand one can copy the proofs of to obtain the well known biduality theorem. Before we state this theorem we have to clarify some terminology. The dual action $`\widehat{\alpha }`$ is an action of $`(M,\mathrm{\Delta })\widehat{\text{}}\text{op}`$ on $`M\text{α}N`$. So we can make the double crossed product $`\widehat{M}_{\widehat{\alpha }}(M\text{α}N)`$ in $`B(HH)N`$ and on this double crossed product there is an action $`\widehat{\alpha }\widehat{\text{}}`$ of $`(M,\mathrm{\Delta })\widehat{\text{}}\text{op}\widehat{\text{}}\text{op}`$. Now $`(M,\mathrm{\Delta })\widehat{\text{}}\text{op}\widehat{\text{}}\text{op}=(M,\mathrm{\Delta })^{}\text{op}`$ and we can define an isomorphism of locally compact quantum groups
$$𝒥:(M,\mathrm{\Delta })(M,\mathrm{\Delta })^{}\text{op}$$
given by $`𝒥(x)=\widehat{J}JxJ\widehat{J}`$ for all $`xM`$.
###### Theorem 2.6 (Biduality theorem).
1. We have $`B(H)N=\left(B(H)\alpha (N)\right)^{\prime \prime }`$.
2. The map $`\mathrm{\Phi }`$ from $`B(H)N`$ to $`B(HH)N`$ defined by
$$\mathrm{\Phi }(z)=(W1)(\iota \alpha )(z)(W^{}1)$$
defines a $``$-isomorphism from $`B(H)N`$ onto $`\widehat{M}_{\widehat{\mathrm{\alpha }}}(M\text{α}N)`$, satisfying
$`\mathrm{\Phi }(\alpha (x))`$ $`=1\alpha (x)`$ for all $`x`$ $`N`$
$`\mathrm{\Phi }(b1)`$ $`=\widehat{\mathrm{\Delta }}\text{op}(b)1`$ for all $`b`$ $`\widehat{M}`$
$`\mathrm{\Phi }(y1)`$ $`=y11`$ for all $`y`$ $`M^{}.`$
In particular $`\mathrm{\Phi }(M\text{α}N)=\widehat{\alpha }(M\text{α}N)`$.
3. When we define
$$\mu =(\sigma \iota )(\iota \alpha ):B(H)NMB(H)N$$
then $`\mu `$ is an action of $`(M,\mathrm{\Delta })`$ on $`B(H)N`$. The unitary $`\mathrm{\Sigma }V^{}\mathrm{\Sigma }1`$ is a $`\mu `$-cocycle and the action $`\gamma `$ of $`(M,\mathrm{\Delta })`$ on $`B(H)N`$ defined by
$$\gamma (z)=(\mathrm{\Sigma }V^{}\mathrm{\Sigma }1)\mu (z)(\mathrm{\Sigma }V\mathrm{\Sigma }1)\text{for all}zB(H)N$$
is isomorphic to the bidual action $`\widehat{\alpha }\widehat{\text{}}`$ of $`(M,\mathrm{\Delta })\widehat{\text{}}\text{op}\widehat{\text{}}\text{op}`$ on $`\widehat{M}_{\widehat{\mathrm{\alpha }}}(M\text{α}N)`$ in the following way:
$$\widehat{\alpha }\widehat{\text{}}(\mathrm{\Phi }(z))=(𝒥\mathrm{\Phi })\gamma (z)\text{for all}zB(H)N.$$
With the help of the biduality theorem Enock and Schwartz were able to prove the following crucial results, which remain true for actions of locally compact quantum groups.
###### Theorem 2.7.
We have
$`(M\text{α}N)^{\widehat{\alpha }}`$ $`=\alpha (N)`$
$`\alpha (N)`$ $`=\{zMN(\iota \alpha )(z)=(\mathrm{\Delta }\iota )(z)\}.`$
## 3 The unitary implementation of a locally compact quantum group action
In this section we will define in a canonical way the unitary implementation of a locally compact quantum group action. This will be a unitary corepresentation of the quantum group, implementing the action and satisfying some other properties. A same kind of result was obtained for Kac algebra actions by Sauvageot in , but the proof of the fact that the implementation is a corepresentation, is wrong. More precisely, Sauvageot’s crucial lemma 4.1 is false. I would like to thank prof. Sauvageot for the discussions on this topic.
We will use a different technique to prove that the implementation is a corepresentation. In the same time we will obtain some interesting results concerning the dual weight on the crossed product $`M\text{α}N`$ given a weight on $`N`$. We will also settle a problem which was left open in .
For integrable actions – and in particular for dual actions – we already obtained an implementation in proposition 2.4, as it was done by Enock and Schwartz. Nevertheless it is desirable to have an implementation without the integrability condition, first of all for reasons of elegance. But, more importantly, one will need this general implementation result in several applications. We refer to the introduction for a discussion.
Fix an action $`\alpha `$ of a locally compact quantum group $`(M,\mathrm{\Delta })`$ on a von Neumann algebra $`N`$. In definition 2.1 and proposition 2.2 we defined the crossed product $`M\text{α}N`$ and the dual action $`\widehat{\alpha }:M\text{α}N\widehat{M}(M\text{α}N)`$. We already observed in proposition 2.5 that $`\widehat{\alpha }`$ is integrable. So we can define the n.s.f. operator valued weight $`T`$ from $`M\text{α}N`$ to $`(M\text{α}N)^{\widehat{\alpha }}`$ by
$$T(z)=(\widehat{\phi }\iota \iota )\widehat{\alpha }(z)\text{for all}z(M\text{α}N)^+.$$
For this, observe that $`\widehat{\alpha }`$ is an action of $`(\widehat{M},\widehat{\mathrm{\Delta }}\text{op})`$ and that $`\widehat{\phi }`$ is the right invariant weight on $`(\widehat{M},\widehat{\mathrm{\Delta }}\text{op})`$. By theorem 2.7 we know that $`(M\text{α}N)^{\widehat{\alpha }}=\alpha (N)`$. So $`T`$ is an operator valued weight from $`M\text{α}N`$ to $`\alpha (N)`$.
With this operator valued weight at hand, we can easily define the dual weights on $`M\text{α}N`$. Nevertheless, to make dual weights a workable tool, we need a concrete GNS-construction for them. The structure of this section is then as follows. First we will restrict the dual weight to a weight for which we can give a GNS-construction (definition 3.4), then we use the restricted weight to obtain the unitary implementation for the action (definition 3.6 and proposition 3.7) and finally we prove that the restricted weight is in fact not a restriction, but equal to the original dual weight (proposition 3.10).
###### Definition 3.1.
Let $`\alpha `$ be an action of $`(M,\mathrm{\Delta })`$ on $`N`$. Denote with $`T`$ the n.s.f. operator valued weight from $`M\text{α}N`$ to $`\alpha (N)`$ given by the formula above. For every n.s.f. weight $`\theta `$ on $`N`$, we define the dual weight $`\stackrel{~}{\theta }`$ on $`M\text{α}N`$ by the formula:
$$\stackrel{~}{\theta }=\theta {}_{}{}^{}\alpha _{}^{1}{}_{}{}^{}T.$$
For the rest of this section we fix a n.s.f. weight $`\theta `$ on $`N`$. One can prove easily the following lemma.
###### Lemma 3.2.
For all $`a𝒩_{\widehat{\phi }}`$ and $`xN`$ we have
$$\stackrel{~}{\theta }\left(\alpha (x^{})(a^{}a1)\alpha (x)\right)=\theta (x^{}x)\widehat{\phi }(a^{}a).$$
###### Proof..
We have
$$\widehat{\alpha }\left(\alpha (x^{})(a^{}a1)\alpha (x)\right)=(1\alpha (x^{}))(\widehat{\mathrm{\Delta }}\text{op}(a^{}a)1)(1\alpha (x)).$$
Choose $`\omega (M\text{α}N)_{}^+`$. Define $`\mu \widehat{M}_{}^+`$ by $`\mu (b)=\omega \left(\alpha (x^{})(b1)\alpha (x)\right)`$ for all $`b\widehat{M}`$. Then
$$T\left(\alpha (x^{})(a^{}a1)\alpha (x)\right),\omega =\widehat{\phi }\left((\iota \mu )\widehat{\mathrm{\Delta }}\text{op}(a^{}a)\right)=\widehat{\phi }(a^{}a)\mu (1)=\widehat{\phi }(a^{}a)\omega (\alpha (x^{}x))$$
by invariance of $`\widehat{\phi }`$. So we may conclude that
$$T\left(\alpha (x^{})(a^{}a1)\alpha (x)\right)=\widehat{\phi }(a^{}a)\alpha (x^{}x).$$
Then the result of the lemma follows immediately. ∎
From now on we will suppose that $`N`$ acts on the GNS-space of the n.s.f. weight $`\theta `$, such that $`(K,\iota ,\mathrm{\Lambda }_\theta )`$ is a GNS-construction for $`\theta `$. We will restrict the weight $`\stackrel{~}{\theta }`$ in the sense of proposition 7.4 of the appendix in order to obtain a concrete GNS-construction. Fix a GNS-construction $`(K_1,\pi _1,\mathrm{\Lambda }_1)`$ for $`\stackrel{~}{\theta }`$. Because of the previous lemma we can define a unique isometry
$$\text{V}:HKK_1\text{such that}\text{V}(\widehat{\mathrm{\Lambda }}(a)\mathrm{\Lambda }_\theta (x))=\mathrm{\Lambda }_1\left((a1)\alpha (x)\right)$$
for all $`a𝒩_{\widehat{\phi }}`$ and $`x𝒩_\theta `$.
Further we define
$$𝒟_0=\mathrm{span}\{(a1)\alpha (x)a𝒩_{\widehat{\phi }},x𝒩_\theta \}.$$
Because we have the isometry V at our disposal there is a well defined linear map
$$\stackrel{~}{\mathrm{\Lambda }}_0:𝒟_0HK:\stackrel{~}{\mathrm{\Lambda }}_0\left((a1)\alpha (x)\right)=\widehat{\mathrm{\Lambda }}(a)\mathrm{\Lambda }_\theta (x)\text{for all}a𝒩_{\widehat{\phi }},x𝒩_\theta .$$
Because $`\mathrm{\Lambda }_1`$ is $`\sigma `$-strong–norm closed, we can close $`\stackrel{~}{\mathrm{\Lambda }}_0`$ for the $`\sigma `$-strong–norm topology, and then we obtain a linear map $`\stackrel{~}{\mathrm{\Lambda }}:𝒟HK`$ satisfying $`𝒟𝒩_{\stackrel{~}{\theta }}`$ and $`\text{V}\stackrel{~}{\mathrm{\Lambda }}(z)=\mathrm{\Lambda }_1(z)`$ for all $`z𝒟`$.
In order to apply proposition 7.4, we need the following lemma.
###### Lemma 3.3.
1. $`𝒟`$ is a weakly dense left ideal in $`M\text{α}N`$.
2. For all $`zM\text{α}N`$ and $`y𝒟`$ we have $`\stackrel{~}{\mathrm{\Lambda }}(zy)=z\stackrel{~}{\mathrm{\Lambda }}(y)`$.
###### Proof..
Choose $`\xi H`$ and $`b𝒯_\phi `$. Let $`(e_i)_{iI}`$ be an orthonormal basis for $`H`$. Choose $`xN`$. Because $`(\mathrm{\Delta }\iota )\alpha (x)=(\iota \alpha )\alpha (x)`$ we have
$$(1\alpha (x))(W1)=(W1)(\iota \alpha )\alpha (x).$$
Hence applying $`\omega _{\xi ,\mathrm{\Lambda }(b)}\iota \iota `$ gives
$$\alpha (x)(\lambda (\omega _{\xi ,\mathrm{\Lambda }(b)})1)=\underset{iI}{}(\lambda (\omega _{e_i,\mathrm{\Lambda }(b)})1)\alpha \left((\omega _{\xi ,e_i}\iota )\alpha (x)\right)$$
in the $`\sigma `$-strong topology.
Choose now $`y𝒩_\theta `$. For every finite subset $`I_0I`$ we have by proposition 7.1 that the element
$`z_{I_0}`$ $`:={\displaystyle \underset{iI_0}{}}(\lambda (\omega _{e_i,\mathrm{\Lambda }(b)})1)\alpha \left((\omega _{\xi ,e_i}\iota )\alpha (x)y\right)`$
belongs to $`𝒟_0`$ and
$`\stackrel{~}{\mathrm{\Lambda }}_0(z_{I_0})`$ $`={\displaystyle \underset{iI_0}{}}\widehat{\mathrm{\Lambda }}\left(\lambda (\omega _{e_i,\mathrm{\Lambda }(b)})\right)(\omega _{\xi ,e_i}\iota )\alpha (x)\mathrm{\Lambda }_\theta (y)`$
$`={\displaystyle \underset{iI_0}{}}J\sigma _{i/2}(b)Je_i(\omega _{\xi ,e_i}\iota )\alpha (x)\mathrm{\Lambda }_\theta (y)`$
$`=(J\sigma _{i/2}(b)J1)(P_{I_0}1)\alpha (x)(\xi \mathrm{\Lambda }_\theta (y))`$
where $`P_{I_0}`$ is the projection on $`\mathrm{span}\{e_iiI_0\}`$. So we get that the net $`(z_{I_0})`$ converges $`\sigma `$-strong to the element
$`z`$ $`:=\alpha (x)(\lambda (\omega _{\xi ,\mathrm{\Lambda }(b)})1)\alpha (y)`$
and the net $`(\stackrel{~}{\mathrm{\Lambda }}(z_{I_0}))`$ converges in norm to
$`(J\sigma _{i/2}(b)J1)\alpha (x)(\xi \mathrm{\Lambda }_\theta (y))`$ $`=\alpha (x)\left(J\sigma _{i/2}(b)J\xi \mathrm{\Lambda }_\theta (y)\right)`$
$`=\alpha (x)\left(\widehat{\mathrm{\Lambda }}(\lambda (\omega _{\xi ,\mathrm{\Lambda }(b)}))\mathrm{\Lambda }_\theta (y)\right)=\alpha (x)\stackrel{~}{\mathrm{\Lambda }}\left((\lambda (\omega _{\xi ,\mathrm{\Lambda }(b)})1)\alpha (y)\right).`$
Then we may conclude that $`z𝒟`$ and
$$\stackrel{~}{\mathrm{\Lambda }}(z)=\alpha (x)\stackrel{~}{\mathrm{\Lambda }}\left((\lambda (\omega _{\xi ,\mathrm{\Lambda }(b)})1)\alpha (y)\right).$$
Because the considered elements $`\lambda (\omega _{\xi ,\mathrm{\Lambda }(b)})`$ form a $`\sigma `$-strong–norm core for $`\widehat{\mathrm{\Lambda }}`$ we conclude that for every $`xN`$ and $`z𝒟`$ we have $`\alpha (x)z𝒟`$ and $`\stackrel{~}{\mathrm{\Lambda }}(\alpha (x)z)=\alpha (x)\stackrel{~}{\mathrm{\Lambda }}(z)`$.
It is easy to prove that for every $`a\widehat{M}`$ and $`z𝒟`$ we have $`(a1)z𝒟`$ and $`\stackrel{~}{\mathrm{\Lambda }}\left((a1)z\right)=(a1)\stackrel{~}{\mathrm{\Lambda }}(z)`$. From this follows the lemma. ∎
We can now apply proposition 7.4.
###### Definition 3.4.
There is a unique n.s.f. weight $`\stackrel{~}{\theta }_0`$ on $`M\text{α}N`$ such that $`𝒩_{\stackrel{~}{\theta }_0}=𝒟`$ and such that $`(HK,\iota ,\stackrel{~}{\mathrm{\Lambda }})`$ is a GNS-construction for $`\stackrel{~}{\theta }_0`$.
Later on we will prove that in fact $`\stackrel{~}{\theta }_0=\stackrel{~}{\theta }`$. This question was left open in the Kac algebra case considered by Sauvageot. In applications the equality $`\stackrel{~}{\theta }_0=\stackrel{~}{\theta }`$ is indispensable, e.g. proposition 5.7 cannot be proved without knowing the GNS-construction of $`\stackrel{~}{\theta }`$, which amounts to the equality $`\stackrel{~}{\theta }_0=\stackrel{~}{\theta }`$.
Let us fix some modular notations.
###### Definition 3.5.
We denote with $`\stackrel{~}{J}`$ and $`\stackrel{~}{}`$ the modular conjugation and modular operator of $`\stackrel{~}{\theta }_0`$ in the GNS-construction $`(HK,\iota ,\stackrel{~}{\mathrm{\Lambda }})`$. We denote with $`\stackrel{~}{\sigma }`$ the modular automorphism group of $`\stackrel{~}{\theta }_0`$ and we put $`\stackrel{~}{T}=\stackrel{~}{J}\stackrel{~}{}^{1/2}`$.
We denote with $`J_\theta `$ and $`_\theta `$ the modular conjugation and modular operator of $`\theta `$ in the GNS-construction $`(K,\iota ,\mathrm{\Lambda }_\theta )`$, and with $`\sigma ^\theta `$ the modular automorphism group of $`\theta `$.
With this notations at hand we will now define the unitary implementation of the action $`\alpha `$. Of course this terminology will only be justified after the proofs of 3.7, 3.12 and 4.4.
###### Definition 3.6.
Define $`U=\stackrel{~}{J}(\widehat{J}J_\theta )`$. Then $`U`$ is a unitary in $`B(HK)`$ and it is called the unitary implementation of $`\alpha `$.
We will first prove the following result.
###### Proposition 3.7.
We have the following formulas:
1. $`\alpha (x)=U(1x)U^{}\text{for all}xN`$.
2. $`\stackrel{~}{\sigma }_t{}_{}{}^{}\alpha =\alpha {}_{}{}^{}\sigma _{t}^{\theta }\text{for all}t`$.
3. $`U(\widehat{J}J_\theta )=(\widehat{J}J_\theta )U^{}`$.
Before we can prove this proposition we need the following lemma.
###### Lemma 3.8.
For all $`y𝒟(\sigma _{i/2}^\theta )`$ we have $`\alpha (y)𝒟(\stackrel{~}{\sigma }_{i/2})`$ and
$$\stackrel{~}{J}\stackrel{~}{\sigma }_{i/2}(\alpha (y))^{}\stackrel{~}{J}=1J_\theta \sigma _{i/2}^\theta (y)^{}J_\theta .$$
###### Proof..
Choose $`a𝒩_{\widehat{\phi }}`$ and $`x𝒩_\theta `$. Then $`xy𝒩_\theta `$ and hence $`(a1)\alpha (x)\alpha (y)𝒩_{\stackrel{~}{\theta }_0}`$ with
$$\stackrel{~}{\mathrm{\Lambda }}\left((a1)\alpha (x)\alpha (y)\right)=\widehat{\mathrm{\Lambda }}(a)\mathrm{\Lambda }_\theta (xy)=(1J_\theta \sigma _{i/2}^\theta (y)^{}J_\theta )\stackrel{~}{\mathrm{\Lambda }}\left((a1)\alpha (x)\right).$$
Because $`𝒟_0`$ is a $`\sigma `$-strong–norm core for $`\stackrel{~}{\mathrm{\Lambda }}`$ we may conclude that for every $`z𝒩_{\stackrel{~}{\theta }_0}`$ we have $`z\alpha (y)𝒩_{\stackrel{~}{\theta }_0}`$ and
$$\stackrel{~}{\mathrm{\Lambda }}(z\alpha (y))=(1J_\theta \sigma _{i/2}^\theta (y)^{}J_\theta )\stackrel{~}{\mathrm{\Lambda }}(z).$$
Then the lemma follows immediately. ∎
###### Proof of proposition 3.7..
Because $`\sigma _{i/2}^\theta (y)^{}=\sigma _{i/2}^\theta (y^{})`$ it follows from the previous lemma that for every $`y𝒟(\sigma _{i/2}^\theta )`$ we have $`\alpha (y)𝒟(\stackrel{~}{\sigma }_{i/2})`$ and
$$\stackrel{~}{\sigma }_{i/2}(\alpha (y))=U(1\sigma _{i/2}^\theta (y))U^{}.$$
Taking the adjoint we may replace $`i/2`$ by $`i/2`$ in the formula above.
Let now $`y𝒟(\sigma _i^\theta )`$. Then we have $`\alpha (y)𝒟(\stackrel{~}{\sigma }_{i/2})`$ and
$$\stackrel{~}{\sigma }_{i/2}(\alpha (y))=U(1\sigma _{i/2}^\theta (y))U^{}.$$
Because $`\sigma _i^\theta (y)𝒟(\sigma _{i/2}^\theta )`$ we also have $`\alpha (\sigma _i^\theta (y))𝒟(\stackrel{~}{\sigma }_{i/2})`$ and
$$\stackrel{~}{\sigma }_{i/2}\left(\alpha (\sigma _i^\theta (y))\right)=U\left(1\sigma _{i/2}^\theta (\sigma _i^\theta (y))\right)U^{}=U(1\sigma _{i/2}^\theta (y))U^{}.$$
So we get $`\stackrel{~}{\sigma }_{i/2}(\alpha (y))=\stackrel{~}{\sigma }_{i/2}\left(\alpha (\sigma _i^\theta (y))\right)`$ and so $`\alpha (y)𝒟(\stackrel{~}{\sigma }_i)`$ with $`\stackrel{~}{\sigma }_i(\alpha (y))=\alpha (\sigma _i^\theta (y))`$. It now follows from the results of \[13, 4.3 and 4.4\] that $`\stackrel{~}{\sigma }_t{}_{}{}^{}\alpha =\alpha {}_{}{}^{}\sigma _{t}^{\theta }`$ for every $`t`$.
But then it follows that for all $`y𝒟(\sigma _{i/2}^\theta )`$ we have $`\stackrel{~}{\sigma }_{i/2}(\alpha (y))=\alpha (\sigma _{i/2}^\theta (y))`$. Combining this with the formula above we get
$$\alpha (\sigma _{i/2}^\theta (y))=\stackrel{~}{\sigma }_{i/2}(\alpha (y))=U(1\sigma _{i/2}^\theta (y))U^{}.$$
By the density of such elements $`\sigma _{i/2}^\theta (y)`$ we get that $`\alpha (x)=U(1x)U^{}`$ for all $`xN`$.
From the definition of $`U`$ follows immediately the final formula we had to prove. ∎
Now we have gathered enough material to prove that $`\stackrel{~}{\theta }_0=\stackrel{~}{\theta }`$. For this we need the following lemma (cfr. \[5, VI.4\]).
###### Lemma 3.9.
Let $`\alpha `$ be an action of $`(M,\mathrm{\Delta })`$ on $`N`$. Let $`\theta _1`$ and $`\theta _2`$ be two $`\delta ^1`$-invariant n.s.f. weights on $`N`$. Then $`[D\theta _2:D\theta _1]_tN^\alpha `$ for all $`t`$.
###### Proof..
Denote with $`M_2`$ the von Neumann algebra of $`2\times 2`$-matrices over $``$. Denote with $`e_{ij}`$ the matrix units. Define
$$\gamma :NM_2MNM_2:\gamma =\alpha \iota .$$
Then $`\gamma `$ is an action of $`(M,\mathrm{\Delta })`$ on $`NM_2`$. Denote with $`\theta `$ the balanced weight on $`NM_2`$ (see e.g. \[26, 3.1\]) given by
$$\theta \left(\begin{array}{cc}x_{11}& x_{12}\\ x_{21}& x_{22}\end{array}\right)=\theta _1(x_{11})+\theta _2(x_{22}).$$
It is immediately clear that $`\theta `$ is $`\delta ^1`$-invariant for the action $`\gamma `$.
Let $`t`$. Denote with $`\mu _t`$ the automorphism of $`M`$ defined by $`\mu _t=\sigma _t^{}{}_{}{}^{}\sigma _{t}^{}{}_{}{}^{}\tau _{t}^{}`$. Here $`(\sigma _t^{})_t`$ denotes the modular automorphism group of $`\psi `$. Then $`\mu _t`$ is implemented by $`Q^{it}=\delta ^{it}\widehat{}^{it}`$. It follows from proposition 2.4 that $`\gamma {}_{}{}^{}\sigma _{t}^{\theta }=(\mu _t\sigma _t^\theta ){}_{}{}^{}\gamma `$ for all $`t`$. In particular we have
$`\alpha ([D\theta _2:D\theta _1]_t)e_{21}`$ $`=\gamma ([D\theta _2:D\theta _1]_te_{21})=\gamma \left(\sigma _t^\theta (1e_{21})\right)`$
$`=(\mu _t\sigma _t^\theta )\gamma (1e_{21})=(\mu _t\sigma _t^\theta )(11e_{21})=1[D\theta _2:D\theta _1]_te_{21}.`$
So we get $`[D\theta _2:D\theta _1]_tN^\alpha `$ for all $`t`$. ∎
Now we can prove the following interesting result. It is important for technical reasons and we will need it in section 5.
###### Proposition 3.10.
Let $`\theta `$ be a n.s.f. weight on $`N`$. Then the weights $`\stackrel{~}{\theta }`$ and $`\stackrel{~}{\theta }_0`$ on $`M\text{α}N`$, defined in 3.1 and 3.4 are equal.
###### Proof..
Recall that the dual action $`\widehat{\alpha }`$ is an action of $`(\widehat{M},\widehat{\mathrm{\Delta }}\text{op})`$ on $`M\text{α}N`$. We claim that the weight $`\stackrel{~}{\theta }_0`$ is $`\widehat{\delta }`$-invariant. Observe that $`\widehat{\delta }^1`$ is the modular element of $`(\widehat{M},\widehat{\mathrm{\Delta }}\text{op})`$ and that is the reason to have $`\widehat{\delta }`$-invariance rather than $`\widehat{\delta }^1`$-invariance.
To prove our claim, choose $`a𝒩_{\widehat{\phi }}`$, $`x𝒩_\theta `$, $`\xi 𝒟(\widehat{\delta }^{1/2})`$ and $`\eta H`$. Then define
$$z:=(\omega _{\xi ,\eta }\iota \iota )\widehat{\alpha }\left((a1)\alpha (x)\right)=\left((\omega _{\xi ,\eta }\iota )\widehat{\mathrm{\Delta }}\text{op}(a)1\right)\alpha (x).$$
It follows from proposition 7.2 of the appendix that
$$(\omega _{\xi ,\eta }\iota )\widehat{\mathrm{\Delta }}\text{op}(a)=(\iota \omega _{\xi ,\eta })\widehat{\mathrm{\Delta }}(a)𝒩_{\widehat{\phi }}\text{and}\widehat{\mathrm{\Lambda }}\left((\omega _{\xi ,\eta }\iota )\widehat{\mathrm{\Delta }}\text{op}(a)\right)=(\iota \omega _{\widehat{\delta }^{1/2}\xi ,\eta })(\widehat{V})\widehat{\mathrm{\Lambda }}(a).$$
So we may conclude that $`z𝒩_{\stackrel{~}{\theta }_0}`$ and
$$\stackrel{~}{\mathrm{\Lambda }}(z)=\left((\iota \omega _{\widehat{\delta }^{1/2}\xi ,\eta })(\widehat{V})1\right)\left(\widehat{\mathrm{\Lambda }}(a)\mathrm{\Lambda }_\theta (x)\right)=\left((\iota \omega _{\widehat{\delta }^{1/2}\xi ,\eta })(\widehat{V})1\right)\stackrel{~}{\mathrm{\Lambda }}\left((a1)\alpha (x)\right).$$
Because $`𝒟_0`$ is a $`\sigma `$-strong–norm core for $`\stackrel{~}{\mathrm{\Lambda }}`$ we conclude that $`(\omega _{\xi ,\eta }\iota \iota )\widehat{\alpha }(y)𝒩_{\stackrel{~}{\theta }_0}`$ for all $`y𝒩_{\stackrel{~}{\theta }_0}`$ and
$$\stackrel{~}{\mathrm{\Lambda }}\left((\omega _{\xi ,\eta }\iota \iota )\widehat{\alpha }(y)\right)=\left((\iota \omega _{\widehat{\delta }^{1/2}\xi ,\eta })(\widehat{V})1\right)\stackrel{~}{\mathrm{\Lambda }}(y).$$
Because $`\widehat{V}`$ is unitary we immediately get that $`\stackrel{~}{\theta }_0`$ is $`\widehat{\delta }`$-invariant.
From proposition 2.5 it follows that $`\stackrel{~}{\theta }`$ is $`\widehat{\delta }`$-invariant. Then we conclude from lemma 3.9 that $`[D\stackrel{~}{\theta }_0:D\stackrel{~}{\theta }]_t(M\text{α}N)^{\widehat{\alpha }}`$ for all $`t`$. So by theorem 2.7 we can take unitaries $`u_tN`$ such that $`[D\stackrel{~}{\theta }_0:D\stackrel{~}{\theta }]_t=\alpha (u_t)`$ for all $`t`$. From the theory of operator valued weights we know that $`\sigma _t^{\stackrel{~}{\theta }}{}_{}{}^{}\alpha =\alpha {}_{}{}^{}\sigma _{t}^{\theta }`$. Because $`([D\stackrel{~}{\theta }_0:D\stackrel{~}{\theta }]_t)`$ is a $`\sigma ^{\stackrel{~}{\theta }}`$-cocycle, we get that $`(u_t)`$ is a $`\sigma ^\theta `$-cocycle. By \[26, 5.1\] we can take a (uniquely determined) n.s.f. weight $`\rho `$ on $`N`$ such that $`[D\rho :D\theta ]_t=u_t`$ for all $`t`$. With $`\rho `$ we can define the n.s.f. weight $`\stackrel{~}{\rho }`$ on $`M\text{α}N`$ in the sense of definition 3.1. Then it follows from the theory of operator valued weights that
$$[D\stackrel{~}{\rho }:D\stackrel{~}{\theta }]_t=\alpha ([D\rho :D\theta ]_t)=\alpha (u_t)=[D\stackrel{~}{\theta }_0:D\stackrel{~}{\theta }]_t$$
for all $`t`$. So $`\stackrel{~}{\rho }=\stackrel{~}{\theta }_0`$. Because $`\stackrel{~}{\theta }_0`$ is a restriction of $`\stackrel{~}{\theta }`$ we get that $`\stackrel{~}{\rho }`$ is a restriction of $`\stackrel{~}{\theta }`$.
Fix $`a_{\widehat{\phi }}^+`$ with $`\widehat{\phi }(a)=1`$. Choose $`x𝒩_\rho `$. Then it follows from lemma 3.2 that $`\alpha (x^{})(a1)\alpha (x)_{\stackrel{~}{\rho }}^+`$ and
$$\stackrel{~}{\rho }\left(\alpha (x^{})(a1)\alpha (x)\right)=\rho (x^{}x).$$
Because $`\stackrel{~}{\rho }`$ is a restriction of $`\stackrel{~}{\theta }`$ we get that $`\alpha (x^{})(a1)\alpha (x)_{\stackrel{~}{\theta }}^+`$ and
$$\stackrel{~}{\theta }\left(\alpha (x^{})(a1)\alpha (x)\right)=\rho (x^{}x).$$
Then it follows from lemma 3.2 that $`\theta (x^{}x)=\rho (x^{}x)`$. This means that $`\rho `$ is a restriction of $`\theta `$.
Further we have, using the theory of operator valued weights in the first equality and proposition 3.7 in the last one,
$$\alpha {}_{}{}^{}\sigma _{t}^{\rho }=\sigma _t^{\stackrel{~}{\rho }}{}_{}{}^{}\alpha =\sigma _t^{\stackrel{~}{\theta }_0}{}_{}{}^{}\alpha =\alpha {}_{}{}^{}\sigma _{t}^{\theta }.$$
So $`\sigma _t^\rho =\sigma _t^\theta `$ for all $`t`$. Because $`\rho `$ is a restriction of $`\theta `$ we may conclude that $`\rho =\theta `$ and then $`\stackrel{~}{\theta }=\stackrel{~}{\rho }=\stackrel{~}{\theta }_0`$. ∎
We want to conclude this section with the proof of the fact that $`UMB(K)`$. First we state the following lemma, which is easily proved because $`\stackrel{~}{\mathrm{\Lambda }}`$ is the closure of $`\stackrel{~}{\mathrm{\Lambda }}_0`$. Recall that $`\stackrel{~}{T}=\stackrel{~}{J}\stackrel{~}{}^{1/2}`$.
###### Lemma 3.11.
Defining $`\widehat{T}=\widehat{J}\widehat{}^{1/2}`$, we have that the linear space
$$\mathrm{span}\{\alpha (x^{})(\eta \mathrm{\Lambda }_\theta (y))x,y𝒩_\theta ,\eta 𝒟(\widehat{T})\}$$
is a core for $`\stackrel{~}{T}`$ and
$$\stackrel{~}{T}\alpha (x^{})(\eta \mathrm{\Lambda }_\theta (y))=\alpha (y^{})(\widehat{T}\eta \mathrm{\Lambda }_\theta (x))$$
for all $`x,y𝒩_\theta `$ and $`\eta 𝒟(\widehat{T})`$.
###### Proposition 3.12.
We have $`UMB(K)`$.
###### Proof..
Let $`t`$. Because $`\widehat{}^{it}`$ implements the automorphism $`\tau _t`$ on $`M`$ we get that $`\mathrm{Ad}\widehat{}^{it}`$ will also leave $`M^{}`$ invariant. So we can define the automorphism group $`(\mu _t)`$ on $`M`$ by
$$\mu _t(x)=J\widehat{}^{it}JxJ\widehat{}^{it}J\text{for all}xM,t.$$
So, for every $`a𝒟(\mu _{i/2})`$ we have $`JaJ\widehat{}^{1/2}\widehat{}^{1/2}J\mu _{i/2}(a)J`$. Further we have
$`\mu _t(R(a))`$ $`=J\widehat{}^{it}J\widehat{J}a^{}\widehat{J}J\widehat{}^{it}J=J\widehat{}^{it}\widehat{J}Ja^{}J\widehat{J}\widehat{}^{it}J`$
$`=\widehat{J}\mu _t(a^{})\widehat{J}=R(\mu _t(a))`$
for all $`t`$ and $`aM`$. Here we used the formula $`\widehat{J}J=\nu ^{i/4}J\widehat{J}`$ stated in the beginning of the paper.
Let now $`a𝒟(\mu _{i/2})`$, $`x,y𝒩_\theta `$ and $`\eta 𝒟(\widehat{T})`$, where $`\widehat{T}=\widehat{J}\widehat{}^{1/2}`$. Then
$`(JaJ1)\stackrel{~}{T}\alpha (x^{})(\eta \mathrm{\Lambda }_\theta (y))=(JaJ1)\alpha (y^{})(\widehat{T}\eta \mathrm{\Lambda }_\theta (x))`$
$`=\alpha (y^{})\left(JaJ\widehat{J}\widehat{}^{1/2}\eta \mathrm{\Lambda }_\theta (x)\right)`$
$`=\alpha (y^{})\left(\widehat{J}JR(a^{})J\widehat{}^{1/2}\eta \mathrm{\Lambda }_\theta (x)\right).`$
Now $`a^{}𝒟(\mu _{i/2})`$ and $`R`$ and $`\mu _t`$ commute. So $`R(a^{})𝒟(\mu _{i/2})`$ and $`\mu _{i/2}(R(a^{}))=R(\mu _{i/2}(a)^{})`$. Then we get
$$JR(a^{})J\widehat{}^{1/2}\widehat{}^{1/2}JR(\mu _{i/2}(a)^{})J.$$
Hence we may conclude that $`JR(\mu _{i/2}(a)^{})J\eta 𝒟(\widehat{T}^{1/2})`$ and
$`(JaJ1)\stackrel{~}{T}\alpha (x^{})(\eta \mathrm{\Lambda }_\theta (y))`$ $`=\alpha (y^{})\left(\widehat{T}JR(\mu _{i/2}(a)^{})J\eta \mathrm{\Lambda }_\theta (x)\right)`$
$`=\stackrel{~}{T}\alpha (x^{})\left(JR(\mu _{i/2}(a)^{})J\eta \mathrm{\Lambda }_\theta (y)\right)`$
$`=\stackrel{~}{T}(JR(\mu _{i/2}(a)^{})J1)\alpha (x^{})(\eta \mathrm{\Lambda }_\theta (y)).`$
Because of the previous lemma we get
$$(JaJ1)\stackrel{~}{T}\stackrel{~}{T}(JR(\mu _{i/2}(a)^{})J1)$$
(3.1)
for all $`a𝒟(\mu _{i/2})`$. By taking the adjoint we get
$$(JR(\mu _{i/2}(a))J1)\stackrel{~}{T}^{}\stackrel{~}{T}^{}(Ja^{}J1)$$
for all $`a𝒟(\mu _{i/2})`$. So for all $`a𝒟(\mu _i)`$
$$(JaJ1)\stackrel{~}{}=(JaJ1)\stackrel{~}{T}^{}\stackrel{~}{T}\stackrel{~}{T}^{}(JR(\mu _{i/2}(a)^{})J1)\stackrel{~}{T}\stackrel{~}{}(J\mu _i(a)J1).$$
Denoting with $`\gamma _t`$ the automorphism $`\mathrm{Ad}\stackrel{~}{}^{it}`$ of $`B(HK)`$ we get that for every $`a𝒟(\mu _i)`$ we have $`JaJ1𝒟(\gamma _i)`$ and $`\gamma _i(JaJ1)=J\mu _i(a)J1`$. Then the results of \[13, 4.3 and 4.4\] allow us to conclude that $`\gamma _t(JaJ1)=J\mu _t(a)J1`$ for every $`t`$ and $`aM`$. This gives
$$(JaJ1)\stackrel{~}{}^{1/2}\stackrel{~}{}^{1/2}(J\mu _{i/2}(a)J1)$$
for all $`a𝒟(\mu _{i/2})`$. Combining this with equation 3.1 we get for every $`a𝒟(\mu _{i/2})`$
$`(JaJ1)\stackrel{~}{T}\stackrel{~}{}^{1/2}`$ $`\stackrel{~}{T}(JR(\mu _{i/2}(a)^{})J1)\stackrel{~}{}^{1/2}`$
$`\stackrel{~}{T}\stackrel{~}{}^{1/2}(JR(a^{})J1)\stackrel{~}{J}(J\widehat{J}a\widehat{J}J1).`$
So we get
$$(JaJ1)\stackrel{~}{J}=\stackrel{~}{J}(J\widehat{J}a\widehat{J}J1)=\stackrel{~}{J}(\widehat{J}JaJ\widehat{J}1)$$
for every $`a𝒟(\mu _{i/2})`$, and hence for every $`aM`$. Rewriting this we get $`(JaJ1)U=U(JaJ1)`$ for every $`aM`$. This gives $`UMB(K)`$. ∎
Finally we want to prove that $`U`$ is a unitary corepresentation of $`(M,\mathrm{\Delta })`$, namely $`(\mathrm{\Delta }\iota )(U)=U_{23}U_{13}`$. This will be done in an indirect way in the next section. Nevertheless the results we use to prove that $`U`$ is a corepresentation are interesting in themselves.
## 4 The unitary implementation is a corepresentation
The main aim of this section is to prove that the unitary implementation $`U`$ is a corepresentation (theorem 4.4). On our way towards the proof of theorem 4.4 we will solve three problems which appear naturally in applications (see section 5 and ). First we will see what happens when we choose a different weight $`\theta `$ on $`N`$, next we will show how $`U`$ changes when the action $`\alpha `$ is deformed with an $`\alpha `$-cocycle and finally we will show that in the presence of a $`\delta ^1`$-invariant weight our implementation agrees with the one of Enock and Schwartz given by proposition 2.4.
In the proof of the first proposition we will make use of Connes’ relative modular theory (see e.g. \[26, 3.11,3.12 and 3.16\]). When $`\theta _i`$ are n.s.f. weights on $`N`$ with GNS-constructions $`(K_i,\pi _i,\mathrm{\Lambda }_i)`$ $`(i=1,2)`$, we denote with $`J_{2,1}`$ the relative modular conjugation, which is a anti-unitary from $`H_1`$ to $`H_2`$. Recall that $`J_{1,2}=J_{2,1}^{}`$. If we denote with $`J_i`$ the modular conjugation of the weight $`\theta _i`$ we have $`J_{2,1}J_1=J_2J_{2,1}`$ and we denote this unitary with $`u`$. Then $`u`$ is the unique unitary from $`K_1`$ to $`K_2`$ which satisfies $`u\pi _1(x)u^{}=\pi _2(x)`$ for all $`xN`$ and which maps the positive cone of $`K_1`$ (determined by the GNS-construction $`(K_1,\pi _1,\mathrm{\Lambda }_1)`$) onto the positive cone of $`K_2`$. We will say that $`u`$ intertwines the two standard representations of $`N`$.
Finally we introduce the one-parameter group $`\sigma ^{2,1}`$ of isometries of $`N`$ given by
$$\sigma _t^{2,1}(x)=[D\theta _2:D\theta _1]_t\sigma _t^{\theta _1}(x)$$
for all $`xN`$ and $`t`$.
###### Proposition 4.1.
Let $`\theta _i`$ be n.s.f. weights on $`N`$ with GNS-constructions $`(K_i,\pi _i,\mathrm{\Lambda }_i)`$ $`(i=1,2)`$. Let $`u`$ be the unitary from $`K_1`$ to $`K_2`$ intertwining the two standard representations of $`N`$. Denote for every $`i=1,2`$ with $`\stackrel{~}{\theta }_i`$ the dual weight of $`\theta _i`$ on $`M\text{α}N`$, with GNS-construction $`(HK_i,\iota \pi _i,\stackrel{~}{\mathrm{\Lambda }}_i)`$. Denote with $`U_iMB(K_i)`$ the unitary implementation of $`\alpha `$ obtained with $`\theta _i`$, as defined in definition 3.6.
Then $`1u`$ is the unitary intertwining the two standard representations of $`M\text{α}N`$. In particular
$$U_2=(1u)U_1(1u^{}).$$
###### Proof..
Let $`a𝒩_{\widehat{\phi }}`$ and $`x𝒩_{\theta _1}`$. Let $`y𝒟(\sigma _{i/2}^{2,1})`$. Then, by \[26, 3.12\], $`xy^{}𝒩_{\theta _2}`$ and
$$\mathrm{\Lambda }_2(xy^{})=J_{2,1}\pi _1(\sigma _{i/2}^{2,1}(y))J_1\mathrm{\Lambda }_1(x).$$
So $`(a1)\alpha (x)\alpha (y)^{}=(a1)\alpha (xy^{})𝒩_{\stackrel{~}{\theta }_2}`$ and
$`\stackrel{~}{\mathrm{\Lambda }}_2\left((a1)\alpha (x)\alpha (y)^{}\right)`$ $`=\widehat{\mathrm{\Lambda }}(a)\mathrm{\Lambda }_2(xy^{})=\left(1J_{2,1}\pi _1(\sigma _{i/2}^{2,1}(y))J_1\right)\left(\widehat{\mathrm{\Lambda }}(a)\mathrm{\Lambda }_1(x)\right)`$
$`=\left(1J_{2,1}\pi _1(\sigma _{i/2}^{2,1}(y))J_1\right)\stackrel{~}{\mathrm{\Lambda }}_1\left((a1)\alpha (x)\right).`$
Because the elements $`(a1)\alpha (x)`$ span a core for $`\stackrel{~}{\mathrm{\Lambda }}_1`$ and because $`\stackrel{~}{\mathrm{\Lambda }}_2`$ is closed (both in the $`\sigma `$-strong–norm topology), we have for all $`z𝒩_{\stackrel{~}{\theta }_1}`$ that $`z\alpha (y)^{}𝒩_{\stackrel{~}{\theta }_2}`$ and
$$\stackrel{~}{\mathrm{\Lambda }}_2(z\alpha (y)^{})=\left(1J_{2,1}\pi _1(\sigma _{i/2}^{2,1}(y))J_1\right)\stackrel{~}{\mathrm{\Lambda }}_1(z).$$
Denoting with $`\stackrel{~}{J}_{2,1}`$ and $`(\stackrel{~}{\sigma }_t^{2,1})`$ the relative modular apparatus of the weights $`\stackrel{~}{\theta }_2`$ and $`\stackrel{~}{\theta }_1`$, it follows from \[26, 3.12\] that $`\alpha (y)𝒟(\stackrel{~}{\sigma }_{i/2}^{2,1})`$ and
$$\stackrel{~}{J}_{2,1}(\iota \pi _1)\left(\stackrel{~}{\sigma }_{i/2}^{2,1}(\alpha (y))\right)\stackrel{~}{J}_1=1J_{2,1}\pi _1(\sigma _{i/2}^{2,1}(y))J_1.$$
Because $`[D\stackrel{~}{\theta }_2:D\stackrel{~}{\theta }_1]_t=\alpha ([D\theta _2:D\theta _1]_t)`$ for every $`t`$ we see that $`\stackrel{~}{\sigma }_t^{2,1}{}_{}{}^{}\alpha =\alpha {}_{}{}^{}\sigma _{t}^{2,1}`$. So we have $`\stackrel{~}{\sigma }_{i/2}^{2,1}(\alpha (y))=\alpha (\sigma _{i/2}^{2,1}(y))`$. Combining this with the equation above we get
$$(\iota \pi _1)\alpha \left(\sigma _{i/2}^{2,1}(y)\right)=\stackrel{~}{J}_{2,1}^{}(\widehat{J}J_{2,1})\left(1\pi _1(\sigma _{i/2}^{2,1}(y))\right)U_1^{}.$$
The last formula is valid for all $`y𝒟(\sigma _{i/2}^{2,1})`$. Because $`U_1`$ implements $`\alpha `$ we may then conclude that $`U_1=\stackrel{~}{J}_{2,1}^{}(\widehat{J}J_{2,1})`$.
Then we get
$$1=U_1U_1^{}=\stackrel{~}{J}_{2,1}^{}(\widehat{J}J_{2,1})(\widehat{J}J_1)\stackrel{~}{J}_1$$
and so $`\stackrel{~}{J}_{2,1}\stackrel{~}{J}_1=1J_{2,1}J_1`$. Now $`u=J_{2,1}J_1`$ and $`\stackrel{~}{J}_{2,1}\stackrel{~}{J}_1`$ is the unitary intertwining the two standard representations of $`M\text{α}N`$. This proves the first claim of the proposition. In particular we get
$$(1u)U_1(1u^{})=(1u)\stackrel{~}{J}_1(\widehat{J}J_1)(1u^{})=\stackrel{~}{J}_2(1u)(\widehat{J}J_1)(1u^{})=\stackrel{~}{J}_2(\widehat{J}J_2)=U_2.$$
This proves the proposition. ∎
In the next proposition we will show how the unitary implementation of an action $`\alpha `$ changes when $`\alpha `$ is deformed with an $`\alpha `$-cocycle.
###### Proposition 4.2.
Let $`\alpha `$ be an action of $`(M,\mathrm{\Delta })`$ on $`N`$ and let $`\text{V}MN`$ be an $`\alpha `$-cocycle in the sense of definition 1.5. Define the action $`\beta `$ of $`(M,\mathrm{\Delta })`$ on $`N`$ by $`\beta (x)=\text{V}\alpha (x)\text{V}^{}`$ for all $`xN`$. If $`\theta `$ is a n.s.f. weight on $`N`$ with GNS-construction $`(K,\iota ,\mathrm{\Lambda }_\theta )`$, the unitary implementations $`U_\alpha `$ and $`U_\beta `$ of $`\alpha `$ and $`\beta `$ obtained with $`\theta `$ satisfy
$$U_\beta =\text{V}U_\alpha (\widehat{J}J_\theta )\text{V}^{}(\widehat{J}J_\theta ).$$
In particular $`U_\beta `$ is a corepresentation if and only if $`U_\alpha `$ is a corepresentation.
###### Proof..
Because $`(\mathrm{\Delta }\iota )(\text{V})=(1\text{V})(\iota \alpha )(\text{V})`$ we have
$$(1\text{V}^{})(W1)(1\text{V})=(W1)(\iota \alpha )(\text{V}^{}).$$
So for every $`\xi ,\eta H`$ and with $`(e_i)_{iI}`$ an orthonormal basis of $`H`$ we have, by applying $`\omega _{\xi ,\eta }\iota \iota `$
$$\text{V}^{}(\lambda (\omega _{\xi ,\eta })1)\text{V}=\underset{iI}{}(\lambda (\omega _{e_i,\eta })1)\alpha \left((\omega _{\xi ,e_i}\iota )(\text{V}^{})\right)$$
in the $`\sigma `$-strong topology. From this it follows that $`\text{V}^{}(a1)\text{V}M\text{α}N`$ for all $`a\widehat{M}`$. But $`\text{V}^{}\beta (x)\text{V}=\alpha (x)`$ for all $`xN`$. So
$$\rho :M\text{β}NM\text{α}N:z\text{V}^{}z\text{V}$$
is a well-defined $``$-homomorphism. By symmetry $`\rho `$ will be surjective and hence it is a $``$-isomorphism.
Consider now the dual weights $`\stackrel{~}{\theta }_\alpha `$ and $`\stackrel{~}{\theta }_\beta `$ on $`M\text{α}N`$ and $`M\text{β}N`$, with canonical GNS-constructions $`(HK,\iota ,\stackrel{~}{\mathrm{\Lambda }}_\alpha )`$ and $`(HK,\iota ,\stackrel{~}{\mathrm{\Lambda }}_\beta )`$. Take $`\xi H`$, $`b𝒯_\phi `$ and $`x𝒩_\theta `$. Then
$$\text{V}^{}(\lambda (\omega _{\xi ,\mathrm{\Lambda }(b)})1)\beta (x)\text{V}=\underset{iI}{}(\lambda (\omega _{e_i,\mathrm{\Lambda }(b)})1)\alpha \left((\omega _{\xi ,e_i}\iota )(\text{V}^{})x\right)$$
in the $`\sigma `$-strong topology. For every finite subset $`I_0`$ of $`I`$ we define
$$z_{I_0}:=\underset{iI_0}{}(\lambda (\omega _{e_i,\mathrm{\Lambda }(b)})1)\alpha \left((\omega _{\xi ,e_i}\iota )(\text{V}^{})x\right).$$
By proposition 7.1 of the appendix we get that $`z_{I_0}`$ belongs to $`𝒩_{\stackrel{~}{\theta }_\alpha }`$ and
$`\stackrel{~}{\mathrm{\Lambda }}_\alpha (z_{I_0})`$ $`={\displaystyle \underset{iI_0}{}}\widehat{\mathrm{\Lambda }}(\lambda (\omega _{e_i,\mathrm{\Lambda }(b)}))(\omega _{\xi ,e_i}\iota )(\text{V}^{})\mathrm{\Lambda }_\theta (x)`$
$`={\displaystyle \underset{iI_0}{}}J\sigma _{i/2}(b)Je_i(\omega _{\xi ,e_i}\iota )(\text{V}^{})\mathrm{\Lambda }_\theta (x)=(J\sigma _{i/2}(b)J1)(P_{I_0}1)\text{V}^{}(\xi \mathrm{\Lambda }_\theta (x))`$
where $`P_{I_0}`$ denotes the projection onto $`\mathrm{span}\{e_iiI_0\}`$. Now define $`z:=\text{V}^{}(\lambda (\omega _{\xi ,\mathrm{\Lambda }(b)})1)\beta (x)\text{V}`$. Then we see that $`z_{I_0}z`$ $`\sigma `$-strong and
$$\stackrel{~}{\mathrm{\Lambda }}_\alpha (z_{I_0})(J\sigma _{i/2}(b)J1)\text{V}^{}(\xi \mathrm{\Lambda }_\theta (x))\text{in norm}.$$
So we get that $`z𝒩_{\stackrel{~}{\theta }_\alpha }`$ and
$`\stackrel{~}{\mathrm{\Lambda }}_\alpha (z)`$ $`=(J\sigma _{i/2}(b)J1)\text{V}^{}(\xi \mathrm{\Lambda }_\theta (x))`$
$`=\text{V}^{}\left(J\sigma _{i/2}(b)J\xi \mathrm{\Lambda }_\theta (x)\right)=\text{V}^{}\stackrel{~}{\mathrm{\Lambda }}_\beta \left((\lambda (\omega _{\xi ,\mathrm{\Lambda }(b)})1)\beta (x)\right).`$
Because the elements $`(\lambda (\omega _{\xi ,\mathrm{\Lambda }(b)})1)\beta (x)`$ span a core for $`\stackrel{~}{\mathrm{\Lambda }}_\beta `$ we have $`\rho (y)𝒩_{\stackrel{~}{\theta }_\alpha }`$ for every $`y𝒩_{\stackrel{~}{\theta }_\beta }`$ and $`\stackrel{~}{\mathrm{\Lambda }}_\alpha (\rho (y))=\text{V}^{}\stackrel{~}{\mathrm{\Lambda }}_\beta (y)`$ in that case.
By symmetry $`\rho (y)𝒩_{\stackrel{~}{\theta }_\alpha }`$ if and only if $`y𝒩_{\stackrel{~}{\theta }_\beta }`$. But then it is clear that $`\stackrel{~}{J}_\beta =\text{V}\stackrel{~}{J}_\alpha \text{V}^{}`$ and so
$$U_\beta =\stackrel{~}{J}_\beta (\widehat{J}J_\theta )=\text{V}\stackrel{~}{J}_\alpha \text{V}^{}(\widehat{J}J_\theta )=\text{V}U_\alpha (\widehat{J}J_\theta )\text{V}^{}(\widehat{J}J_\theta ).$$
Now suppose that $`U_\alpha `$ is a corepresentation, meaning that $`(\mathrm{\Delta }\iota )(U_\alpha )=U_{\alpha \mathrm{\hspace{0.17em}23}}U_{\alpha \mathrm{\hspace{0.17em}13}}`$. Then
$$(\mathrm{\Delta }\iota )(U_\beta )=(\mathrm{\Delta }\iota )(\text{V})U_{\alpha \mathrm{\hspace{0.17em}23}}U_{\alpha \mathrm{\hspace{0.17em}13}}(\mathrm{\Delta }\iota )(RL_\theta )(\text{V})$$
where $`L_\theta `$ is the $``$-anti-isomorphism from $`N`$ to $`N^{}`$ defined by $`L_\theta (x)=J_\theta x^{}J_\theta `$ for all $`xN`$. Then we can compute
$`(\mathrm{\Delta }\iota )(U_\beta )`$ $`=\text{V}_{23}(\iota \alpha )(\text{V})U_{\alpha \mathrm{\hspace{0.17em}23}}U_{\alpha \mathrm{\hspace{0.17em}13}}(RRL_\theta )(\mathrm{\Delta }\text{op}\iota )(\text{V})`$
$`=\text{V}_{23}U_{\alpha \mathrm{\hspace{0.17em}23}}\text{V}_{13}U_{\alpha \mathrm{\hspace{0.17em}23}}^{}U_{\alpha \mathrm{\hspace{0.17em}23}}U_{\alpha \mathrm{\hspace{0.17em}13}}(RRL_\theta )(\text{V}_{13}(\iota \alpha )(\text{V})_{213})`$
$`=\text{V}_{23}U_{\alpha \mathrm{\hspace{0.17em}23}}\text{V}_{13}U_{\alpha \mathrm{\hspace{0.17em}13}}(\widehat{J}\widehat{J}J_\theta )U_{\alpha \mathrm{\hspace{0.17em}13}}\text{V}_{23}^{}U_{\alpha \mathrm{\hspace{0.17em}13}}^{}\text{V}_{13}^{}(\widehat{J}\widehat{J}J_\theta )`$
$`=\text{V}_{23}U_{\alpha \mathrm{\hspace{0.17em}23}}\text{V}_{13}(\widehat{J}\widehat{J}J_\theta )\text{V}_{23}^{}(\widehat{J}\widehat{J}J_\theta )(\widehat{J}\widehat{J}J_\theta )U_{\alpha \mathrm{\hspace{0.17em}13}}^{}\text{V}_{13}^{}(\widehat{J}\widehat{J}J_\theta )`$
$`=\text{V}_{23}U_{\alpha \mathrm{\hspace{0.17em}23}}\left((\widehat{J}J_\theta )\text{V}^{}(\widehat{J}J_\theta )\right)_{23}\text{V}_{13}U_{\alpha \mathrm{\hspace{0.17em}13}}\left((\widehat{J}J_\theta )\text{V}^{}(\widehat{J}J_\theta )\right)_{13}`$
$`=U_{\beta \mathrm{\hspace{0.17em}23}}U_{\beta \mathrm{\hspace{0.17em}13}}.`$
So, when $`U_\alpha `$ is a corepresentation then $`U_\beta `$ is a corepresentation. By symmetry also the converse implication holds. ∎
In proposition 2.4 we saw how to construct, with the methods of Enock and Schwartz, an implementation of an action $`\alpha `$ in the presence of a $`\delta ^1`$-invariant weight. We will show now that this implementation coincides with the unitary implementation given in definition 3.6.
###### Proposition 4.3.
Let $`\alpha `$ be an action of $`(M,\mathrm{\Delta })`$ on $`N`$. Let $`\theta `$ be a n.s.f. and $`\delta ^1`$-invariant weight on $`N`$ with GNS-construction $`(K,\iota ,\mathrm{\Lambda }_\theta )`$. When $`V_\theta `$ is the unitary defined in proposition 2.4 and when $`U`$ is the unitary implementation of $`\alpha `$ defined in definition 3.6, then $`U=V_\theta `$.
###### Proof..
Recall that
$$(\omega _{\xi ,\eta }\iota )(V_\theta )\mathrm{\Lambda }_\theta (x)=\mathrm{\Lambda }_\theta \left((\omega _{\delta ^{1/2}\xi ,\eta }\iota )\alpha (x)\right)$$
for all $`\xi 𝒟(\delta ^{1/2})`$ and $`\eta H`$. Because the positive operators $`\delta `$ and $`\widehat{}`$ strongly commute, we can define the closure $`Q`$ of the product $`\delta \widehat{}`$. Denoting with $`\chi _𝒰`$ the characteristic function of a subset $`𝒰`$ we consider the following subspace of $`H`$.
$$𝒟_0=\underset{n,m}{}\chi _{]\frac{1}{n},n[}(\delta )\chi _{]\frac{1}{m},m[}(\widehat{})H.$$
Let now $`\xi H`$, $`\eta 𝒟_0`$, $`x𝒯_\theta `$ and $`y𝒩_\theta 𝒩_\theta ^{}`$. Put again $`\widehat{T}=\widehat{J}\widehat{}^{1/2}`$ and $`\stackrel{~}{T}=\stackrel{~}{J}\stackrel{~}{}^{1/2}`$. Then
$`(\theta _\xi ^{}1)\stackrel{~}{T}\alpha (x^{})(\eta \mathrm{\Lambda }_\theta (y))`$ $`=(\theta _\xi ^{}1)\alpha (y^{})(\widehat{T}\eta \mathrm{\Lambda }_\theta (x))=\mathrm{\Lambda }_\theta \left((\omega _{\widehat{T}\eta ,\xi }\iota )\alpha (y^{})x\right)`$
$`=J_\theta \sigma _{i/2}^\theta (x)^{}J_\theta \mathrm{\Lambda }_\theta \left((\omega _{\widehat{T}\eta ,\xi }\iota )\alpha (y^{})\right)`$
$`=J_\theta \sigma _{i/2}^\theta (x)^{}J_\theta (\omega _{\delta ^{1/2}\widehat{T}\eta ,\xi }\iota )(V_\theta )\mathrm{\Lambda }_\theta (y^{})`$
$`=(\theta _\xi ^{}1)(1J_\theta \sigma _{i/2}^\theta (x)^{}J_\theta )V_\theta (\delta ^{1/2}\widehat{T}\eta \mathrm{\Lambda }_\theta (y^{})).`$
Now
$$\delta ^{1/2}\widehat{T}\eta =\delta ^{1/2}\widehat{J}\widehat{}^{1/2}\eta =\widehat{J}\delta ^{1/2}\widehat{}^{1/2}\eta =\widehat{J}Q^{1/2}\eta .$$
So we may conclude that
$$\stackrel{~}{T}\alpha (x^{})(\eta \mathrm{\Lambda }_\theta (y))=(1J_\theta \sigma _{i/2}^\theta (x)^{}J_\theta )V_\theta (\widehat{J}J_\theta )(Q^{1/2}_\theta ^{1/2})(\eta \mathrm{\Lambda }_\theta (y))$$
for all $`\eta 𝒟_0`$, $`x𝒯_\theta `$ and $`y𝒩_\theta 𝒩_\theta ^{}`$. Because $`\stackrel{~}{T}`$ is closed we can conclude that $`\eta \mathrm{\Lambda }_\theta (y)𝒟(\stackrel{~}{T})`$ and
$$\stackrel{~}{T}(\eta \mathrm{\Lambda }_\theta (y))=V_\theta (\widehat{J}J_\theta )(Q^{1/2}_\theta ^{1/2})(\eta \mathrm{\Lambda }_\theta (y))$$
for all $`\eta 𝒟_0`$ and $`y𝒩_\theta 𝒩_\theta ^{}`$. Because $`𝒟_0`$ is a core for $`Q^{1/2}`$ and $`\mathrm{\Lambda }_\theta (𝒩_\theta 𝒩_\theta ^{})`$ for $`_\theta ^{1/2}`$ we get
$$V_\theta (\widehat{J}J_\theta )(Q^{1/2}_\theta ^{1/2})\stackrel{~}{T}.$$
(4.1)
We now claim that $`(Q^{it}_\theta ^{it})\stackrel{~}{T}=\stackrel{~}{T}(Q^{it}_\theta ^{it})`$ for every $`t`$. Together with the fact that $`𝒟(Q^{1/2}_\theta ^{1/2})𝒟(\stackrel{~}{T})`$ this leads to the conclusion that $`𝒟(Q^{1/2}_\theta ^{1/2})`$ is a core for $`\stackrel{~}{T}`$. Then we get that the inclusion in equation 4.1 is in fact an equality. Uniqueness of the polar decomposition gives us $`V_\theta (\widehat{J}J_\theta )=\stackrel{~}{J}`$ and so $`V_\theta =U`$.
So we only have to prove our claim. For this choose $`x,y𝒩_\theta `$ and $`\xi 𝒟(\widehat{T})`$. Then using proposition 2.4 we get
$`(Q^{it}_\theta ^{it})\stackrel{~}{T}\alpha (x^{})(\xi \mathrm{\Lambda }_\theta (y))`$ $`=(Q^{it}_\theta ^{it})\alpha (y^{})(\widehat{T}\xi \mathrm{\Lambda }_\theta (x))`$
$`=(Q^{it}_\theta ^{it})V_\theta (1y^{})V_\theta ^{}(\widehat{T}\xi \mathrm{\Lambda }_\theta (x))`$
$`=V_\theta (1\sigma _t^\theta (y^{}))V_\theta ^{}(Q^{it}\widehat{T}\xi _\theta ^{it}\mathrm{\Lambda }_\theta (x))`$
because $`Q^{it}_\theta ^{it}`$ and $`V_\theta `$ commute. Now observe that $`Q^{it}`$ and $`\widehat{T}`$ commute, so that $`Q^{it}\xi 𝒟(\widehat{T})`$ and
$`(Q^{it}_\theta ^{it})\stackrel{~}{T}\alpha (x^{})(\xi \mathrm{\Lambda }_\theta (y))`$ $`=\alpha (\sigma _t^\theta (y)^{})(\widehat{T}Q^{it}\xi \mathrm{\Lambda }_\theta (\sigma _t^\theta (x)))`$
$`=\stackrel{~}{T}\alpha (\sigma _t^\theta (x^{}))(Q^{it}\xi \mathrm{\Lambda }_\theta (\sigma _t^\theta (y)))`$
$`=\stackrel{~}{T}(Q^{it}_\theta ^{it})\alpha (x^{})(\xi \mathrm{\Lambda }_\theta (y)).`$
From this immediately follows our claim, and then the proof of the proposition is complete. ∎
With all these results at hand we can now prove the following theorem.
###### Theorem 4.4.
The unitary implementation $`U`$ of an action $`\alpha `$ of $`(M,\mathrm{\Delta })`$ on $`N`$ is a corepresentation in the sense that $`(\mathrm{\Delta }\iota )(U)=U_{23}U_{13}`$.
###### Proof..
Consider the bidual action $`\widehat{\alpha }\widehat{\text{}}`$ of $`(M,\mathrm{\Delta })\widehat{\text{}}\text{op}\widehat{\text{}}\text{op}`$ on $`\widehat{M}_{\widehat{\alpha }}(M\text{α}N)`$. Let $`\theta `$ be a n.s.f. weight on $`N`$ and denote with $`\stackrel{~}{\theta }\stackrel{~}{\text{}}`$ the bidual weight on $`\widehat{M}_{\widehat{\alpha }}(M\text{α}N)`$. It follows from proposition 2.5 that $`\stackrel{~}{\theta }\stackrel{~}{\text{}}`$ is a $`J\delta J`$-invariant weight for the action $`\widehat{\alpha }\widehat{\text{}}`$. With the notation of theorem 2.6 we define $`\rho =\stackrel{~}{\theta }\stackrel{~}{\text{}}{}_{}{}^{}\mathrm{\Phi }`$. Then $`\rho `$ will be a n.s.f. and $`\delta ^1`$-invariant weight on $`B(H)N`$ for the action $`\gamma `$ of $`(M,\mathrm{\Delta })`$ on $`B(H)N`$. Combining proposition 4.3 and proposition 2.4 the unitary implementation of $`\gamma `$ constructed with the weight $`\rho `$ is a corepresentation. By proposition 4.2 the unitary implementation of $`\mu :=(\sigma \iota )(\iota \alpha )`$ constructed with $`\rho `$ is a corepresentation as well. Then it follows from proposition 4.1 that the unitary implementation $`U_\mu `$ of $`\mu `$ constructed with the n.s.f. weight $`\mathrm{Tr}\theta `$ on $`B(H)N`$ will be a corepresentation. Here $`\mathrm{Tr}`$ denotes the usual trace on $`B(H)`$.
Represent $`N`$ on the GNS-space of $`\theta `$ such that $`(K,\iota ,\mathrm{\Lambda }_\theta )`$ is a GNS-construction for $`\theta `$. Let $`(H_{\mathrm{Tr}},\pi _{\mathrm{Tr}},\mathrm{\Lambda }_{\mathrm{Tr}})`$ be a GNS-construction for $`\mathrm{Tr}`$. Then we have a canonical GNS-construction $`(H_{\mathrm{Tr}}K,\pi _{\mathrm{Tr}}\iota ,\mathrm{\Lambda }_{\mathrm{Tr}\theta })`$ for $`\mathrm{Tr}\theta `$. With this we construct the GNS-construction $`(HH_{\mathrm{Tr}}K,\iota \pi _{\mathrm{Tr}}\iota ,\stackrel{~}{\mathrm{\Lambda }}_{\mathrm{Tr}\theta })`$ of the dual weight $`(\mathrm{Tr}\theta )\stackrel{~}{\text{}}`$ on $`M\text{μ}(B(H)N)`$. Denote with $`\stackrel{~}{T}_{\mathrm{Tr}\theta }=\stackrel{~}{J}_{\mathrm{Tr}\theta }\stackrel{~}{}_{\mathrm{Tr}\theta }^{1/2}`$ the modular operator of this dual weight. As before we denote with $`\stackrel{~}{T}=\stackrel{~}{J}\stackrel{~}{}^{1/2}`$ the modular operator of the weight $`\stackrel{~}{\theta }`$ on $`M\text{α}N`$ with GNS-construction $`(HK,\iota ,\stackrel{~}{\mathrm{\Lambda }})`$. It is an easy exercise to check that
$$\mathrm{\Sigma }_{12}\stackrel{~}{T}_{\mathrm{Tr}\theta }\mathrm{\Sigma }_{12}=J_{\mathrm{Tr}}\stackrel{~}{T}$$
where $`\mathrm{\Sigma }_{12}`$ flips the first two legs of $`HH_{\mathrm{Tr}}K`$. By uniqueness of the polar decomposition we get
$$\mathrm{\Sigma }_{12}\stackrel{~}{J}_{\mathrm{Tr}\theta }\mathrm{\Sigma }_{12}=J_{\mathrm{Tr}}\stackrel{~}{J}$$
and hence $`\mathrm{\Sigma }_{12}U_\mu \mathrm{\Sigma }_{12}=1U`$. Because $`U_\mu `$ is a corepresentation, also $`U`$ will be a corepresentation. ∎
## 5 Subfactors and inclusions of von Neumann algebras
It is well known that there is an important link between irreducible, depth 2 inclusions of factors and quantum groups. After a conjecture of Ocneanu the first result in this direction was proved by David in , Longo in and Szymanski in . They were able to prove that every irreducible, depth 2 inclusion of $`II_1`$-factors with finite index has the form $`N^\alpha N`$, where $`N`$ is a $`II_1`$-factor and $`\alpha `$ is an action of a finite Kac algebra (i.e. a finite dimensional locally compact quantum group, or a finite dimensional Hopf $``$-algebra with positive invariant integral). The restriction on type and index has been removed by Enock and Nest in and . There does not appear a finite quantum group but an arbitrary locally compact quantum group.
The theory of Enock and Nest is quite technical, but the results are deep and beautiful. They are important in themselves and serve as a motivation for the concept of a locally compact quantum group.
Before we describe their result we have to explain a little bit the basic theory of infinite index inclusions of factors or von Neumann algebras. So, let us look at an inclusion $`N_0N_1`$ of von Neumann algebras. In this most general setting one can perform the well known basic construction of Jones. For this we have to choose a n.s.f. weight $`\theta `$ on $`N_1`$ and represent $`N_1`$ on the GNS-space of $`\theta `$. Denote with $`J_\theta `$ the modular conjugation of $`\theta `$. Then we define $`N_2=J_\theta N_0^{}J_\theta `$. Because $`N_1^{}=J_\theta N_1J_\theta `$ we have $`N_0N_1N_2`$ and this inclusion of three von Neumann algebras is called the basic construction. One can continue in the same way and represent $`N_2`$ on the GNS-space of some n.s.f. weight. Then we obtain the von Neumann algebra $`N_3`$. Going on we get a tower of von Neumann algebras
$$N_0N_1N_2N_3\mathrm{}$$
which is called the Jones tower.
But there is more. In the theory of inclusions of $`II_1`$-factors an important role is played by conditional expectations. In the more general theory being described now, this role will be taken over by operator valued weights. Before we can explain this, and also because we need it in the proof of theorem 5.3, we have to explain Connes’ spatial modular theory. For this we refer to e.g. \[26, §7\] and \[28, §III\].
Suppose that $`N`$ is a von Neumann algebra acting on a Hilbert space $`K`$. Let $`\phi `$ be a n.s.f. weight on $`N`$ and $`\psi `$ a n.s.f. weight on $`N^{}`$. Let $`(K_\psi ,\pi _\psi ,\mathrm{\Lambda }_\psi )`$ be a GNS-construction for $`\psi `$. For every $`\xi K`$ we define the densely defined operator $`R^\psi (\xi )`$ with domain $`\mathrm{\Lambda }_\psi (𝒩_\psi )K_\psi `$ and range in $`K`$ by $`R^\psi (\xi )\mathrm{\Lambda }_\psi (x)=x\xi `$ for all $`x𝒩_\psi `$. When $`\xi K`$ we can define an operator $`\mathrm{\Theta }^\psi (\xi )`$ in the extended positive part of $`B(K)`$ by
$$\omega _\eta ,\mathrm{\Theta }^\psi (\xi )=\{\begin{array}{cc}R^\psi (\xi )^{}\eta ^2\hfill & \text{if}\eta 𝒟(R^\psi (\xi )^{})\hfill \\ +\mathrm{}\hfill & \text{else}\hfill \end{array}.$$
In fact $`\mathrm{\Theta }^\psi (\xi )=R^\psi (\xi )R^\psi (\xi )^{}+\mathrm{}P`$ where $`P`$ is the projection onto the orthogonal complement of $`𝒟(R^\psi (\xi )^{})`$. Then one can prove that $`\mathrm{\Theta }^\psi (\xi )`$ belongs to $`N^+\text{ext}`$ and it is possible to define a strictly positive, self-adjoint operator $`{\displaystyle \frac{d\phi }{d\psi }}`$ on $`K`$ such that
$$\omega _\xi ,\frac{d\phi }{d\psi }=\phi ,\mathrm{\Theta }^\psi (\xi )$$
for all $`\xi K`$. Here we used the extension of the weight $`\phi `$ to the extended positive part of $`N`$. The operator $`{\displaystyle \frac{d\phi }{d\psi }}`$ is called the spatial derivative of $`\phi `$ with respect to $`\psi `$.
So, let $`N_0N_1`$ be an inclusion of von Neumann algebras and $`T_1`$ a n.s.f. operator valued weight from $`N_1`$ to $`N_0`$. Represent again $`N_1`$ on the GNS-space of a n.s.f. weight $`\theta `$. Let $`N_0N_1N_2`$ be the basic construction. Then there exists a unique n.s.f. operator valued weight $`T_2`$ from $`N_2`$ to $`N_1`$ such that
$$\frac{d(\mu {}_{}{}^{}T_{2}^{})}{d\nu ^{}}=\frac{d\mu }{d\left((\nu {}_{}{}^{}T_{1}^{})^{}\right)}$$
for all n.s.f. weights $`\mu `$ on $`N_1`$ and $`\nu `$ on $`N_0`$. Here we denote with $`\eta ^{}`$ the n.s.f. weight on either $`N_2^{}=J_\theta N_0J_\theta `$ or $`N_1^{}=J_\theta N_1J_\theta `$, given by the formula $`\eta ^{}(x)=\eta (J_\theta xJ_\theta )`$ for all positive $`x`$, whenever $`\eta `$ is a n.s.f. weight on either $`N_0`$ or $`N_1`$. The existence of $`T_2`$ follows from \[26, 12.11\]. One can continue in the same way and construct n.s.f. operator valued weights $`T_i`$ from $`N_i`$ to $`N_{i1}`$ anywhere in the Jones tower.
Next recall that an inclusion of von Neumann algebras $`N_0N_1`$ is said to be
* irreducible, when $`N_1N_0^{}=`$.
* of depth 2, when $`N_1N_0^{}N_2N_0^{}N_3N_0^{}`$ is the basic construction.
Finally we describe the notion of regularity as it was introduced by Enock and Nest in \[7, 11.12\]. Let $`N_0N_1`$ be an inclusion of von Neumann algebras. Suppose that $`T_1`$ is a n.s.f. operator valued weight from $`N_1`$ to $`N_0`$. Let $`N_0N_1N_2N_3\mathrm{}`$ be the Jones tower and construct as above the operator valued weights $`T_2`$ from $`N_2`$ to $`N_1`$ and $`T_3`$ from $`N_3`$ to $`N_2`$. Then $`T_1`$ is called regular when the restrictions of $`T_2`$ to $`N_2N_0^{}`$ and of $`T_3`$ to $`N_3N_1^{}`$ are both semifinite.
Then we can give the main result of Enock and Nest. Recall that for a locally compact quantum group $`(M,\mathrm{\Delta })`$ we denoted with $`(M,\mathrm{\Delta })^{}`$ the commutant locally compact quantum group, as described in the introduction.
###### Theorem 5.1 (Enock and Nest).
Let $`N_0N_1`$ be an irreducible, depth 2 inclusion of factors and let $`T_1`$ be a regular n.s.f. operator valued weight from $`N_1`$ to $`N_0`$. Then the von Neumann algebra $`M=N_3N_1^{}`$ can be given the structure of a locally compact quantum group $`(M,\mathrm{\Delta })`$, such that there exists an outer action $`\alpha `$ of $`(M,\mathrm{\Delta })^{}`$ on $`N_1`$ satisfying $`N_0=N_1^\alpha `$ and such that the inclusions $`N_0N_1N_2`$ and $`N_1^\alpha \alpha (N_1)M^{}\text{α}N_1`$ are isomorphic.
The definition of an outer action is given in definition 5.5. Further we want to mention that in it is proved that $`(M,\mathrm{\Delta })`$, together with invariant weights and antipode, is in fact a Woronowicz algebra. But it should be stressed that there is a small mistake in the proof that the Haar weight is invariant under the scaling group, so that in fact $`(M,\mathrm{\Delta })`$ is an arbitrary locally compact quantum group, possibly with scaling constant different from 1.
The main aim of this section is to clarify the conditions of Enock and Nest’s theorem (in particular the regularity condition) and to prove a converse result: when $`\alpha `$ is an integrable and outer action of $`(M,\mathrm{\Delta })`$ on $`N`$, then the inclusion $`N^\alpha N`$ is irreducible, of depth 2 and the operator valued weight $`(\psi \iota )\alpha `$ from $`N`$ to $`N^\alpha `$ is regular. The same result is stated in \[7, 11.14\] for the special case of dual Kac algebra actions, but the proof is incomplete. The crucial point, our proposition 5.7 identifying two operator valued weights, is not proved in . We also remark that it will follow from corollary 5.6 that the actions appearing in Enock and Nest’s theorem are integrable. Further we refer to section 6 for the link between outer and minimal actions.
First of all we study the following problem. Let $`\alpha `$ be an action of $`(M,\mathrm{\Delta })`$ on $`N`$. Let $`N^\alpha NN_2`$ be the basic construction. When $`(M,\mathrm{\Delta })`$ is finite dimensional, it is known that $`N_2`$ is a quotient of the crossed product $`M\text{α}N`$ (a proof can be found in \[14, 4.1.3\], but the result was undoubtedly known before). More precisely, there exists a surjective $``$-homomorphism $`\rho `$ from $`M\text{α}N`$ to $`N_2`$ sending $`\alpha (x)`$ to $`x`$ for all $`xN`$. So, when $`M\text{α}N`$ is a factor, the inclusion $`N^\alpha \alpha (N)M\text{α}N`$ is the basic construction. More specifically, when $`N`$ is a $`II_1`$-factor and $`\alpha `$ is an outer action (or equivalently a free action) of a finite group $`G`$ on $`N`$ it is well known that the crossed product $`G\text{α}N`$ can be identified with $`(N\{u_ttG\})^{\prime \prime }`$, where $`N`$ is represented standardly and $`(u_t)_{tG}`$ is the canonical implementation of $`\alpha `$. This can be found in e.g. .
More generally we look at the following problem. Suppose that a locally compact group $`G`$ acts on a von Neumann algebra $`N`$ with action $`\alpha `$. Then we can construct the crossed product $`G\text{α}N`$ as follows. We represent $`N`$ on a Hilbert space $`K`$ and define operators on $`L^2(G)KL^2(G,K)`$ by putting
$`(\alpha (x)\xi )(g)`$ $`=\alpha _{g^1}(x)\xi (g)`$ $`\text{for all}gG\text{and}xN,\xi L^2(G,K)`$
$`(\lambda _g\xi )(h)`$ $`=\xi (g^1h)`$ $`\text{for all}g,hG\text{and}\xi L^2(G,K).`$
Then we define $`G\text{α}N=\left(\alpha (N)\{\lambda _ggG\}\right)^{\prime \prime }`$. But, when we represent $`N`$ standardly on $`K`$ and denote with $`(u_g)_{gG}`$ the canonical unitary implementation of $`\alpha `$, we can also define
$$N_2=\left(N\{u_ggG\}\right)^{\prime \prime }.$$
Purely algebraicly one would expect to be able to define a $``$-homomorphism $`\rho :G\text{α}NN_2`$ satisfying $`\rho (\alpha (x))=x`$ for all $`xN`$ and $`\rho (\lambda _g)=u_g`$ for all $`gG`$. When the group $`G`$ is finite, this can be done easily. In theorem 5.3 we will prove that the construction of such a $`\rho `$ is possible if and only if the action is integrable, and this will be proved for arbitrary locally compact quantum group actions. To see the link with the group case, recall that now the role of the regular representation $`(\lambda _g)`$ is taken over by $`\lambda (\omega )=(\omega \iota )(W)`$ for all $`\omega M_{}`$. So we work in fact with the regular representation of the $`L^1`$-functions.
Before we come to the proof of our main theorem 5.3 we characterize the basic construction $`N_2=J_\theta (N^\alpha )^{}J_\theta `$ in terms of the unitary implementation of $`\alpha `$.
###### Proposition 5.2.
Let $`\alpha `$ be an action of $`(M,\mathrm{\Delta })`$ on $`N`$. Fix a n.s.f. weight $`\theta `$ on $`N`$ and let $`N`$ act on the GNS-space of $`\theta `$. Let $`U`$ be the unitary implementation of $`\alpha `$ obtained with $`\theta `$. Let $`N_2=J_\theta (N^\alpha )^{}J_\theta `$ be the basic construction from $`N^\alpha N`$. Then
$$N_2=(N\{(\omega \iota )(U)\omega M_{}\})^{\prime \prime }.$$
###### Proof..
Because $`N_2=J_\theta (N^\alpha )^{}J_\theta `$ we get easily that
$$N_2^{}=N^{}J_\theta \{(\omega \iota )(U^{})\omega M_{}\}^{}J_\theta .$$
But $`(\widehat{J}J_\theta )U^{}(\widehat{J}J_\theta )=U`$, so that we have
$$N_2^{}=N^{}\{(\omega \iota )(U)\omega M_{}\}^{}.$$
Because $`U`$ is a corepresentation the $`\sigma `$-strong closure of $`\{(\omega \iota )(U)\omega M_{}\}`$ is self-adjoint and then the result follows. ∎
Then we prove the following result.
###### Theorem 5.3.
Let $`\alpha `$ be an action of $`(M,\mathrm{\Delta })`$ on $`N`$. Fix a n.s.f. weight $`\theta `$ on $`N`$ and let $`N`$ act on the GNS-space of $`\theta `$. Let $`U`$ be the unitary implementation of $`\alpha `$ obtained with $`\theta `$. Let $`N_2=J_\theta (N^\alpha )^{}J_\theta `$ be the basic construction from $`N^\alpha N`$. Then the following statements are equivalent.
* There exists a normal and surjective $``$-homomorphism $`\rho :M\text{α}NN_2`$ satisfying
$$\rho (\alpha (x))=x\text{for all}xN\text{and}\rho \left((\omega \iota )(W)1\right)=(\omega \iota )(U^{})\text{for all}\omega M_{}.$$
* The action $`\alpha `$ is integrable.
###### Proof of the first implication..
Let us first suppose the first statement is valid. Because $`\mathrm{Ker}\rho `$ is a $`\sigma `$-strong closed, two-sided ideal of $`M\text{α}N`$ we can take a central projection $`PM\text{α}N`$ such that $`\mathrm{Ker}\rho =(M\text{α}N)(1P)`$. Let $`\rho _P`$ be the restriction of $`\rho `$ to $`(M\text{α}N)P`$. Then $`\rho _P`$ is a $``$-isomorphism onto $`N_2`$. When $`\eta `$ is a n.s.f. weight on $`M\text{α}N`$ we have $`\sigma _t^\eta (P)=P`$ for all $`t`$, because $`P`$ is central. So the restriction $`\eta _P`$ of $`\eta `$ to $`(M\text{α}N)P`$ is a n.s.f. weight and $`\sigma _t^{\eta _P}`$ is the restriction of $`\sigma _t^\eta `$ to $`(M\text{α}N)P`$ for all $`t`$.
For every n.s.f. weight $`\mu `$ on $`N`$ we define the n.s.f. weight $`\stackrel{ˇ}{\mu }`$ on $`N_2`$ by $`\stackrel{ˇ}{\mu }=(\stackrel{~}{\mu })_P{}_{}{}^{}\rho _{P}^{1}`$. Here $`\stackrel{~}{\mu }`$ denotes as before the dual weight on $`M\text{α}N`$. For every $`xN`$ we have
$$\sigma _t^{\stackrel{ˇ}{\mu }}(x)=\rho _P\left(\sigma _t^{(\stackrel{~}{\mu })_P}(\rho _P^1(x))\right)=\rho _P\left(\sigma _t^{(\stackrel{~}{\mu })_P}(\alpha (x)P)\right)=\rho _P\left(\sigma _t^{\stackrel{~}{\mu }}(\alpha (x))P\right)=\rho _P\left(\alpha (\sigma _t^\mu (x))P\right)=\sigma _t^\mu (x).$$
When $`\mu `$ and $`\nu `$ are both n.s.f. weights on $`N`$ we have
$$[D\stackrel{ˇ}{\mu }:D\stackrel{ˇ}{\nu }]_t=\rho _P([D(\stackrel{~}{\mu })_P:D(\stackrel{~}{\nu })_P]_t)=\rho _P([D\stackrel{~}{\mu }:D\stackrel{~}{\nu }]_tP)=\rho _P(\alpha ([D\mu :D\nu ]_t)P)=[D\mu :D\nu ]_t$$
for all $`t`$. So, by \[26, 12.7\], there exists a unique n.s.f. operator valued weight $`T_2`$ from $`N_2`$ to $`N`$ such that $`\stackrel{ˇ}{\mu }=\mu {}_{}{}^{}T_{2}^{}`$ for all n.s.f. weights $`\mu `$ on $`N`$. So $`\mu {}_{}{}^{}T_{2}^{}{}_{}{}^{}\rho _{P}^{}=(\stackrel{~}{\mu })_P`$ for all n.s.f. weights $`\mu `$ on $`N`$.
When $`\nu `$ is a n.s.f. weight on either $`N^\alpha `$ or $`N`$ we denote again with $`\nu ^{}`$ the n.s.f. weight on either $`J_\theta N^\alpha J_\theta =N_2^{}`$ or $`J_\theta NJ_\theta =N^{}`$ given by $`\nu ^{}(x)=\nu (J_\theta xJ_\theta )`$ for all positive $`x`$ in either $`N^\alpha `$ or $`N`$. By \[26, 12.11\] there exists a unique n.s.f. operator valued weight $`T_1`$ from $`N`$ to $`N^\alpha `$ such that
$$\frac{d(\mu {}_{}{}^{}T_{2}^{})}{d\nu ^{}}=\frac{d\mu }{d\left((\nu {}_{}{}^{}T_{1}^{})^{}\right)}$$
(5.1)
for all n.s.f. weights $`\mu `$ on $`N`$ and $`\nu `$ on $`N^\alpha `$.
Choose now a n.s.f. weight $`\theta _0`$ on $`N^\alpha `$. Put $`\theta _1=\theta _0{}_{}{}^{}T_{1}^{}`$ and $`\theta _2=\theta _1{}_{}{}^{}T_{2}^{}`$. When we change the weight $`\theta `$ which was chosen on $`N`$ in the beginning of the story, the tower $`N^\alpha NN_2`$ will be transformed into a unitarily equivalent tower. The unitary implementing this transformation is the unique unitary intertwining the two standard representations of $`N`$. This unitary also intertwines the two implementations of $`\alpha `$ by proposition 4.1. Hence also $`\rho `$ can be transformed. So we may suppose that $`\theta =\theta _1`$.
From equation 5.1 follows that
$$\frac{d\theta _2}{d\theta _0^{}}=\frac{d\theta _1}{d\theta _1^{}}.$$
So we also have $`{\displaystyle \frac{d\theta _0^{}}{d\theta _2}}={\displaystyle \frac{d\theta _1^{}}{d\theta _1}}`$. But $`{\displaystyle \frac{d\theta _1^{}}{d\theta _1}}=_\theta ^1`$ because $`K`$ is the GNS-space of $`\theta =\theta _1`$. To compute $`{\displaystyle \frac{d\theta _0^{}}{d\theta _2}}`$ we need a GNS-construction for the weight $`\theta _2`$. But $`\theta _2{}_{}{}^{}\rho _{P}^{}=\theta {}_{}{}^{}T_{2}^{}{}_{}{}^{}\rho _{P}^{}=(\stackrel{~}{\theta })_P`$. So we put $`L=P(HK)`$ and as before we denote with $`(HK,\iota ,\stackrel{~}{\mathrm{\Lambda }})`$ the GNS-construction of $`\stackrel{~}{\theta }`$. For every $`x𝒩_{\theta _2}`$ we define $`\mathrm{\Lambda }_{\theta _2}(x)=\stackrel{~}{\mathrm{\Lambda }}(\rho _P^1(x))`$. Then $`\mathrm{\Lambda }_{\theta _2}(x)L`$ and it is easy to check that $`(L,\rho _P^1,\mathrm{\Lambda }_{\theta _2})`$ is a GNS-construction for $`\theta _2`$. Also observe that for all $`a𝒩_{\widehat{\phi }}`$ and $`x𝒩_\theta `$ we have $`\rho (a1)x𝒩_{\theta _2}`$ and
$$\mathrm{\Lambda }_{\theta _2}(\rho (a1)x)=P(\widehat{\mathrm{\Lambda }}(a)\mathrm{\Lambda }_\theta (x)).$$
Now choose $`z𝒯_\theta `$. Then
$$+\mathrm{}>\theta (z^{}z)=\omega _{\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))},_\theta ^1=\omega _{\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))},\frac{d\theta _0^{}}{d\theta _2}=\theta _0^{},\mathrm{\Theta }^{\theta _2}\left(\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))\right).$$
(5.2)
Choose now a family $`(\xi _i)_{iI}`$ of vectors in $`K`$ such that
$$\theta _0^{}(z)=\underset{iI}{}\omega _{\xi _i}(z)$$
for all $`z(J_\theta N^\alpha J_\theta )^+`$. Fix $`iI`$. Then
$$\omega _{\xi _i},\mathrm{\Theta }^{\theta _2}\left(\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))\right)<+\mathrm{}$$
and so $`\xi _i𝒟\left(R^{\theta _2}\left(\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))\right)^{}\right)`$. Further
$$\omega _{\xi _i},\mathrm{\Theta }^{\theta _2}\left(\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))\right)=R^{\theta _2}\left(\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))\right)^{}\xi _i^2.$$
(5.3)
We will compute the final expression. For this we choose $`\omega `$ and $`x𝒩_\theta `$. Recall that the subset $`M_{}`$ was introduced in the introduction. Observe that
$$R^{\theta _2}\left(\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))\right)^{}\xi _iL.$$
So we have
$`R^{\theta _2}\left(\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))\right)^{}\xi _i,\widehat{\mathrm{\Lambda }}\left((\omega \iota )(W)\right)\mathrm{\Lambda }_\theta (x)`$ $`=R^{\theta _2}\left(\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))\right)^{}\xi _i,P\left(\widehat{\mathrm{\Lambda }}\left((\omega \iota )(W)\right)\mathrm{\Lambda }_\theta (x)\right)`$
$`=R^{\theta _2}\left(\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))\right)^{}\xi _i,\mathrm{\Lambda }_{\theta _2}\left(\rho \left((\omega \iota )(W)1\right)x\right)`$
$`=\xi _i,R^{\theta _2}\left(\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))\right)\mathrm{\Lambda }_{\theta _2}\left((\omega \iota )(U^{})x\right)`$
$`=\xi _i,(\omega \iota )(U^{})x\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))`$
$`=\xi _i,(\omega \iota )(U^{})J_\theta z^{}J_\theta \mathrm{\Lambda }_\theta (x)`$
$`=\overline{\omega }\left((\iota \omega _{\xi _i,\mathrm{\Lambda }_\theta (x)})((1J_\theta zJ_\theta )U)\right).`$
By continuity we get that
$$R^{\theta _2}\left(\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))\right)^{}\xi _i,\widehat{\mathrm{\Lambda }}\left((\omega \iota )(W)\right)\eta =\overline{\omega }\left((\iota \omega _{\xi _i,\eta })((1J_\theta zJ_\theta )U)\right)$$
for all $`\omega `$, $`\eta K`$ and $`z𝒯_\theta `$. By proposition 7.3 of the appendix, it follows that
$$(\iota \omega _{\xi _i,\eta })\left((1J_\theta zJ_\theta )U\right)𝒩_\phi $$
and
$$\mathrm{\Lambda }\left((\iota \omega _{\xi _i,\eta })((1J_\theta zJ_\theta )U)\right)=(1\theta _\eta ^{})R^{\theta _2}\left(\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))\right)^{}\xi _i$$
for all $`\eta K`$ and $`z𝒯_\theta `$. Fix an orthonormal basis $`(e_j)_{jJ}`$ for $`K`$. Then we may conclude that
$`R^{\theta _2}\left(\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))\right)^{}\xi _i^2`$ $`={\displaystyle \underset{jJ}{}}(1\theta _{e_j}^{})R^{\theta _2}\left(\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))\right)^{}\xi _i^2`$
$`={\displaystyle \underset{jJ}{}}\phi \left((\iota \omega _{\xi _i,e_j})((1J_\theta zJ_\theta )U)^{}(\iota \omega _{\xi _i,e_j})((1J_\theta zJ_\theta )U)\right)`$
$`=\phi \left((\iota \omega _{\xi _i})(U^{}(1J_\theta z^{}zJ_\theta )U)\right)`$
$`=\phi (\widehat{J}(\iota \omega _{J_\theta \xi _i})\alpha (z^{}z)\widehat{J})`$
$`=\psi \left((\iota \omega _{J_\theta \xi _i})\alpha (z^{}z)\right)`$
$`=\omega _{J_\theta \xi _i},(\psi \iota )\alpha (z^{}z).`$
Combining this with equation 5.3 we get that
$$\omega _{\xi _i},\mathrm{\Theta }^{\theta _2}\left(\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))\right)=\omega _{J_\theta \xi _i},(\psi \iota )\alpha (z^{}z)$$
for all $`z𝒯_\theta `$ and $`iI`$. Summing over $`i`$ we get
$$\theta _0^{},\mathrm{\Theta }^{\theta _2}\left(\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))\right)=\theta _0,(\psi \iota )\alpha (z^{}z)$$
for all $`z𝒯_\theta `$. Using equation 5.2 we get that
$$\theta (z^{}z)=\theta _0,(\psi \iota )\alpha (z^{}z)$$
for all $`z𝒯_\theta `$. Hence the normal faithful weight $`\theta _0{}_{}{}^{}(\psi \iota )\alpha `$ is semifinite. From \[26, 11.7\] it follows that $`(\psi \iota )\alpha `$ is semifinite. So $`\alpha `$ is integrable.
Proof of the second implication. The second implication can be proved along the same lines as in the case of an Abelian group action, see . So let us suppose that $`\alpha `$ is integrable. Choose a n.s.f. weight $`\theta _0`$ on $`N^\alpha `$ and put $`\theta =\theta _0{}_{}{}^{}(\psi \iota )\alpha `$. Then $`\theta `$ is a n.s.f. weight on $`N`$. Represent $`N`$ on the GNS-space $`K`$ of $`\theta `$ such that $`(K,\iota ,\mathrm{\Lambda }_\theta )`$ is a GNS-construction for $`\theta `$. Choose a family of vectors $`(\xi _i)_{iI}`$ in $`K`$ such that
$$\theta _0(x)=\underset{iI}{}\omega _{\xi _i}(x)\text{for all}x(N^\alpha )^+.$$
Define $`L=_{iI}HK`$ and let $`\pi `$ be the $`I`$-fold direct sum of the identical representation $`\iota `$ of $`M\text{α}N`$ on $`HK`$. Recall that for any operator valued weight $`T`$ we define $`𝒩_T`$ as the left ideal of elements $`x`$ for which $`T(x^{}x)`$ is bounded. Also recall that we introduced the canonical GNS-construction $`(H,\iota ,\mathrm{\Gamma })`$ for $`\psi `$ in the introduction. When $`z𝒩_{\psi \iota }`$ we define $`(\mathrm{\Gamma }\iota )(z)B(K,HK)`$ by $`(\mathrm{\Gamma }\iota )(z)\mathrm{\Lambda }_\theta (x)=(\mathrm{\Gamma }\mathrm{\Lambda }_\theta )(z(1x))`$ for all $`x𝒩_\theta `$, where $`\mathrm{\Gamma }\mathrm{\Lambda }_\theta `$ denotes the canonical GNS-map of $`\psi \theta `$. One can check easily that $`(\mathrm{\Gamma }\iota )(z)^{}(\mathrm{\Gamma }\iota )(z)=(\psi \iota )(z^{}z)`$. For this see e.g. \[7, 10.6\]. Put $`T=(\psi \iota )\alpha `$. For all $`x𝒩_T𝒩_\theta `$ we define
$$\text{V}\mathrm{\Lambda }_\theta (x)=\underset{iI}{}(\mathrm{\Gamma }\iota )\alpha (x)\xi _i.$$
Observe that V is well-defined:
$$\underset{iI}{}(\mathrm{\Gamma }\iota )\alpha (x)\xi _i^2=\underset{iI}{}\omega _{\xi _i}\left((\psi \iota )\alpha (x^{}x)\right)=\theta _0,(\psi \iota )\alpha (x^{}x)=\theta (x^{}x)<\mathrm{}.$$
Because $`𝒩_T𝒩_\theta `$ is a $`\sigma `$-strong–norm core for $`\mathrm{\Lambda }_\theta `$ we get that $`\mathrm{\Lambda }_\theta (𝒩_T𝒩_\theta )`$ is dense in $`K`$. So V can be extended uniquely to an isometry from $`K`$ to $`L`$.
We now want to prove that the range of V is invariant under $`\pi (M\text{α}N)`$. So we first choose $`yN`$. Then for every $`x𝒩_T𝒩_\theta `$ we have
$`\pi (\alpha (y))\text{V}\mathrm{\Lambda }_\theta (x)`$ $`={\displaystyle \underset{iI}{}}\alpha (y)(\mathrm{\Gamma }\iota )\alpha (x)\xi _i`$
$`={\displaystyle \underset{iI}{}}(\mathrm{\Gamma }\iota )\alpha (yx)\xi _i=\text{V}\mathrm{\Lambda }_\theta (yx)=\text{V}y\mathrm{\Lambda }_\theta (x).`$
Next we will look at the invariance under $`\pi (\widehat{M})`$. Analogously as in proposition 7.2 of the appendix we have that for every $`x𝒩_\psi `$, $`\xi 𝒟(\delta ^{1/2})`$ and $`\eta H`$, $`(\omega _{\delta ^{1/2}\xi ,\eta }\iota )\mathrm{\Delta }(x)𝒩_\psi `$ and
$$\mathrm{\Gamma }\left((\omega _{\delta ^{1/2}\xi ,\eta }\iota )\mathrm{\Delta }(x)\right)=(\omega _{\xi ,\eta }\iota )(W^{})\mathrm{\Gamma }(x).$$
Then it follows easily that for all $`z𝒩_{\psi \iota }`$ we have
$$x:=(\omega _{\delta ^{1/2}\xi ,\eta }\iota \iota )(\mathrm{\Delta }\iota )(z)𝒩_{\psi \iota }\text{and}(\mathrm{\Gamma }\iota )(x)=\left((\omega _{\xi ,\eta }\iota )(W^{})1\right)(\mathrm{\Gamma }\iota )(z).$$
So for all $`\xi 𝒟(\delta ^{1/2})`$, $`\eta H`$ and $`x𝒩_T𝒩_\theta `$ we have
$`\pi \left((\omega _{\xi ,\eta }\iota )(W^{})1\right)\text{V}\mathrm{\Lambda }_\theta (x)`$ $`={\displaystyle \underset{iI}{}}\left((\omega _{\xi ,\eta }\iota )(W^{})1\right)(\mathrm{\Gamma }\iota )\alpha (x)\xi _i`$
$`={\displaystyle \underset{iI}{}}(\mathrm{\Gamma }\iota )\left((\omega _{\delta ^{1/2}\xi ,\eta }\iota \iota )(\mathrm{\Delta }\iota )\alpha (x)\right)\xi _i`$
$`={\displaystyle \underset{iI}{}}(\mathrm{\Gamma }\iota )\left(\alpha ((\omega _{\delta ^{1/2}\xi ,\eta }\iota )\alpha (x))\right)\xi _i`$
$`=\text{V}\mathrm{\Lambda }_\theta \left((\omega _{\delta ^{1/2}\xi ,\eta }\iota )\alpha (x)\right)=\text{V}(\omega _{\xi ,\eta }\iota )(U)\mathrm{\Lambda }_\theta (x)`$
by propositions 2.4 and 4.3. So the range of V is invariant under $`\pi (M\text{α}N)`$. Then we can define a $``$-homomorphism
$$\rho :M\text{α}NB(K):\rho (z)=\text{V}^{}\pi (z)\text{V}.$$
By the computations above we get that $`\rho (\alpha (x))=x`$ for all $`xN`$ and $`(\iota \rho )(W1)=U^{}`$. Then it follows from proposition 5.2 that $`\rho (M\text{α}N)=N_2`$ and so the theorem is proved. ∎
One can also prove the following more general kind of result, where we do not specify what $`\rho `$ should be.
###### Proposition 5.4.
Let $`\alpha `$ be an action of $`(M,\mathrm{\Delta })`$ on $`N`$. Fix a n.s.f. weight $`\theta `$ on $`N`$ and let $`N`$ act on the GNS-space $`K`$ of $`\theta `$. Consider the inclusions
$$N^\alpha \alpha (N)M\text{α}N\text{and}N^\alpha NN_2=J_\theta (N^\alpha )^{}J_\theta .$$
Then the following statements are equivalent.
* There exists a surjective $``$-homomorphism $`\rho `$ from $`M\text{α}N`$ to $`N_2`$ such that $`\rho `$ is an isomorphism of $`\alpha (N)`$ onto $`N`$ and of $`N^\alpha `$ onto $`N^\alpha `$.
* The action $`\alpha `$ is cocycle-equivalent with an integrable action $`\beta `$ satisfying $`N^\beta =N^\alpha `$.
###### Proof of the first implication..
Suppose the first statement is true. Because $`N`$ is represented on the GNS-space of $`\theta `$, there exists a unitary $`u`$ on $`K`$ such that $`\rho (\alpha (x))=uxu^{}`$ for all $`xN`$ and $`uJ_\theta =J_\theta u`$. Define $`\stackrel{~}{\rho }`$ from $`M\text{α}N`$ to $`B(K)`$ by $`\stackrel{~}{\rho }(z)=u^{}\rho (z)u`$ for all $`zM\text{α}N`$. Then $`\stackrel{~}{\rho }(\alpha (x))=x`$ for all $`xN`$. Further
$$u^{}N^\alpha u=u^{}\rho (N^\alpha )u=u^{}\rho (\alpha (N^\alpha ))u=N^\alpha .$$
Because $`uJ_\theta =J_\theta u`$, we get $`u^{}N_2u=N_2`$, which leads to $`\stackrel{~}{\rho }(M\text{α}N)=N_2`$.
So we may suppose from the beginning that $`\rho (\alpha (x))=x`$ for all $`xN`$. Define the unitary $`XMB(K)`$ by
$$X=(\widehat{J}J_\theta )(\iota \rho )(W1)(\widehat{J}J_\theta ).$$
Put $`\text{V}=XU^{}`$. Clearly $`\text{V}MB(K)`$ and
$$(\widehat{J}J_\theta )\text{V}(\widehat{J}J_\theta )=(\iota \rho )(W1)U.$$
For every $`xN`$ we have
$`(\iota \rho )(W1)U(1x)`$ $`=(\iota \rho )(W1)\alpha (x)U=(\iota \rho )\left((W1)(\iota \alpha )\alpha (x)\right)U`$
$`=(\iota \rho )\left((1\alpha (x))(W1)\right)U=(1x)(\iota \rho )(W1)U.`$
So we get $`(\iota \rho )(W1)UMN^{}`$ and hence $`\text{V}MN`$. In the next computation we denote again with $`L_\theta `$ the $``$-anti-automorphism of $`B(K)`$ given by $`L_\theta (x)=J_\theta x^{}J_\theta `$ for all $`xB(K)`$. Then we have
$`(\mathrm{\Delta }\iota )(\text{V})`$ $`=(\mathrm{\Delta }\iota )(RL_\theta )(\iota \rho )(W^{}1)(\mathrm{\Delta }\iota )(U^{})=(RRL_\theta )(\mathrm{\Delta }\text{op}\rho )(W^{}1)U_{13}^{}U_{23}^{}`$
$`=(RRL_\theta )(\iota \iota \rho )(W_{13}^{}W_{23}^{})U_{13}^{}U_{23}^{}`$
$`=\left((RL_\theta )(\iota \rho )(W^{}1)\right)_{23}U_{23}^{}U_{23}\left((RL_\theta )(\iota \rho )(W^{}1)U^{}\right)_{13}U_{23}^{}`$
$`=\text{V}_{23}(\iota \alpha )(\text{V}).`$
So V is a $`\alpha `$-cocycle. Define the action $`\beta `$ of $`(M,\mathrm{\Delta })`$ on $`N`$ given by $`\beta (x)=\text{V}\alpha (x)\text{V}^{}`$ for all $`xN`$. Then $`\beta (x)=X(1x)X^{}`$ for all $`xN`$. Because the $`\sigma `$-strong closure of $`\{(\omega \iota )(X)\omega M_{}\}`$ equals $`J_\theta \rho (\widehat{M})J_\theta `$ we get that $`N^\beta =J_\theta \rho (M\text{α}N)^{}J_\theta =N^\alpha `$.
To conclude the proof of the first implication we have to show that $`\beta `$ is integrable. For this we will use the previous theorem. From proposition 4.2 it follows that the unitary implementation $`U_\beta `$ of $`\beta `$ is given by
$$U_\beta =\text{V}U(\widehat{J}J_\theta )\text{V}^{}(\widehat{J}J_\theta )=\text{V}(\iota \rho )(W^{}1).$$
From the proof of proposition 4.2 we also get that $`z\text{V}^{}z\text{V}`$ gives an isomorphism from $`M\text{β}N`$ onto $`M\text{α}N`$. So we can define
$$\stackrel{~}{\rho }:M\text{β}NN_2:\stackrel{~}{\rho }(z)=\rho (\text{V}^{}z\text{V}).$$
Then $`\stackrel{~}{\rho }`$ is a surjective $``$-homomorphism onto $`N_2=J_\theta (N^\alpha )^{}J_\theta =J_\theta (N^\beta )^{}J_\theta `$ and clearly $`\stackrel{~}{\rho }(\beta (x))=x`$ for all $`xN`$. Because V is an $`\alpha `$-cocycle we get that
$$(1\text{V}^{})(W^{}1)(1\text{V})=(\iota \alpha )(\text{V})(W^{}1).$$
From this it follows that
$$(\iota \stackrel{~}{\rho })(W^{}1)=(\iota \rho )\left((1\text{V}^{})(W^{}1)(1\text{V})\right)=\text{V}(\iota \rho )(W^{}1)=U_\beta .$$
By the previous theorem we get that $`\beta `$ is integrable.
Proof of the second implication. Conversely suppose that the second statement is valid and take such an action $`\beta `$. Let V be an $`\alpha `$-cocycle such that $`\beta (x)=\text{V}\alpha (x)\text{V}^{}`$ for all $`xN`$. It follows from the proof of proposition 4.2 that
$$\mathrm{\Phi }:M\text{α}NM\text{β}N:\mathrm{\Phi }(z)=\text{V}z\text{V}^{}$$
is an isomorphism and $`\mathrm{\Phi }(\alpha (x))=\beta (x)`$ for all $`xN`$. By the previous theorem we can find a surjective $``$-homomorphism $`\stackrel{~}{\rho }`$ from $`M\text{β}N`$ onto $`J_\theta (N^\beta )^{}J_\theta `$ satisfying $`\stackrel{~}{\rho }(\beta (x))=x`$ for all $`xN`$. Putting $`\rho =\stackrel{~}{\rho }{}_{}{}^{}\mathrm{\Phi }`$ and observing that $`N^\alpha =N^\beta `$ we get the first statement. ∎
We do not know an example of a non-integrable action $`\alpha `$ which is cocycle-equivalent with an integrable action $`\beta `$ satisfying $`N^\alpha =N^\beta `$, but it seems to be natural that this kind of actions will exist. We will now specify a case in which it cannot exist. This should be compared with the example of a finite group acting outerly on a factor as described above.
###### Definition 5.5.
An action $`\alpha `$ of a locally compact quantum group $`(M,\mathrm{\Delta })`$ on $`N`$ is called outer when
$$M\text{α}N\alpha (N)^{}=.$$
###### Corollary 5.6.
Let $`\alpha `$ be an outer action of $`(M,\mathrm{\Delta })`$ on $`N`$. Choose again a n.s.f. weight $`\theta `$ on $`N`$ and represent $`N`$ on the GNS-space of $`\theta `$. Let $`N_2=J_\theta (N^\alpha )^{}J_\theta `$ be the basic construction. Then the inclusions
$$N^\alpha \alpha (N)M\text{α}N\text{and}N^\alpha NN_2$$
are isomorphic if and only if $`\alpha `$ is integrable.
###### Proof..
When $`\alpha `$ is integrable, one can use theorem 5.3 and then observe that the $``$-homomorphism $`\rho `$ is faithful because $`M\text{α}N`$ is a factor.
Next suppose that the inclusions stated above are isomorphic. By proposition 5.4 there exists an integrable action $`\beta `$ which is cocycle equivalent with $`\alpha `$ and satisfies $`N^\beta =N^\alpha `$. Let $`\text{V}MN`$ be an $`\alpha `$-cocycle such that $`\beta (x)=\text{V}\alpha (x)\text{V}^{}`$ for all $`xN`$. Then for all $`xN^\alpha =N^\beta `$ we have
$$1x=\beta (x)=\text{V}\alpha (x)\text{V}^{}=\text{V}(1x)\text{V}^{}.$$
Hence we get $`\text{V}M(N(N^\alpha )^{})`$. From our assumption and the fact that $`\alpha `$ is outer it follows that $`N_2N^{}=`$. But then also
$$(N^\alpha )^{}N=J_\theta (N_2N^{})J_\theta =.$$
So we can take $`uM`$ such that $`\text{V}=u1`$. Because V is an $`\alpha `$-cocycle we get that $`\mathrm{\Delta }(u)=uu`$. By the unicity of right invariant weights on $`(M,\mathrm{\Delta })`$ there exists a number $`\lambda >0`$ such that $`\psi (u^{}au)=\lambda \psi (a)`$ for all $`aM^+`$. Then we get that for all $`xN^+`$ we have $`(\psi \iota )\alpha (x)=\lambda (\psi \iota )\beta (x)`$. Because $`\beta `$ is integrable it follows that $`\alpha `$ is integrable. ∎
There exist outer actions which are not integrable: see 6.3. Combining the previous result with theorem 5.1 we get that all actions coming out of Enock and Nest’s construction are integrable.
Next we turn towards the notion of a regular operator valued weight. Suppose $`\alpha `$ is an integrable action of $`(M,\mathrm{\Delta })`$ on $`N`$ and suppose that the $``$-homomorphism $`\rho `$ given by theorem 5.3 is faithful. This will of course be the case whenever $`M\text{α}N`$ is a factor, but also when $`\alpha `$ is a dual action or a semidual action. The latter follows from proposition 5.12. Then we can prove that the operator valued weight $`(\psi \iota )\alpha `$ from $`N`$ to $`N^\alpha `$ is regular. More precisely, we will do the following. By our assumption the basic construction $`N^\alpha NN_2`$ is isomorphic with $`N^\alpha \alpha (N)M\text{α}N`$ through the isomorphism $`\rho `$. Let us denote with $`T_1`$ the operator valued weight $`(\psi \iota )\alpha `$ from $`N`$ to $`N^\alpha `$. Then we can construct the operator valued weight $`T_2`$ from $`N_2`$ to $`N`$ by modular theory and the basic construction, as described above. Through the isomorphism $`\rho `$ the operator valued weight $`T_2`$ is transformed to an operator valued weight from $`M\text{α}N`$ to $`\alpha (N)`$. In the next proposition we prove that this operator valued weight is equal to the canonical operator valued weight $`T=(\widehat{\phi }\iota \iota )\widehat{\alpha }`$ from $`M\text{α}N`$ to $`\alpha (N)`$.
###### Proposition 5.7.
Let $`\alpha `$ be an integrable action of $`(M,\mathrm{\Delta })`$ on $`N`$. Suppose that the $``$-homomorphism $`\rho `$ constructed in theorem 5.3 is faithful. Denote with $`T_2`$ and $`T`$ the operator valued weights defined above. Then $`\rho {}_{}{}^{}T=T_2{}_{}{}^{}\rho `$.
###### Proof..
For clarity we stress that $`T_1`$ is the operator valued weight $`(\psi \iota )\alpha `$ from $`N`$ to $`N^\alpha `$, that $`T_2`$ is obtained out of $`T_1`$ by modular theory and the basic construction, and it goes from $`N_2`$ to $`N`$. Finally $`T`$ is the canonical operator valued weight $`(\widehat{\phi }\iota \iota )\widehat{\alpha }`$ from $`M\text{α}N`$ to $`\alpha (N)`$, giving the dual weights by the formula $`\stackrel{~}{\theta }=\theta {}_{}{}^{}\alpha _{}^{1}{}_{}{}^{}T`$ for all n.s.f. weights $`\theta `$ on $`N`$.
Choose a n.s.f. weight $`\theta _0`$ on $`N^\alpha `$. Put $`\theta =\theta _0{}_{}{}^{}T_{1}^{}`$ and let $`\theta _2=\stackrel{~}{\theta }{}_{}{}^{}\rho _{}^{1}`$. We will prove that $`\theta _2=\theta {}_{}{}^{}T_{2}^{}`$. As in the proof of theorem 5.3 we may suppose that $`N`$ is represented on the GNS-space of $`\theta `$ such that $`(K,\iota ,\mathrm{\Lambda }_\theta )`$ is a GNS-construction for $`\theta `$. Let $`(HK,\iota ,\stackrel{~}{\mathrm{\Lambda }})`$ be the canonical GNS-construction for $`\stackrel{~}{\theta }`$ and put $`\mathrm{\Lambda }_{\theta _2}=\stackrel{~}{\mathrm{\Lambda }}{}_{}{}^{}\rho _{}^{1}`$. We now make a kind of converse reasoning of the proof of theorem 5.3. Denote again with $`\theta _0^{}`$ the n.s.f. weight on $`J_\theta N^\alpha J_\theta =N_2^{}`$ given by $`\theta _0^{}(x)=\theta _0(J_\theta xJ_\theta )`$ for all positive $`x`$. We claim that for all $`z𝒯_\theta `$
$$\omega _{\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))},\frac{d\theta _0^{}}{d\theta _2}=\theta _0,(\psi \iota )\alpha (z^{}z).$$
(5.4)
So choose $`z𝒯_\theta `$. Take a family of vectors $`(\xi _i)_{iI}`$ in $`K`$ such that
$$\theta _0(x)=\underset{iI}{}xJ_\theta \xi _i,J_\theta \xi _i$$
for all $`x(N^\alpha )^+`$. Because $`\theta _0,(\psi \iota )\alpha (z^{}z)=\theta (z^{}z)<\mathrm{}`$ we have
$$\omega _{J_\theta \xi _i},(\psi \iota )\alpha (z^{}z)<\mathrm{}$$
for all $`iI`$. Fix $`iI`$. Then we conclude from the previous formula that
$$\phi \left((\iota \omega _{\xi _i})(U^{}(1J_\theta z^{}zJ_\theta )U)\right)<\mathrm{}.$$
So, when $`(e_j)_{jJ}`$ is an orthonormal basis for $`K`$ we can define the element $`\eta HK`$ by
$$\eta :=\underset{jJ}{}\mathrm{\Lambda }\left((\iota \omega _{\xi _i,e_j})((1J_\theta zJ_\theta )U)\right)e_j.$$
It is easy to check that for all $`\mu K`$ we have $`(\iota \omega _{\xi _i,\mu })((1J_\theta zJ_\theta )U)𝒩_\phi `$ and
$$\mathrm{\Lambda }\left((\iota \omega _{\xi _i,\mu })((1J_\theta zJ_\theta )U)\right)=(1\theta _\mu ^{})\eta .$$
Using the notation $`M_{}`$ introduced in the introduction, we get for all $`\omega `$ and $`x𝒩_\theta `$ that
$`\xi _i,R^{\theta _2}\left(\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))\right)\mathrm{\Lambda }_{\theta _2}\left((\omega \iota )(U^{})x\right)`$ $`=\xi _i,(\omega \iota )(U^{})J_\theta z^{}J_\theta \mathrm{\Lambda }_\theta (x)`$
$`=\overline{\omega }\left((\iota \omega _{\xi _i,\mathrm{\Lambda }_\theta (x)})((1J_\theta zJ_\theta )U)\right)`$
$`=\eta ,\widehat{\mathrm{\Lambda }}\left((\omega \iota )(W)\right)\mathrm{\Lambda }_\theta (x)`$
$`=\eta ,\mathrm{\Lambda }_{\theta _2}\left((\omega \iota )(U^{})x\right).`$
From this we get that
$$\xi _i,R^{\theta _2}\left(\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))\right)\mathrm{\Lambda }_{\theta _2}(y)=\eta ,\mathrm{\Lambda }_{\theta _2}(y)$$
for all $`y𝒩_{\theta _2}`$. Hence $`\xi _i𝒟\left(R^{\theta _2}\left(\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))\right)^{}\right)`$ and
$$R^{\theta _2}\left(\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))\right)^{}\xi _i^2=\eta ^2=\omega _{J_\theta \xi _i},(\psi \iota )\alpha (z^{}z).$$
This means that
$$\omega _{\xi _i},\mathrm{\Theta }^{\theta _2}\left(\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))\right)=\omega _{J_\theta \xi _i},(\psi \iota )\alpha (z^{}z).$$
Summing over $`i`$ we get our claim stated in equation 5.4. But now
$$\theta _0,(\psi \iota )\alpha (z^{}z)=\theta (z^{}z)=\omega _{\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))},_\theta ^1,$$
and so
$$\omega _{\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))},_\theta ^1=\omega _{\mathrm{\Lambda }_\theta (\sigma _{i/2}^\theta (z))},\frac{d\theta _0^{}}{d\theta _2}$$
for all $`z𝒯_\theta `$. Next we claim that $`{\displaystyle \frac{d\theta _0^{}}{d\theta _2}}`$ and $`_\theta ^1`$ commute strongly. Then we will be able to conclude that $`{\displaystyle \frac{d\theta _0^{}}{d\theta _2}}=_\theta ^1`$. But then $`{\displaystyle \frac{d\theta _2}{d\theta _0^{}}}=_\theta `$, and so we will get
$$\frac{d(\theta {}_{}{}^{}T_{2}^{})}{d\theta _0^{}}=\frac{d\theta }{d\left((\theta _0{}_{}{}^{}T_{1}^{})^{}\right)}=\frac{d\theta }{d\theta ^{}}=_\theta =\frac{d\theta _2}{d\theta _0^{}}.$$
So we may conclude that $`\theta _2=\theta {}_{}{}^{}T_{2}^{}`$. By definition of $`\stackrel{~}{\theta }`$ we have $`\stackrel{~}{\theta }=\theta {}_{}{}^{}\rho {}_{}{}^{}T`$ and then $`\theta {}_{}{}^{}\rho {}_{}{}^{}T=\stackrel{~}{\theta }=\theta _2{}_{}{}^{}\rho =\theta {}_{}{}^{}T_{2}^{}{}_{}{}^{}\rho `$. By \[26, 11.13\] we get that $`\rho {}_{}{}^{}T=T_2{}_{}{}^{}\rho `$.
So we only have to prove our claim. Hence we want to prove that $`{\displaystyle \frac{d\theta _0^{}}{d\theta _2}}`$ and $`_\theta ^{it}`$ commute for every $`t`$. For this it is sufficient to prove that $`\mathrm{Ad}_\theta ^{it}`$ leaves both $`J_\theta N^\alpha J_\theta `$ and $`N_2`$ invariant and
$$\theta _0^{}{}_{}{}^{}\mathrm{Ad}_\theta ^{it}=\theta _0^{}\text{and}\theta _2{}_{}{}^{}\mathrm{Ad}_\theta ^{it}=\theta _2$$
for all $`t`$. When $`xN^\alpha `$ we have
$$_\theta ^{it}J_\theta xJ_\theta _\theta ^{it}=J_\theta \sigma _t^\theta (x)J_\theta =J_\theta \sigma _t^{\theta _0}(x)J_\theta J_\theta N^\alpha J_\theta .$$
Then it is immediately clear that $`\theta _0^{}{}_{}{}^{}\mathrm{Ad}_\theta ^{it}=\theta _0^{}`$.
Because $`N_2=(J_\theta N^\alpha J_\theta )^{}`$ we have that $`\mathrm{Ad}_\theta ^{it}`$ leaves $`N_2`$ invariant. Recall that we denoted with $`(\stackrel{~}{\sigma }_t)`$ the modular group of $`\stackrel{~}{\theta }`$ on $`M\text{α}N`$. Then we have, for all $`xN`$
$$_\theta ^{it}\rho (\alpha (x))_\theta ^{it}=_\theta ^{it}x_\theta ^{it}=\sigma _t^\theta (x)=\rho \left(\alpha (\sigma _t^\theta (x))\right)=\rho \left(\stackrel{~}{\sigma }_t(\alpha (x))\right).$$
(5.5)
Finally, for all $`\omega B(H)_{}`$ we have by proposition 4.3 and 2.4 that
$`_\theta ^{it}\rho \left((\omega \iota )(W)1\right)_\theta ^{it}`$ $`=_\theta ^{it}(\omega \iota )(U^{})_\theta ^{it}=(Q^{it}\omega Q^{it}\iota )(U^{})`$
$`=\rho \left((Q^{it}\omega Q^{it}\iota )(W)1\right)`$
$`=\rho \left((Q^{it}_\theta ^{it})\left((\omega \iota )(W)1\right)(Q^{it}_\theta ^{it})\right)`$
where we used that $`W(Q^{it}Q^{it})=(Q^{it}Q^{it})W`$. From the proof of proposition 4.3 it follows that $`\stackrel{~}{}^{it}=Q^{it}_\theta ^{it}`$ and so we see that
$$_\theta ^{it}\rho (a1)_\theta ^{it}=\rho (\stackrel{~}{\sigma }_t(a1))$$
for all $`a\widehat{M}`$ and $`tR`$. Combining this with equation 5.5 we get that $`_\theta ^{it}\rho (z)_\theta ^{it}=\rho (\stackrel{~}{\sigma }_t(z))`$ for all $`zM\text{α}N`$ and $`t`$. Then we get immediately that $`\theta _2{}_{}{}^{}\mathrm{Ad}_\theta ^{it}=\theta _2`$ for all $`t`$.
This proves our claim and ends the proof of the proposition. ∎
###### Corollary 5.8.
Under the same assumptions as in proposition 5.7, the operator valued weight $`(\psi \iota )\alpha `$ from $`N`$ to $`N^\alpha `$ is regular.
###### Proof..
Using the notations introduced above we will identify the inclusions $`N^\alpha NN_2`$ and $`N^\alpha \alpha (N)M\text{α}N`$. Then we get that $`T_2=(\widehat{\phi }\iota \iota )\widehat{\alpha }`$. Now it is obvious that $`\widehat{M}M\text{α}N(N^\alpha )^{}`$ and $`𝒩_{\widehat{\phi }}𝒩_{T_2}`$. So the restriction of $`T_2`$ to $`N_2(N^\alpha )^{}`$ is semifinite.
Next observe that $`\alpha (N)=(M\text{α}N)^{\widehat{\alpha }}`$. Applying the first part of the proof to the dual action $`\widehat{\alpha }`$, which is integrable and for which the $``$-homomorphism $`\rho `$ is faithful by proposition 5.12, we get that the restriction of $`T_3`$ to $`N_3N_1^{}`$ is semifinite. ∎
As a final ingredient for the converse of Enock and Nest’s theorem we look at depth 2 inclusions. The assumption of the following proposition may seem strange, but one can immediately look at the corollary for a more clear result.
###### Proposition 5.9.
Let $`\alpha `$ be an action of $`(M,\mathrm{\Delta })`$ on $`N`$ such that $`N^\alpha \alpha (N)M\text{α}N`$ is the basic construction. Then the inclusion $`N^\alpha N`$ has depth 2.
###### Proof..
Choose a n.s.f. weight $`\theta `$ on $`N`$ and let $`\stackrel{~}{\theta }`$ be the dual weight on $`M\text{α}N`$. Represent $`N`$ on the GNS-space of $`\theta `$ such that $`(K,\iota ,\mathrm{\Lambda }_\theta )`$ is a GNS-construction for $`\theta `$. Let $`(HK,\iota ,\stackrel{~}{\mathrm{\Lambda }})`$ be the canonical GNS-construction for $`\stackrel{~}{\theta }`$ and denote with $`\stackrel{~}{J}`$ the modular conjugation of $`\stackrel{~}{\theta }`$. Then it follows from definition 3.6 that $`U=\stackrel{~}{J}(\widehat{J}J_\theta )`$ is the unitary implementation of $`\alpha `$. The basic construction from $`\alpha (N)M\text{α}N`$ is then given by
$$\stackrel{~}{J}\alpha (N)^{}\stackrel{~}{J}=\stackrel{~}{J}U(B(H)N^{})U^{}\stackrel{~}{J}=B(H)N.$$
To prove that $`N^\alpha N`$ has depth 2, we have to show that
$$\alpha (N(N^\alpha )^{})(M\text{α}N)(N^\alpha )^{}B(H)(N(N^\alpha )^{})$$
is the basic construction. But it is immediately clear that the restriction of $`\alpha `$ to $`N(N^\alpha )^{}`$ is an action $`\beta `$ of $`(M,\mathrm{\Delta })`$ on $`N(N^\alpha )^{}`$. So by the first part of the proof it is sufficient to prove that
$$M\text{β}(N(N^\alpha )^{})=(M\text{α}N)(N^\alpha )^{}.$$
(5.6)
Now it follows from theorems 2.6 and 2.7 that
$`M\text{β}(N(N^\alpha )^{})`$ $`=\{zB(H)(N(N^\alpha )^{})(\iota \beta )(z)=V_{12}z_{13}V_{12}^{}\}`$
and
$`M\text{α}N`$ $`=\{zB(H)N(\iota \alpha )(z)=V_{12}z_{13}V_{12}^{}\}.`$
From this we can immediately deduce equation 5.6, and that concludes the proof. ∎
Although the following result is an immediate corollary of the previous one, we include it for completeness. The first statement is clear and the next two statements follow from the first, using proposition 5.12 for the last one.
###### Corollary 5.10.
Let $`\alpha `$ be an action of $`(M,\mathrm{\Delta })`$ on $`N`$.
* If $`\alpha `$ is integrable and the $``$-homomorphism in theorem 5.3 is faithful, then the inclusion $`N^\alpha N`$ has depth 2.
* If $`\alpha `$ is integrable and $`M\text{α}N`$ is a factor, then the inclusion $`N^\alpha N`$ has depth 2.
* The inclusion $`\alpha (N)M\text{α}N`$ has depth 2.
We now prove the announced result giving a converse to the theorem of Enock and Nest.
###### Proposition 5.11.
Let $`\alpha `$ be an integrable outer action of $`(M,\mathrm{\Delta })`$ on $`N`$. Then the operator valued weight $`(\psi \iota )\alpha `$ from $`N`$ to $`N^\alpha `$ is regular. Further the inclusion $`N^\alpha N`$ is irreducible and has depth 2.
###### Proof..
Because $`M\text{α}N`$ is a factor the $``$-homomorphism $`\rho `$ from theorem 5.3 is faithful. Then we apply corollary 5.8 to obtain the regularity of $`(\psi \iota )\alpha `$ and corollary 5.10 to get that $`N^\alpha N`$ has depth 2. It is clear that $`N^\alpha N`$ is irreducible, because
$$N(N^\alpha )^{}=J_\theta (N_2N^{})J_\theta =.$$
As a complement to theorem 5.3 we prove the following easy result. The terminology is taken from .
###### Proposition 5.12.
Let $`\alpha `$ be an action of $`(M,\mathrm{\Delta })`$ on $`N`$. Then we call $`\alpha `$ semidual when there exists a unitary $`vB(H)N`$ satisfying $`(\iota \alpha )(v)=v_{13}V_{12}^{}`$.
* Every dual action is semidual.
* Every semidual action is integrable and the $``$-homomorphism $`\rho `$ from theorem 5.3 is faithful.
###### Proof..
Let us first prove the first statement. Denote with $`\widehat{\alpha }`$ the dual action, which is an action of $`(\widehat{M},\widehat{\mathrm{\Delta }}\text{op})`$ on $`M\text{α}N`$. Because $`\widehat{\phi }`$ is the right Haar weight of $`(\widehat{M},\widehat{\mathrm{\Delta }}\text{op})`$, the role of $`V`$ is played by $`\mathrm{\Sigma }\widehat{W}^{}\mathrm{\Sigma }=W`$. So we have to find a unitary $`vB(H)(M\text{α}N)`$ satisfying $`(\iota \widehat{\alpha })(v)=v_{13}W_{12}^{}`$. Then it is clear that we can take $`v=W^{}1`$ and so $`\widehat{\alpha }`$ is semidual.
To prove the second part suppose that $`vB(H)N`$ is unitary and $`(\iota \alpha )(v)=v_{13}V_{12}^{}`$. Define the isomorphism $`\mathrm{\Psi }:B(H)NB(H)N`$ by $`\mathrm{\Psi }(z)=vzv^{}`$. Using the notation of theorem 2.6 we get that $`\mu (\mathrm{\Psi }(z))=(\iota \mathrm{\Psi })\gamma (z)`$ for all $`zB(H)N`$. So the action $`\mu `$ of $`(M,\mathrm{\Delta })`$ on $`B(H)N`$ is isomorphic with the action $`\gamma `$, which is integrable because it is isomorphic with the bidual action $`\widehat{\alpha }\widehat{\text{}}`$. Hence $`\mu `$ is integrable, and so $`\alpha `$ is integrable.
Fix now a n.s.f. weight $`\theta `$ on $`N`$ and represent $`N`$ on the GNS-space of $`\theta `$ such that $`(K,\iota ,\mathrm{\Lambda }_\theta )`$ is a GNS-construction. Let $`N_2=J_\theta (N^\alpha )^{}J_\theta `$ be the basic construction from $`N^\alpha N`$ and let $`\rho :M\text{α}NN_2`$ be the $``$-homomorphism from theorem 5.3. Then define $`w=(\widehat{J}J_\theta )v(\widehat{J}J_\theta )`$ and define
$$\eta :N_2B(HK):\eta (z)=Uw^{}(1z)wU^{}\text{for all}zN_2.$$
Because $`wB(H)N^{}`$ we have
$$\eta (x)=U(1x)U^{}=\alpha (x)$$
for all $`xN`$. Further we have $`(\iota \alpha )(v)=v_{13}V_{12}^{}`$ and so $`U_{23}v_{13}U_{23}^{}=v_{13}V_{12}^{}`$. Putting $`\widehat{J}\widehat{J}J_\theta `$ around this equation and using that $`V=(\widehat{J}\widehat{J})\mathrm{\Sigma }W^{}\mathrm{\Sigma }(\widehat{J}\widehat{J})`$ (see \[21, 2.15\]), we get
$$U_{23}^{}w_{13}U_{23}=w_{13}(\mathrm{\Sigma }W\mathrm{\Sigma })_{12}.$$
Flipping the first two legs of this equation and rewriting it we get
$$w_{23}^{}U_{13}^{}w_{23}=W_{12}U_{13}^{}.$$
From this it follows that
$$U_{23}w_{23}^{}U_{13}^{}w_{23}U_{23}^{}=U_{23}W_{12}U_{13}^{}U_{23}^{}=U_{23}W_{12}(\mathrm{\Delta }\iota )(U^{})=U_{23}U_{23}^{}W_{12}=W_{12}.$$
Then we get for all $`\omega M_{}`$ that
$$\eta \left((\omega \iota )(U^{})\right)=(\omega \iota \iota )(U_{23}w_{23}^{}U_{13}^{}w_{23}U_{23}^{})=(\omega \iota )(W)1.$$
Hence we may conclude that $`\eta {}_{}{}^{}\rho =\iota `$ and so $`\rho `$ is faithful. ∎
## 6 Minimal actions and outer actions
In definition 5.5 we already defined the notion of an outer action. In the literature one usually encounters the notion of outer action when dealing with discrete group actions and one encounters the notion of minimal action when dealing with compact group actions. In this section we will prove how both notions can be linked in a locally compact quantum group setting. We will also prove a generalization of the main theorem of Yamanouchi, : when working on separable Hilbert spaces, we prove that every integrable outer action with infinite fixed point algebra is a dual action.
The following definition appears in \[15, 4.3\] when dealing with actions of compact Kac algebras.
###### Definition 6.1.
An action $`\alpha `$ of $`(M,\mathrm{\Delta })`$ on $`N`$ is called minimal when
$$N(N^\alpha )^{}=\text{and}\{(\iota \omega )\alpha (x)\omega N_{},xN\}^{\prime \prime }=M.$$
We will prove the following result.
###### Proposition 6.2.
Let $`\alpha `$ be an action of $`(M,\mathrm{\Delta })`$ on $`N`$.
* If $`\alpha `$ is minimal, then $`\alpha `$ is outer.
* If $`\alpha `$ is outer and integrable, then $`\alpha `$ is minimal.
###### Proof..
Let $`\alpha `$ be minimal. Let $`z(M\text{α}N)\alpha (N)^{}`$. Then certainly $`z(B(H)N)(N^\alpha )^{}`$ and hence $`zB(H)`$ by minimality. We now claim that for $`xB(H)`$ we have $`x1M\text{α}N`$ if and only if $`x\widehat{M}`$. Suppose $`x1M\text{α}N`$. It is clear that for every $`zM\text{α}N`$ we have $`(\iota \alpha )(z)=V_{12}z_{13}V_{12}^{}`$. So we get $`(x1)V=V(x1)`$. From this it follows that $`x\widehat{M}`$. So we may conclude that $`z=x1`$, where $`x\widehat{M}`$. Because $`z\alpha (N)^{}`$ we get that $`(x1)\alpha (y)=\alpha (y)(x1)`$ for all $`yN`$. By minimality we get $`xM^{}`$. But then $`xM^{}\widehat{M}=`$ and so $`z`$. Hence $`\alpha `$ is outer.
Let now $`\alpha `$ be outer and integrable. Choose a n.s.f. weight $`\theta `$ on $`N`$ and represent $`N`$ on the GNS-space of $`\theta `$. Let $`J_\theta `$ denote the modular conjugation of $`\theta `$ and let $`N_2=J_\theta (N^\alpha )^{}J_\theta `$ be the basic construction from $`N^\alpha N`$. Let $`\rho `$ be the $``$-homomorphism given in theorem 5.3. Then $`\rho `$ is faithful because $`M\text{α}N`$ is a factor. Because $`\rho `$ is an isomorphism we get $`N_2N^{}=`$ and so
$$N(N^\alpha )^{}=J_\theta (N_2N^{})J_\theta =.$$
Next we claim that $`\left(\alpha (N)N^{}\right)^{\prime \prime }=MB(K)`$. Because, by theorem 2.6, $`B(H)N=\left(M\text{α}NM\right)^{\prime \prime }`$, we get
$`B(H)\left(\alpha (N)N^{}\right)^{\prime \prime }`$ $`=\left((\iota \alpha )(B(H)N)N^{}\right)^{\prime \prime }`$
$`=\left((\iota \alpha )(M\text{α}N)MN^{}\right)^{\prime \prime }`$
$`=V_{12}\left((M\text{α}N)_{13}V^{}(M)VN^{}\right)^{\prime \prime }V_{12}^{}.`$
When $`\stackrel{~}{J}`$ denotes the modular conjugation of the dual weight $`\stackrel{~}{\theta }`$, we already observed in the proof of proposition 5.9 that $`B(H)N=\stackrel{~}{J}\alpha (N)^{}\stackrel{~}{J}`$. Then the outerness of $`\alpha `$ implies that $`B(H)N(M\text{α}N)^{}=`$ and so
$$\left(N^{}M\text{α}N\right)^{\prime \prime }=B(H)B(K).$$
Then we may conclude from the previous computation that
$`B(H)\left(\alpha (N)N^{}\right)^{\prime \prime }`$ $`=V_{12}\left(B(H)B(K)V^{}(M)V\right)^{\prime \prime }V_{12}^{}`$
$`=\left(V(B(H))V^{}B(K)M\right)^{\prime \prime }`$
$`=\left(\mathrm{\Delta }(M)B(K)(\widehat{M}M)\right)^{\prime \prime }`$
$`=B(H)MB(K)`$
where we have used that $`V\widehat{M}^{}M`$, $`(\widehat{M}M)^{\prime \prime }=B(H)`$ and $`\left(\mathrm{\Delta }(M)B(H)\right)^{\prime \prime }=B(H)M`$. Then our claim follows and hence it is clear that
$$\{(\iota \omega )\alpha (x)\omega N_{},xN\}^{\prime \prime }=M.$$
So $`\alpha `$ is minimal. ∎
We will now give an example of an outer action which is not minimal.
###### Counterexample 6.3.
There exists an action $`\alpha `$ of $``$ on a $`II_1`$-factor $`N`$ such that $`\alpha `$ is outer and $`N^\alpha =`$. Then $`\alpha `$ is clearly not minimal, and neither can $`N^\alpha \alpha (N)M\text{α}N`$ be the basic construction.
###### Proof..
Let $`G`$ be the free group with a countably infinite number of generators $`\{a_nn\}`$. It is well known that the free group factor $`N=(G)`$ is a $`II_1`$-factor. Let $`\beta `$ be the automorphism of $`G`$ satisfying $`\beta (a_n)=a_{n+1}`$ for all $`n`$. Let $`\alpha `$ be the automorphism of $`N`$ satisfying $`\alpha (\lambda _g)=\lambda _{\beta (g)}`$ for all $`gG`$. Define the automorphism group $`(\alpha _n)_n`$ in the usual way by $`\alpha _n=\alpha ^n`$ for all $`n`$. It is easy to verify that $`\alpha `$ is a free action and hence $`\alpha `$ is outer (see \[16, def. 1.4.2 and prop. 1.4.4\]). Further it is easy to check that $`N^\alpha =`$. ∎
We conclude this section with a generalization of the main theorem of Yamanouchi . It is remarkable that the proof of our result is much more easy then Yamanouchi’s proof. In the following result is proved for minimal actions of compact Kac algebras, which are automatically integrable because the Haar weight is finite.
###### Proposition 6.4.
Let $`\alpha `$ be an action of $`(M,\mathrm{\Delta })`$ on $`N`$. Suppose that both $`M`$ and $`N`$ are $`\sigma `$-finite von Neumann algebras (i.e. with separable preduals). If the action $`\alpha `$ is minimal and integrable and if $`N^\alpha `$ is infinite, then $`\alpha `$ is a dual action.
###### Proof..
Consider the action $`\beta `$ of $`(M,\mathrm{\Delta })`$ on $`\stackrel{~}{N}=B(H)NM_2()`$ given by
$$\beta \left(\begin{array}{cc}x& y\\ z& r\end{array}\right)=\left(\begin{array}{cc}\mu (x)& \mu (y)\stackrel{~}{V}^{}\\ \stackrel{~}{V}\mu (z)& \gamma (r)\end{array}\right),$$
for $`x,y,z,rB(H)N`$. Here we used the notations of theorem 2.6: $`\mu (x)=(\sigma \iota )(\iota \alpha )(x)`$, $`\gamma (r)=\stackrel{~}{V}\mu (r)\stackrel{~}{V}^{}`$ and $`\stackrel{~}{V}=\mathrm{\Sigma }V^{}\mathrm{\Sigma }1`$. Let us define now
$$𝒥=\{xB(H)N(\iota \alpha )(x)=x_{13}V_{12}^{}\}.$$
Using matrix notation and referring to theorem 2.6 and 2.7, it is then clear that $`x\stackrel{~}{N}^\beta `$ if and only if $`x_{11}B(H)N^\alpha `$, $`x_{22}M\text{α}N`$ and $`x_{12},x_{21}^{}𝒥`$.
Choose a n.s.f. weight $`\theta `$ on $`N`$ and represent $`N`$ on the GNS-space of $`\theta `$ such that $`(K,\iota ,\mathrm{\Lambda }_\theta )`$ is a GNS-construction. Then we fix $`z𝒩_{(\psi \iota )\alpha }`$ and $`\xi H`$ and we claim that the element $`xB(HK)`$ defined by
$$x:=(\mathrm{\Gamma }\iota )\alpha (z)(\theta _\xi ^{}1)$$
belongs to $`𝒥^{}`$. Here we used the notation $`\mathrm{\Gamma }\iota `$ introduced in the proof of theorem 5.3. To prove our claim we observe that for all $`b𝒩_\psi `$, $`y𝒩_\theta `$ and $`\eta H`$
$$(\theta _{\mathrm{\Gamma }(b)}^{}1)x(\eta \mathrm{\Lambda }_\theta (y))=\eta ,\xi (\theta _{\mathrm{\Gamma }(b)}^{}1)(\mathrm{\Gamma }\mathrm{\Lambda }_\theta )(\alpha (z)(1y))=\eta ,\xi (\psi \iota )\left((b^{}1)\alpha (z)\right)\mathrm{\Lambda }_\theta (y).$$
We can conclude that
$$(\omega _{\eta ,\mathrm{\Gamma }(b)}\iota )(x)=\eta ,\xi (\psi \iota )\left((b^{}1)\alpha (z)\right).$$
So $`xB(H)N`$ and for all $`\eta H,b𝒩_\psi `$ and $`\omega N_{}`$ we have
$`(\omega _{\eta ,\mathrm{\Gamma }(b)}\iota \omega )(\iota \alpha )(x)`$ $`=(\iota \omega )\alpha \left(\eta ,\xi (\psi \iota )\left((b^{}1)\alpha (z)\right)\right)`$
$`=\eta ,\xi (\psi \iota )\left((b^{}1)\mathrm{\Delta }((\iota \omega )\alpha (z))\right)`$
$`=\eta ,\xi (\omega _{\mathrm{\Gamma }((\iota \omega )\alpha (z)),\mathrm{\Gamma }(b)}\iota )(V).`$
Next we observe that for all $`y𝒩_\psi `$
$`\eta ,\xi \mathrm{\Gamma }\left((\iota \omega )\alpha (z)\right),\mathrm{\Gamma }(y)`$ $`=\eta ,\xi \omega (\psi \iota )\left((y^{}1)\alpha (z)\right)`$
$`=\omega \left((\omega _{\eta ,\mathrm{\Gamma }(y)}\iota )(x)\right)=(\iota \omega )(x)\eta ,\mathrm{\Gamma }(y).`$
Inserting this in the computation above we get that
$$(\omega _{\eta ,\mathrm{\Gamma }(b)}\iota \omega )(\iota \alpha )(x)=(\omega _{(\iota \omega )(x)\eta ,\mathrm{\Gamma }(b)}\iota )(V)=(\omega _{\eta ,\mathrm{\Gamma }(b)}\iota \omega )(V_{12}x_{13}).$$
Then it follows that $`x𝒥^{}`$.
So we see that $`𝒥\{0\}`$. Because $`\alpha `$ is minimal we also have that $`\alpha `$ is outer by proposition 6.2. In particular $`M\text{α}N`$ is a factor. Also $`N^\alpha `$ is a factor. Because $`𝒥\{0\}`$ we then get immediately that $`\stackrel{~}{N}^\beta `$ is a factor. Because $`M`$ is supposed to be $`\sigma `$-finite, the Hilbert space $`H`$ is separable. So $`\stackrel{~}{N}^\beta `$ is $`\sigma `$-finite. Denoting with $`e_{ij}`$ the matrix units in $`M_2()`$ we see that the projections $`1e_{11}`$ and $`1e_{22}`$ both belong to $`\stackrel{~}{N}^\beta `$. Because $`N^\alpha M\text{α}N`$ both projections are infinite. Hence they are equivalent in the $`\sigma `$-finite factor $`\stackrel{~}{N}^\beta `$. Take $`w\stackrel{~}{N}^\beta `$ such that $`w^{}w=1e_{22}`$ and $`ww^{}=1e_{11}`$. Then there exists a unitary $`v𝒥`$ such that $`w=ve_{12}`$.
Now we can consider the isomorphism
$$\mathrm{\Psi }:B(H)NB(H)N:\mathrm{\Psi }(z)=v^{}zv.$$
It is easy to check that $`(\iota \mathrm{\Psi })\mu (z)=\gamma (\mathrm{\Psi }(z))`$ for all $`zB(H)N`$. So the actions $`\mu `$ and $`\gamma `$ are isomorphic. Because $`\gamma `$ is isomorphic to the bidual action $`\widehat{\alpha }\widehat{\text{}}`$ by theorem 2.6, we get that $`\mu `$ is a dual action. Because $`N^\alpha `$ is properly infinite and because $`H`$ is a separable Hilbert space we get that the action $`\alpha `$ on $`N`$ is isomorphic with the action $`\mu `$ on $`B(H)N`$. So $`\alpha `$ is a dual action. ∎
## 7 Appendix
In this appendix we collect four technical results which do not have anything to do with actions. The first three results are general results on locally compact quantum groups and the last one deals with n.s.f. weights on a von Neumann algebra. We will use freely the notations introduced in the introduction.
###### Proposition 7.1.
Let $`(M,\mathrm{\Delta })`$ be a locally compact quantum group. For every $`\xi H`$ and $`b𝒯_\phi `$ we have
$$\lambda (\omega _{\xi ,\mathrm{\Lambda }(b)})𝒩_{\widehat{\phi }}\text{and}\widehat{\mathrm{\Lambda }}(\lambda (\omega _{\xi ,\mathrm{\Lambda }(b)}))=J\sigma _{i/2}(b)J\xi .$$
Moreover $`\mathrm{span}\{\lambda (\omega _{\xi ,\mathrm{\Lambda }(b)})\xi H,b𝒯_\phi \}`$ is a $`\sigma `$-strong–norm core for $`\widehat{\mathrm{\Lambda }}`$.
###### Proof..
The first statement follows easily from the definition of $`\widehat{\phi }`$. Let $`x𝒩_\phi `$, then
$$\omega _{\xi ,\mathrm{\Lambda }(b)}(x^{})=x^{}\xi ,\mathrm{\Lambda }(b)=J\sigma _{i/2}(b)J\xi ,\mathrm{\Lambda }(x).$$
So we get the first statement. To prove the second one we define
$$=\{a𝒩_\phi \text{there exists}\omega M_{}\text{such that}\omega (x)=\phi (xa)\text{for all}x𝒩_\phi ^{}\}.$$
It is clear that for $`a`$ such a $`\omega M_{}`$ is necessarily unique. We denote it with $`a\phi `$. Then for every $`a`$ we have $`\lambda (a\phi )𝒩_{\widehat{\phi }}`$ and $`\widehat{\mathrm{\Lambda }}(\lambda (a\phi ))=\mathrm{\Lambda }(a)`$. Define $`𝒟_0=\{\lambda (a\phi )a\}`$. We claim that $`𝒟_0`$ is a $`\sigma `$-strong–norm core for $`\widehat{\mathrm{\Lambda }}`$.
Denote with $`𝒟`$ the domain of the $`\sigma `$-strong–norm closure of the restriction of $`\widehat{\mathrm{\Lambda }}`$ to $`𝒟_0`$.
Let $`a`$ and $`t`$. Define $`b=\tau _t(a)\delta ^{it}`$. Then $`b𝒩_\phi `$ and for all $`x𝒩_\phi ^{}`$ we have
$$\phi (xb)=\phi (x\tau _t(a)\delta ^{it})=\nu ^t\phi (\delta ^{it}x\tau _t(a))=\phi (\delta ^{it}\tau _t(x)a)=(a\phi )(\delta ^{it}\tau _t(x))=(\rho _t(a\phi ))(x)$$
where we used the notation of \[19, 8.7\]. So $`b`$ and $`b\phi =\rho _t(a\phi )`$. Hence
$$\widehat{\sigma }_t(\lambda (a\phi ))=\lambda (\rho _t(a\phi ))=\lambda (b\phi )𝒟_0.$$
So we get that $`𝒟_0`$ is invariant under $`\widehat{\sigma }_t`$. Then it is easy to conclude that $`𝒟`$ is invariant under $`\widehat{\sigma }_t`$ for all $`t`$.
Let now $`\omega M_{}`$ and suppose that there exists a $`\mu M_{}`$ such that $`\mu (x)=\omega (S^1(x))`$ for all $`x𝒟(S^1)`$. Let $`a`$. Define $`b=(\mu \iota )\mathrm{\Delta }(a)`$. Then $`b𝒩_\phi `$ and for all $`x𝒩_\phi ^{}`$ we have
$`\phi (xb)`$ $`=\phi \left(x(\mu \iota )\mathrm{\Delta }(a)\right)=\mu \left((\iota \phi )\left((1x)\mathrm{\Delta }(a)\right)\right)=\omega \left((\iota \phi )\left(\mathrm{\Delta }(x)(1a)\right)\right)`$
$`=\phi \left((\omega \iota )\mathrm{\Delta }(x)a\right)=(\omega a\phi )\mathrm{\Delta }(x).`$
So we see that $`b`$ and $`b\phi =(\omega a\phi )\mathrm{\Delta }`$. Then we may conclude that
$$\lambda (\omega )\lambda (a\phi )=\lambda (b\phi )𝒟_0.$$
Because such elements $`\lambda (\omega )`$ form a $`\sigma `$-strong dense subset of $`\widehat{M}`$ it is easy to conclude that $`𝒟`$ is a left ideal in $`\widehat{M}`$.
Because $`𝒟`$ is a $`\sigma `$-strong dense left ideal of $`\widehat{M}`$, invariant under $`\widehat{\sigma }`$ and because $`𝒟𝒩_{\widehat{\phi }}`$, we may conclude that $`𝒟`$ is a $`\sigma `$-strong–norm core for $`\widehat{\mathrm{\Lambda }}`$. But then $`𝒟=𝒩_{\widehat{\phi }}`$ and we have proven our claim.
Then it follows easily that also
$$\mathrm{span}\{\lambda (ab\phi )a,b𝒯_\phi \}$$
is a $`\sigma `$-strong–norm core for $`\widehat{\mathrm{\Lambda }}`$. This last space equals
$$\mathrm{span}\{\lambda (\omega _{\mathrm{\Lambda }(a),\mathrm{\Lambda }(b)})a,b𝒯_\phi \}$$
and so the proposition is proven. ∎
For completeness we also include the following easy result.
###### Proposition 7.2.
Let $`(M,\mathrm{\Delta })`$ be a locally compact quantum group. For every $`a𝒩_\phi `$, $`\xi 𝒟(\delta ^{1/2})`$ and $`\eta H`$ we have $`(\iota \omega _{\xi ,\eta })\mathrm{\Delta }(a)𝒩_\phi `$ and
$$\mathrm{\Lambda }\left((\iota \omega _{\xi ,\eta })\mathrm{\Delta }(a)\right)=(\iota \omega _{\delta ^{1/2}\xi ,\eta })(V)\mathrm{\Lambda }(a).$$
###### Proof..
Let $`(e_n)`$ be the sequence of operators defined in the proof of \[19, 7.6\]. Because $`\mathrm{\Delta }(\delta )=\delta \delta `$ it is clear that
$$\left((\iota \omega _{\xi ,\eta })\mathrm{\Delta }(ae_n)\right)\delta ^{1/2}(\iota \omega _{\delta ^{1/2}\xi ,\eta })\mathrm{\Delta }\left(a(\delta ^{1/2}e_n)\right).$$
Because $`a(\delta ^{1/2}e_n)𝒩_\psi `$ we have $`(\iota \omega _{\delta ^{1/2}\xi ,\eta })\mathrm{\Delta }\left(a(\delta ^{1/2}e_n)\right)𝒩_\psi `$. We know that $`\phi =\psi _{\delta ^1}`$, so that $`(\iota \omega _{\xi ,\eta })\mathrm{\Delta }(ae_n)𝒩_\phi `$ and
$`\mathrm{\Lambda }\left((\iota \omega _{\xi ,\eta })\mathrm{\Delta }(ae_n)\right)`$ $`=\mathrm{\Gamma }\left((\iota \omega _{\delta ^{1/2}\xi ,\eta })\mathrm{\Delta }\left(a(\delta ^{1/2}e_n)\right)\right)`$
$`=(\iota \omega _{\delta ^{1/2}\xi ,\eta })(V)\mathrm{\Gamma }\left(a(\delta ^{1/2}e_n)\right)`$
$`=(\iota \omega _{\delta ^{1/2}\xi ,\eta })(V)\mathrm{\Lambda }(ae_n).`$
Because $`\mathrm{\Lambda }`$ is $`\sigma `$-strong–norm closed, the conclusion follows. ∎
We also need the following technical result.
###### Proposition 7.3.
Let $`(M,\mathrm{\Delta })`$ be a locally compact quantum group and let $`xM`$. Suppose that there exists a vector $`\eta H`$ such that
$$\omega (x^{})=\xi (\omega ),\eta $$
for all $`\omega `$. Then $`x𝒩_\phi `$ and $`\mathrm{\Lambda }(x)=\eta `$.
###### Proof..
Let $`\omega `$ and $`y𝒯_\phi `$. Then we have for all $`a𝒩_\phi `$ that
$$(\omega y)(a^{})=\omega \left((ay^{})^{}\right)=\xi (\omega ),\mathrm{\Lambda }(ay^{})=J\sigma _{i/2}(y^{})J\xi (\omega ),\mathrm{\Lambda }(a).$$
So we get that $`\omega y`$ and $`\xi (\omega y)=J\sigma _{i/2}(y^{})J\xi (\omega )`$. Take now a net $`(e_\alpha )`$ in $`𝒯_\phi `$ such that $`\sigma _z(e_\alpha )1`$ in the $`\sigma `$-strong topology for all $`z`$. Then we have for all $`\omega `$
$$\xi (\omega ),\mathrm{\Lambda }(xe_\alpha )=\omega (e_\alpha ^{}x^{})=(\omega e_\alpha ^{})(x^{})=\xi (\omega e_\alpha ^{}),\eta =\xi (\omega ),J\sigma _{i/2}(e_\alpha )^{}J\eta .$$
Hence $`\mathrm{\Lambda }(xe_\alpha )=J\sigma _{i/2}(e_\alpha )^{}J\eta `$ for all $`\alpha `$. Because $`\mathrm{\Lambda }`$ is $`\sigma `$-strong–norm closed we get $`x𝒩_\phi `$ and $`\mathrm{\Lambda }(x)=\eta `$. ∎
The following result is probably well known, but we could not find it in the literature.
###### Proposition 7.4.
Let $`\theta `$ be a n.s.f. weight on a von Neumann algebra $`N`$ with GNS-construction $`(H,\pi ,\mathrm{\Lambda })`$. Suppose that
* $`𝒟`$ is a weakly dense left ideal in $`N`$ with $`𝒟𝒩_\theta `$.
* $`K`$ is a Hilbert space and $`\mathrm{\Lambda }_0:𝒟K`$ is a linear map such that $`\mathrm{\Lambda }_0(𝒟)`$ is dense in $`K`$.
* $`\pi _0`$ is normal representation of $`N`$ on $`K`$ such that $`\pi _0(x)\mathrm{\Lambda }_0(y)=\mathrm{\Lambda }_0(xy)`$ for all $`xN`$ and $`y𝒟`$.
* V is an isometry from $`K`$ to $`H`$ such that $`\text{V}\mathrm{\Lambda }_0(x)=\mathrm{\Lambda }(x)`$ for all $`x𝒟`$.
* $`\mathrm{\Lambda }_0`$ is $`\sigma `$-strong–norm closed.
Then there exists a unique n.s.f. weight $`\mu `$ on $`N`$ such that $`𝒩_\mu =𝒟`$ and $`(K,\pi _0,\mathrm{\Lambda }_0)`$ is a GNS-construction for $`\mu `$. In particular $`\mu `$ is a restriction of $`\theta `$, which means that for every $`x_\mu ^+`$ we have $`x_\theta ^+`$ and $`\mu (x)=\theta (x)`$.
###### Proof..
Because V is an isometry, $`\mathrm{\Lambda }_0`$ is injective. Define $`𝒰=\mathrm{\Lambda }_0(𝒟𝒟^{})`$. Then $`𝒰`$ is a dense subspace of $`K`$. We make $`𝒰`$ into a $``$-algebra by using $`\mathrm{\Lambda }_0`$ and the $``$-algebra structure on $`𝒟𝒟^{}`$. We claim that $`𝒰`$ is a left Hilbert algebra. The only non-trivial point is to prove that the map $`\mathrm{\Lambda }_0(x)\mathrm{\Lambda }_0(x^{})`$ for $`x𝒟𝒟^{}`$ is closable. But, suppose that $`(x_n)`$ is a sequence in $`𝒟𝒟^{}`$ such that $`\mathrm{\Lambda }_0(x_n)0`$ and $`\mathrm{\Lambda }_0(x_n^{})\xi K`$. Applying V we get $`\mathrm{\Lambda }(x_n)0`$ and $`\mathrm{\Lambda }(x_n^{})\text{V}\xi `$. Because $`\theta `$ is a n.s.f. weight we get that $`\text{V}\xi =0`$ and so $`\xi =0`$. This gives our claim.
It is clear that the von Neumann algebra generated by the left Hilbert algebra $`𝒰`$ is $`\pi _0(N)`$. Because $`\mathrm{\Lambda }_0`$ is injective we have that $`\pi _0`$ is injective. So the n.s.f. weight on $`\pi _0(N)`$ which is canonically associated with $`𝒰`$ can be composed with $`\pi _0`$ to obtain a n.s.f. weight $`\mu `$ on $`N`$. Let $`(K,\pi _0,\mathrm{\Lambda }_1)`$ be the canonically associated GNS-construction.
Then, by definition of $`\mu `$, every $`x𝒟𝒟^{}`$ will belong to $`𝒩_\mu `$ and $`\mathrm{\Lambda }_1(x)=\mathrm{\Lambda }_0(x)`$. Let now $`x𝒟`$. Take a net $`(e_\alpha )`$ in $`𝒟`$ such that $`e_\alpha 1`$ $`\sigma `$-strong. Then we have $`e_\alpha ^{}xx`$ $`\sigma `$-strong, $`e_\alpha ^{}x𝒟𝒟^{}`$ and
$$\mathrm{\Lambda }_1(e_\alpha ^{}x)=\mathrm{\Lambda }_0(e_\alpha ^{}x)=\pi _0(e_\alpha ^{})\mathrm{\Lambda }_0(x)\mathrm{\Lambda }_0(x).$$
Because $`\mathrm{\Lambda }_1`$ is $`\sigma `$-strong–norm closed we get $`x𝒩_\mu `$ and $`\mathrm{\Lambda }_1(x)=\mathrm{\Lambda }_0(x)`$.
Conversely, suppose $`x𝒩_\mu `$. By the main theorem of there exists a net $`(x_\alpha )`$ in $`𝒟𝒟^{}`$ such that $`x_\alpha x`$ for all $`\alpha `$, $`x_\alpha x`$ $`\sigma `$-strong and $`\mathrm{\Lambda }_1(x_\alpha )\mathrm{\Lambda }_1(x)`$ in norm. But $`\mathrm{\Lambda }_1(x_\alpha )=\mathrm{\Lambda }_0(x_\alpha )`$ for every $`\alpha `$. Because $`\mathrm{\Lambda }_0`$ is $`\sigma `$-strong–norm closed we get that $`x𝒟`$. This concludes our proof. ∎
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# Operations, Disturbance, and Simultaneous Measurability
## I Introduction
The probability distribution of the outcome of a measurement is determined by the observable to be measured and the state at the time of the measurement, but the joint probability distribution of the outcomes of two successive measurements on the same object depends on how the first measurement disturbs the object. The disturbance depends not only on the observable and the state but also the apparatus to be used. Thus, the joint probability distribution of successive measurements will be closely related to how the apparatus disturbs the object. It would be an interesting and significant problem to investigate the relation between the disturbance and the joint probability distribution, although there has been no systematic approach to the problem. This paper investigates in particular the relation between the disturbance and the joint probability formula for simultaneous measurements.
In quantum mechanics observables are represented by linear operators, for which the product operation is not necessarily commutative. Two observables are represented by commuting operators if and only if they are simultaneously measurable, and then quantum mechanics predicts the joint probability distribution of the outcomes of their simultaneous measurement \[1, p. 228\]. But, it has not been answered fully what measurement can be considered as a simultaneous measurement of those observables.
The current formulation has two arguments to show how commuting observables can be measured simultaneously. The first argument is based on the fact that any commuting observables $`A`$ and $`B`$ have the third observable $`C`$ for which $`A`$ and $`B`$ are functions of $`C`$ \[1, p. 173\]. In this case, the measurement of $`C`$ gives also the outcomes of the $`A`$-measurement and the $`B`$-measurement simultaneously \[1, p. 228\]. This argument gives one special instance of simultaneous measurement, but it is quite open from this argument how a pair of measuring apparatuses for $`A`$ and $`B`$ makes a simultaneous measurement of $`A`$ and $`B`$.
The second argument assumes the projection postulate proposed by Lüders . The projection postulate determines uniquely the state after the measurement conditional upon the outcome of the measurement, so that for the successive measurement of any pair of discrete observables $`A`$ and $`B`$, the joint probability distribution of their outcomes is determined. According to this probability distribution, if $`A`$ and $`B`$ commute, we have the standard joint probability formula for the simultaneous measurement of $`A`$ and $`B`$. Thus, under the projection postulate, the successive measurements of $`A`$ and $`B`$ are considered effectively as their simultaneous measurement.
If we would restrict the class of measurements to those satisfying the projection postulate, any successive measurements of commuting observables could be considered as a simultaneous measurement. This approach has, however, the following limitations: (i) Some of the most familiar measuring apparatuses such as photon counters do not satisfy the projection postulate. (ii) When the observable has continuous spectrum, no measurements satisfy the repeatability hypothesis , so that the projection postulate cannot be formulated properly for measurements of continuous observables. (iii) The measurement of a function of an observable $`C`$ such as $`A=f(C)`$ using the apparatus measuring $`C`$ does not satisfy the projection postulate in general .
In fact, Einstein, Podolsky, and Rosen (EPR) have derived the joint probability formula for the outcomes of measurements of two observables pertaining to entangled subsystems based on the projection postulate. This EPR-correlation is indeed one instance of the joint probability formula for the simultaneous measurement. The EPR-correlation has been experimentally tested by optical experiments . However, those optical experiments use the photon counting and violate EPR’s original assumption that the measurements satisfy the projection postulate. The recent realizations of quantum teleportation are also optical realizations of the EPR-correlation that violates EPR’s assumption of the projection postulate. Thus, if we should be restricted to measurements satisfying the projection postulate, the scope of measurement theory would exclude most of the recent results in quantum information processing using entanglement .
The modern measurement theory extends the scope from the measurements satisfying the projection postulate to more general measurements described by operations and effects or more generally by operation valued measures introduced from an axiomatic motivation . Thus, it is natural to expect to have a well-defined way to calculate probabilities in any combinations of measuring apparatuses, once we have identified the operation valued measure in the general theory with the given model of measuring apparatus. However, the determination of the operation valued measure so far has relied on the projection postulate or the joint probability formula . Thus, the foundations of the modern approach in the present status involve the same difficulty as establishing the joint probability formula without assuming the projection postulate.
In this paper we shall abandon the projection postulate as a universal quantum rule and consider the following problem: under what condition can a successive measurement of two or more observables be considered as a simultaneous measurement of those observables?
The prospective solution could be stated in the intuitive language that the preceding measurement does not disturb the observable to be measured later. However, in the current quantum mechanics very little has been known about the disturbance caused by general measurement beyond the projection postulate. In order to answer the question in the rigorous language, this paper will attempt to develop a theory of disturbance in general measurements as well as to determine the possible state changes caused by general measurements of observables. The justification of joint probability formula and the EPR-correlation will then follow without assuming the projection postulate.
Section II defines the simultaneous measurement for a pair of measuring apparatus. Sections III and IV discuss simultaneous measurements under the repeatability hypothesis and the projection postulate with indicating their limitations. The following three sections develop theory of general measurement. Section V introduces the nonselective operations and their duals. Section VI discusses the Davies-Lewis postulate for the existence of operation valued measures corresponding to apparatuses and shows that the two justifications of their postulate known so far involve the same difficulty as the one in establishing the joint probability formula without assuming the projection postulate. Section VII gives a new justification of the Davies-Lewis postulate without assuming the projection postulate or the joint probability formula and proves the factoring property of operation valued measures. Section VIII formulates the disturbance in measurement and establishes in the rigorous language the relation between the disturbance and the joint probability formula. Sections IX and X apply the above result to the EPR-correlation and the minimum disturbing measurement. Section XI concludes the paper with some remarks on the uncertainty principle.
## II Statistical formula for simultaneous measurements
### A The Born statistical formula
To formulate the problem precisely, let $`𝐒`$ be a quantum system with the Hilbert space $``$ of state vectors. We shall distinguish measuring apparatuses by their own output variables , denoting by $`𝐀(𝐱)`$ the apparatus measuring the system $`𝐒`$ with the output variable $`𝐱`$, which, we assume, takes values in the real line $`𝐑`$. We shall denote by “$`𝐱(t)\mathrm{\Delta }`$” the probabilistic event that the outcome of the measurement using the apparatus $`𝐀(𝐱)`$ at the time $`t`$ is in a Borel set $`\mathrm{\Delta }`$ in the real line $`𝐑`$. (Throughout this paper, “Borel set” can be replaced by “interval” for simplifying the presentation without any loss of generality.)
Let $`A`$ be an observable of $`𝐒`$. The spectral projection of $`A`$ corresponding to a Borel set $`\mathrm{\Delta }`$ is denoted by $`E^A(\mathrm{\Delta })`$. According to the Born statistical formula, an apparatus $`𝐀(𝐚)`$ with output variable $`𝐚`$ is said to measure an observable $`A`$ at the time $`t`$ if the relation
$$\mathrm{Pr}\{𝐚(t)\mathrm{\Delta }\}=\text{Tr}[E^A(\mathrm{\Delta })\rho (t)]$$
(1)
holds for the state $`\rho (t)`$ of the system $`𝐒`$ at the time $`t`$. The state $`\rho (t)`$ is called the input state to the apparatus $`𝐀(𝐚)`$ and is taken to be an arbitrary density operator.
The relation between the present formulation based on spectral projections due to von Neumann and Dirac’s formulation is as follows. If the observable $`A`$ has the Dirac type spectral decomposition
$$A=\underset{\nu }{}\underset{\mu }{}\mu |\mu ,\nu \mu ,\nu |+\underset{\nu }{}\lambda |\lambda ,\nu \lambda ,\nu |d\lambda ,$$
where $`\mu `$ varies over the discrete eigenvalues, $`\lambda `$ varies over the continuous eigenvalues, and $`\nu `$ is the degeneracy parameter, then we have
$$E^A(\mathrm{\Delta })=\underset{\nu }{}\underset{\mu \mathrm{\Delta }}{}|\mu ,\nu \mu ,\nu |+\underset{\nu }{}_\mathrm{\Delta }|\lambda ,\nu \lambda ,\nu |d\lambda .$$
In this case, we have
$`\text{Tr}[E^A(\mathrm{\Delta })\rho (t)]`$
$`=`$ $`{\displaystyle \underset{\nu }{}}{\displaystyle \underset{\mu \mathrm{\Delta }}{}}\mu ,\nu |\rho (t)|\mu ,\nu +{\displaystyle \underset{\nu }{}}{\displaystyle _\mathrm{\Delta }}\lambda ,\nu |\rho (t)|\lambda ,\nu 𝑑\lambda .`$
### B Simultaneous measurements using one apparatus
Any commuting observables $`A`$ and $`B`$ are simultaneously measurable and the joint probability distribution of the outcomes of their simultaneous measurement is given by
$$\mathrm{Pr}\{𝐚(t)\mathrm{\Delta },𝐛(t)\mathrm{\Delta }^{}\}=\text{Tr}[E^A(\mathrm{\Delta })E^B(\mathrm{\Delta }^{})\rho (t)],$$
(2)
where $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }^{}`$ are arbitrary Borel sets and $`𝐚`$ and $`𝐛`$ denote the output variables of the apparatuses measuring $`A`$ and $`B`$ at the time $`t`$, respectively.
A well-known proof of this formula from (1) runs as follows \[1, p. 228\]. Since $`A`$ and $`B`$ are commutable, there exist an observable $`C`$ and real-valued functions $`f`$ and $`g`$ such that $`A=f(C)`$ and $`B=g(C)`$ \[1, p. 173\]. Their spectral projections satisfy the relations
$`E^A(\mathrm{\Delta })`$ $`=`$ $`E^C(f^1(\mathrm{\Delta })),`$
$`E^B(\mathrm{\Delta }^{})`$ $`=`$ $`E^C(g^1(\mathrm{\Delta }^{})).`$
For the outcome $`c`$ of the $`C`$-measurement, one defines the outcome of the $`A`$-measurement to be $`f(c)`$ and the outcome of the $`B`$-measurement to be $`g(c)`$. Let $`𝐚`$, $`𝐛`$, and $`𝐜`$ be the output variables of the measurements of $`A`$, $`B`$, and $`C`$, respectively. Then, we have
$`\mathrm{Pr}\{𝐚(t)\mathrm{\Delta },𝐛(t)\mathrm{\Delta }^{}\}`$
$`=`$ $`\mathrm{Pr}\{𝐜(t)f^1(\mathrm{\Delta }),𝐜(t)g^1(\mathrm{\Delta }^{})\}`$
$`=`$ $`\mathrm{Pr}\{𝐜(t)f^1(\mathrm{\Delta })g^1(\mathrm{\Delta }^{})\}`$
$`=`$ $`\text{Tr}[E^C(f^1(\mathrm{\Delta })g^1(\mathrm{\Delta }^{}))\rho (t)]`$
$`=`$ $`\text{Tr}[E^C(f^1(\mathrm{\Delta }))E^C(g^1(\mathrm{\Delta }^{}))\rho (t)]`$
$`=`$ $`\text{Tr}[E^A(\mathrm{\Delta })E^B(\mathrm{\Delta }^{})\rho (t)].`$
Thus, their outcomes satisfy (2) so that the measurement of $`C`$ at the time $`t`$ gives a simultaneous measurement of $`A`$ and $`B`$.
### C Simultaneous measurements using two apparatuses
The above proof gives one special instance of simultaneous measurement which uses one measuring apparatus with two output variables, but it is rather open when a pair of measuring apparatuses for $`A`$ and $`B`$ makes a simultaneous measurement of $`A`$ and $`B`$.
In order to formulate this problem precisely, suppose that the observer measures $`A`$ at the time $`t`$ using the apparatus $`𝐀(𝐚)`$. Let $`t+\mathrm{\Delta }t`$ be the time just after the $`𝐀(𝐚)`$-measurement. This means precisely that $`t+\mathrm{\Delta }t`$ is the instant of the time just after the interaction is turned off between $`𝐀(𝐚)`$ and $`𝐒`$ and that after the time $`t+\mathrm{\Delta }t`$ the object $`𝐒`$ is free from the apparatus $`𝐀(𝐚)`$. (Note that the last condition precludes the recoupling of the system with the apparatus.)
Let $`𝐀(𝐛)`$ be another apparatus measuring an observable $`B`$ of $`𝐒`$ with output variable $`𝐛`$. If the measurement using $`𝐀(𝐛)`$ is turned on at the time $`t+\mathrm{\Delta }t`$, the two measurements are called the successive measurement using $`𝐀(𝐚)`$ and $`𝐀(𝐛)`$ (in this order). Then, the successive measurement using $`𝐀(𝐚)`$ and $`𝐀(𝐛)`$ is defined to be a simultaneous measurement of $`A`$ and $`B`$ if and only if the joint probability distribution of their output variables $`𝐚`$ and $`𝐛`$ satisfies the standard joint probability formula
$$\mathrm{Pr}\{𝐚(t)\mathrm{\Delta },𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}\}=\text{Tr}[E^A(\mathrm{\Delta })E^B(\mathrm{\Delta }^{})\rho (t)].$$
(3)
It should be noted that the validity of the above relation depends on how the apparatus $`𝐀(𝐚)`$ disturbs $`𝐒`$ during the first measurement but does not depend on the property of the apparatus $`𝐀(𝐛)`$ as long as $`𝐀(𝐛)`$ measures the observable $`B`$ in any input state. Therefore, the problem to be considered is to find the necessary and sufficient condition on the apparatus $`𝐀(𝐚)`$ measuring $`A`$ in order for the successive measurement using $`𝐀(𝐚)`$ and an arbitrary apparatus $`𝐀(𝐛)`$ measuring $`B`$ to satisfy (3).
In the conventional approach, von Neumann \[1, p. 224\] proved that if two observable are simultaneously measurable, then they are represented by commuting operators under the repeatability hypothesis. (The repeatability hypothesis will be discussed in detail in the next section.) The following theorem, though an easy consequence of the definition, shows that the simultaneous measurability extended by the above definition is still consistent with the old one.
Theorem 1. If the successive measurement of $`A`$ and $`B`$ using apparatuses $`𝐀(𝐚)`$ and $`𝐀(𝐛)`$, respectively, is a simultaneous measurement of observables $`A`$ and $`B`$, then $`A`$ and $`B`$ commute.
Proof. Suppose that (3) holds. By the positivity of probability, the both sides are nonnegative. Since $`\rho (t)`$ is arbitrary, the product $`E^A(\mathrm{\Delta })E^B(\mathrm{\Delta }^{})`$ is a positive self-adjoint operator so that $`E^A(\mathrm{\Delta })`$ and $`E^B(\mathrm{\Delta }^{})`$ commute. Since $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }^{}`$ are arbitrary, it follows that $`A`$ and $`B`$ commute. $`\mathrm{}`$
In the following two sections we shall re-examine the conventional approach from the operational point of view before starting with the general considerations.
## III Simultaneous measurements under the repeatability hypothesis
### A Von Neumann’s formulation
The conventional approach to measurement theory supposes that the measurement leaves the measured system in the eigenstate corresponding to the outcome of the measurement . This assumption is equivalent to the repeatability hypothesis formulated by von Neumann \[1, p. 335\] as follows.
(M) If a physical quantity is measured twice in succession in a system, then we get the same value each time.
Even though the repeatability hypothesis has been posed as a universal law by von Neumann \[1, p. 213\] based on the experiment due to Compton and Simmons, in the modern approach it characterizes merely a class of measuring apparatuses. Thus, in what follows, by saying that the apparatus $`𝐀(𝐚)`$ satisfies the repeatability hypothesis it is meant precisely that the repeatability hypothesis holds for the repeated measurement of $`A`$ using the apparatus $`𝐀(𝐚)`$ for the first $`A`$-measurement.
### B Repeatability hypothesis and joint probability
Suppose that the system $`𝐒`$ is measured at the time $`t`$ by an apparatus $`𝐀(𝐚)`$. Let $`t+\mathrm{\Delta }t`$ be the time just after the measurement. For any Borel set $`\mathrm{\Delta }`$, let $`\rho (t+\mathrm{\Delta }t|𝐚(t)\mathrm{\Delta })`$ be the state at $`t+\mathrm{\Delta }t`$ of $`𝐒`$ conditional upon $`𝐚(t)\mathrm{\Delta }`$. Thus, if the system $`𝐒`$ is sampled randomly from the subensemble of the similar systems that yield the outcome of the $`𝐀(𝐚)`$-measurement in the Borel set $`\mathrm{\Delta }`$, then $`𝐒`$ is in the state $`\rho (t+\mathrm{\Delta }t|𝐚(t)\mathrm{\Delta })`$ at the time $`t+\mathrm{\Delta }t`$. When $`\mathrm{Pr}\{𝐚(t)\mathrm{\Delta }\}=0`$, the state $`\rho (t+\mathrm{\Delta }t|𝐚(t)\mathrm{\Delta })`$ is indefinite, and let $`\rho (t+\mathrm{\Delta }t|𝐚(t)\mathrm{\Delta })`$ be an arbitrarily chosen density operator for mathematical convenience.
Suppose that the apparatus $`𝐀(𝐚)`$ measures a discrete observable $`A`$ with eigenvalues $`a_1,a_2,\mathrm{}`$ and that the $`𝐀(𝐚)`$-measurement at the time $`t`$ is followed immediately by an $`𝐀(𝐛)`$-measurement measuring $`B`$. The conditional probability of $`𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}`$ conditional upon the outcome $`𝐚(t)=a_n`$ is the probability of obtaining the outcome $`𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}`$ in the state $`\rho (t+\mathrm{\Delta }|𝐚(t)=a_n)`$, so that we have
$`\mathrm{Pr}\{𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}|𝐚(t)=a_n\}`$ (4)
$`=`$ $`\text{Tr}[E^B(\mathrm{\Delta }^{})\rho (t+\mathrm{\Delta }t|𝐚(t)=a_n)].`$ (5)
The joint probability distribution of the outcomes of the successive measurement using $`𝐀(𝐚)`$ and $`𝐀(𝐛)`$ is given by the well-known relation
$`\mathrm{Pr}\{𝐚(t)\mathrm{\Delta },𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}\}`$ (6)
$`=`$ $`{\displaystyle \underset{a_n\mathrm{\Delta }}{}}\mathrm{Pr}\{𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}|𝐚(t)=a_n\}\mathrm{Pr}\{𝐚(t)=a_n\}.`$ (7)
Now, suppose that $`A`$ is nondegenerate and that the apparatus $`𝐀(𝐚)`$ satisfies the repeatability hypothesis. Then, the state of the system just after the $`𝐀(𝐚)`$-measurement conditional upon the outcome $`𝐚(t)=a_n`$ is determined uniquely as the normalized eigenstate
$$\rho (t+\mathrm{\Delta }t|𝐚(t)=a_n)=|\varphi _n\varphi _n|$$
(8)
corresponding to the eigenstate $`a_n`$, provided $`\mathrm{Pr}\{𝐚(t)=a_n\}>0`$ \[1, pp. 215–217\].
From (4) and (8), we have
$`{\displaystyle \underset{a_n\mathrm{\Delta }}{}}\mathrm{Pr}\{𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}|𝐚(t)=a_n\}\mathrm{Pr}\{𝐚(t)=a_n\}`$
$`=`$ $`{\displaystyle \underset{a_n\mathrm{\Delta }}{}}\varphi _n|E^B(\mathrm{\Delta }^{})|\varphi _n\varphi _n|\rho (t)|\varphi _n`$
$`=`$ $`{\displaystyle \underset{a_n\mathrm{\Delta }}{}}\text{Tr}[|\varphi _n\varphi _n|E^B(\mathrm{\Delta }^{})|\varphi _n\varphi _n|\rho (t)].`$
Thus, from (6) we have
$`\mathrm{Pr}\{𝐚(t)\mathrm{\Delta },𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}\}`$ (9)
$`=`$ $`{\displaystyle \underset{a_n\mathrm{\Delta }}{}}\text{Tr}[|\varphi _n\varphi _n|E^B(\mathrm{\Delta }^{})|\varphi _n\varphi _n|\rho (t)].`$ (10)
If $`A`$ and $`B`$ commute, we have
$`\mathrm{Pr}\{𝐚(t)\mathrm{\Delta },𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}\}`$
$`=`$ $`{\displaystyle \underset{a_n\mathrm{\Delta }}{}}\text{Tr}[|\varphi _n\varphi _n|E^B(\mathrm{\Delta }^{})\rho (t)]`$
$`=`$ $`\text{Tr}[E^A(\mathrm{\Delta })E^B(\mathrm{\Delta }^{})\rho (t)],`$
and hence we obtain the joint probability formula (3).
By the above, we conclude that if the apparatus $`𝐀(𝐚)`$ measuring a nondegenerate discrete observable $`A`$ satisfies the repeatability hypothesis, the successive measurement using $`𝐀(𝐚)`$ and the apparatus $`𝐀(𝐛)`$ measuring an arbitrary observable $`B`$ commuting with $`A`$ is a simultaneous measurement of $`A`$ and $`B`$.
### C Repeatability hypothesis and measuring processes
Consider a model of measuring process in which an observable $`A=_na_n|\varphi _n\varphi _n|`$ with nondegenerate eigenvalues $`a_n`$ is measured by an apparatus $`𝐀(𝐚)`$ having the nondegenerate probe observable $`B=_na_n|\xi _n\xi _n|`$. The probe observable is generally defined to be the quantum mechanical observable in the apparatus that is to be correlated with the measured observable by the measuring interaction and amplified in the later stage of the apparatus to the directly sensible variable read out eventually by the observer . In the conventional approach, the notion of the “pointer position” is used ambiguously instead, which sometimes means the “probe observable” and sometimes means the “directly sensible variable”. Let $`𝐏`$ be the subsystem of the apparatus $`𝐀(𝐚)`$ that includes the probe observable and actually interacts with the measured system $`𝐒`$. Suppose that $`𝐏`$ is prepared in the state $`|\xi `$ just before measurement. Let $`U`$ be the unitary operator representing the time evolution of the composite system $`𝐒+𝐏`$ during the measuring interaction.
The apparatus $`𝐀(𝐚)`$ satisfies the repeatability hypothesis if and only if $`U`$ satisfies
$$U:|\varphi _n|\xi e^{i\theta _n}|\varphi _n|\xi _n,$$
(11)
where $`e^{i\theta _n}`$ is an arbitrary phase factor.
### D Measurements violating the repeatability hypothesis
A typical model which does not satisfy (11) is the photon counting measurement such that if the measurement is taken place in the number state $`|\varphi _n=|n`$ then the apparatus absorbs all the photons and outputs the amplified classical energy proportional to the number of the absorbed photons. In this case, $`U`$ satisfies
$$U:|n|\xi |0|\xi _n,$$
(12)
and hence $`U`$ does not satisfy (11).
For less idealized models of photon counting measurement, we refer to . For the case of measurements of continuous observables, there have been known also models of exact position measurement that do not satisfy the repeatability hypothesis even approximately . They have been applied to the position monitoring that breaks the standard quantum limit .
### E Significance of the repeatability hypothesis
Given $`A`$, $`B`$, and $`|\xi `$ generally, the apparatus $`𝐀(𝐚)`$ measures the observable $`A`$, or equivalently satisfies (1), if and only if $`U`$ satisfies
$$U:|\varphi _n|\xi |\varphi _n^{}|\xi _n,$$
(13)
where $`\{|\varphi _n^{}\}`$ is an arbitrary family of normalized vectors, not necessarily orthogonal. Thus, if we do not assume the repeatability hypothesis, the measurement is to correlate causally the input state $`|\varphi _n`$ of the object before measurement with the output state $`|\xi _n`$ of the probe after measurement for some orthonormal basis $`\{|\xi _n\}`$. On the other hand, the repeatability hypothesis requires not only that the input state $`|\varphi _n`$ is correlated to the output state $`|\xi _n`$ causally but also that in the composite system after the measurement the input state $`|\varphi _n`$ and the output state $`|\xi _n`$ is entangled to have the complete statistical correlation.
It is said quite often that the measurement of an observable $`A`$ is to change the input state $`|\psi `$ to the eigenstate $`|\varphi _n`$ with the probability $`|\varphi _n|\psi |^2`$. This does not follow from the Born statistical formula (1) but assumes the repeatability hypothesis. Thus, only when the measurement is assumed to satisfy the repeatability hypothesis, we can say that the measurement changes probabilistically the state of the object to one of the eigenstates of the measured observable.
Unless the repeatability hypothesis is assumed, it is, therefore, not a correct description that the measurement is to make a one-to-one correspondence (or to make an entanglement, in the modern language) between the state of the object before the measurement and the state of the probe after the measurement as described by (11), by which the problem of measuring the object is transferred to the problem of measuring the probe .
In the sequel, a measurement is called repeatable if it is carried out by the apparatus satisfying the repeatability hypothesis.
## IV Simultaneous measurements under the projection postulate
### A Von Neumann’s measurements of degenerate observables
For the observables with nondegenerate purely discrete spectrum, the repeatability hypothesis determines the state after measurement uniquely, but when the observable has degenerate eigenvalues, the state after measurement is not determined uniquely but depends on the “actual measuring arrangement” \[1, p. 348\]. Von Neumann \[1, p. 348\] considered the following ways of measurement satisfying the repeatability hypothesis: Let $`\{|\varphi _{n,m}\}`$ be an orthonormal basis and let $`A`$ be an observable represented by
$$A=\underset{n,m}{}a_n|\varphi _{n,m}\varphi _{n,m}|.$$
(14)
Suppose that the observer performs a repeatable measurement of a nondegenerate observable $`A^{}`$ given by
$$A^{}=\underset{n,m}{}a_{n,m}|\varphi _{n,m}\varphi _{n,m}|,$$
(15)
where all $`a_{n,m}`$ are different, and that if the outcome of the $`A^{}`$-measurement is $`a_{n,m}`$ then the outcome of the $`A`$-measurement is taken to be $`a_n`$. Then we have a repeatable measurement of $`A`$.
Suppose that the observable $`A`$ is measured at the time $`t`$ by the above way in a state (vector) $`|\psi `$ using the apparatus $`𝐀(𝐚)`$, then at the time, $`t+\mathrm{\Delta }t`$, just after the measurement the object is left in the state (density operator)
$`\rho (t+\mathrm{\Delta }t|𝐚(t)=a_n)`$ (16)
$`=`$ $`{\displaystyle \frac{1}{_m|c_{n,m}|^2}}{\displaystyle \underset{m}{}}|c_{n,m}|^2|\varphi _{n,m}\varphi _{n,m}|,`$ (17)
where $`c_{n,m}=\varphi _{n,m}|\psi `$. This state depends not only on the observable $`A`$ and the outcome $`a_n`$ but also on the choice of the orthonormal basis $`\{|\varphi _{n,m}\}`$ that satisfies (14). Since there are infinitely many essentially different choices of $`\{|\varphi _{n,m}\}`$, the state change depends on the way of measurement even if the repeatability hypothesis holds.
In this case, the joint probability distribution of the outcomes of the $`𝐀(𝐚)`$-measurement and the immediately following $`B`$-measurement using the apparatus $`𝐀(𝐛)`$ is given by
$`\mathrm{Pr}\{𝐚(t)\mathrm{\Delta },𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}\}`$ (18)
$`=`$ $`{\displaystyle \underset{a_n\mathrm{\Delta }}{}}{\displaystyle \underset{m}{}}\text{Tr}[|\varphi _{n,m}\varphi _{n,m}|E^B(\mathrm{\Delta }^{})|\varphi _{n,m}\varphi _{n,m}|\rho (t)].`$ (19)
From the relation
$$E^A(\mathrm{\Delta })=\underset{a_n\mathrm{\Delta }}{}\underset{m}{}|\varphi _{n,m}\varphi _{n,m}|,$$
the joint probability formula (3) holds for the arbitrary input state $`\rho (t)`$ if and only if $`A^{}`$ and $`B`$ commute. Although $`A`$ and $`B`$ commute, there are many choices of $`\{|\varphi _{n,m}\}`$ such that $`A^{}`$ and $`B`$ do not commute. Thus, even if $`A`$ and $`B`$ commute and $`𝐀(𝐚)`$ satisfies the repeatability hypothesis, the successive measurement using $`𝐀(𝐚)`$ and $`𝐀(𝐛)`$ cannot be a simultaneous measurement of $`A`$ and $`B`$ in general.
### B Lüders’s formulation
The previous argument shows the existence of infinitely many different ways of measuring the same observable which satisfies the repeatability hypothesis but does not satisfy the joint probability formula for the simultaneous measurement. Moreover, Lüders has pointed out that the observable corresponding to the identity operator $`I`$ is considered to be measured without changing the input state, but any one of the above measurement for the identity changes the state unreasonably, and he has suggested that the above measurement for a degenerate observable is always more disturbing than the desirable one.
In order to determine the canonical way of measuring even the degenerate observables, Lüders proposed the following hypothesis.
(P) If an observable $`A`$ is measured in a state $`\rho `$, then at the time just after measurement the object is left in the state
$$\frac{E^A\{a\}\rho E^A\{a\}}{\text{Tr}[E^A\{a\}\rho ]},$$
provided that the object leads to the outcome $`a`$ with $`\text{Tr}[E^A\{a\}\rho ]>0`$.
In particular, if the object is measured in the vector state $`|\psi `$, then the state after measurement is the vector state $`E^A\{a\}|\psi `$ up to normalization. Thus, the eigenstate corresponding to the outcome $`a`$ is uniquely chosen as the projection, and hence the above hypothesis is called the projection postulate.
### C Projection postulate and joint probability
Suppose that a discrete observable $`A`$ of $`𝐒`$ with eigenvalues $`a_1,a_2,\mathrm{}`$ is measured at the time $`t`$ by the apparatus $`𝐀(𝐚)`$ measuring $`A`$ satisfying the projection postulate. Then, the state of the system at the time $`t+\mathrm{\Delta }t`$ just after the $`𝐀(𝐚)`$-measurement conditional upon the outcome $`𝐚(t)=a_n`$ is
$$\rho (t+\mathrm{\Delta }t|𝐚(t)=a_n)=\frac{E^A\{a_n\}\rho (t)E^A\{a_n\}}{\text{Tr}[E^A\{a_n\}\rho (t)]},$$
(21)
provided $`\mathrm{Pr}\{𝐚(t)=a_n\}>0`$.
From (6) and (21), we have
$`\mathrm{Pr}\{𝐚(t)\mathrm{\Delta },𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}\}`$ (22)
$`=`$ $`{\displaystyle \underset{a_n\mathrm{\Delta }}{}}\text{Tr}[E^B(\mathrm{\Delta }^{})E^A\{a_n\}\rho (t)E^A\{a_n\}].`$ (23)
Thus, the joint probability formula (3) holds for the arbitrary input state $`\rho (t)`$ if and only if $`A`$ and $`B`$ commutes. Therefore, we have seen the following theorem .
Theorem 2. The successive measurement of commuting observables $`A`$ and $`B`$ using apparatuses $`𝐀(𝐚)`$ and $`𝐀(𝐛)`$ is a simultaneous measurement of $`A`$ and $`B`$, if $`𝐀(𝐚)`$ satisfies the projection postulate.
We have seen that in order for a successive measurement of $`A`$ and $`B`$ to be a simultaneous measurement, the projection postulate for the apparatus measuring $`A`$ is a sufficient condition. If we restrict our attention to the measuring apparatuses satisfying the projection postulate, any successive measurement of commuting observables is a simultaneous measurement, but we have also seen that there are many kinds of measuring apparatuses which do not satisfy the projection postulate. Now, we shall turn to the problem as to under what condition a successive measurement of two or more observables can be considered as a simultaneous measurement.
## V Nonselective operations of measuring apparatuses
Every measuring process is considered to include an interaction, called the measuring interaction, between the measured system and the measuring apparatus. Let us consider the following description of the measuring interaction arising in the measurement using the apparatus $`𝐀(𝐚)`$. In the following, the probe $`𝐏`$ is a microscopic subsystem of the apparatus $`𝐀(𝐚)`$ that actually interacts with $`𝐒`$. More precisely, we define the probe $`𝐏`$ to be the smallest subsystem of $`𝐀(𝐚)`$ such that the composite system $`𝐒+𝐏`$ is isolated during the measuring interaction. Since we assume naturally that $`𝐒+𝐀(𝐚)`$ is isolated during the measuring interaction, the smallest subsystem exists. The measurement is carried out by the interaction between the system $`𝐒`$ and the probe $`𝐏`$ and by the subsequent measurement on the probe $`𝐏`$. We assume that the probe system is a quantum mechanical system described by the Hilbert space $`𝒦`$ of state vectors. At the time of measurement, $`t`$, the probe $`𝐏`$ is in the fixed state $`\sigma `$, so that the composite system is in the state
$$\rho _{𝐒+𝐏}(t)=\rho (t)\sigma .$$
(24)
The time evolution of the composite system $`𝐒+𝐏`$ during the interaction is described by a unitary operator $`U`$ on $`𝒦`$. Hence, at the time just after the interaction, $`t+\mathrm{\Delta }t`$, the composite system $`𝐒+𝐏`$ is in the state
$$\rho _{𝐒+𝐏}(t+\mathrm{\Delta }t)=U(\rho (t)\sigma )U^{}.$$
(25)
The outcome of this measurement is obtained by the measurement of an observable $`M`$ of $`𝐏`$ at the time $`t+\mathrm{\Delta }t`$. The observable $`M`$ is called the probe observable. Thus, the output probability distribution of the apparatus $`𝐀(𝐚)`$ is
$$\mathrm{Pr}\{𝐚(t)\mathrm{\Delta }\}=\text{Tr}[(IE^M(\mathrm{\Delta }))U(\rho (t)\sigma )U^{}].$$
(26)
We shall call the above description of the measuring process as the indirect measurement model determined by $`(𝒦,\sigma ,U,M)`$. In this model, from (25) the system $`𝐒`$ is in the state
$$\rho (t+\mathrm{\Delta }t)=\mathrm{Tr}_𝒦[U(\rho (t)\sigma )U^{}]$$
(27)
at the time $`t+\mathrm{\Delta }t`$, where $`\mathrm{Tr}_𝒦`$ is the partial trace over $`𝒦`$. The state change $`\rho (t)\rho (t+\mathrm{\Delta }t)`$ is determined independent of the outcome of measurement and called the nonselective state change.
As introduced previously, we have another type of state change $`\rho (t)\rho (t+\mathrm{\Delta }t|𝐚(t)\mathrm{\Delta })`$, where $`\rho (t+\mathrm{\Delta }t|𝐚(t)\mathrm{\Delta })`$ is the state at $`t+\mathrm{\Delta }t`$ of $`𝐒`$ conditional upon $`𝐚(t)\mathrm{\Delta }`$. Since the condition $`𝐚(t)𝐑`$ makes no selection, we have
$$\rho (t+\mathrm{\Delta }t|𝐚(t)𝐑)=\rho (t+\mathrm{\Delta }t).$$
(28)
For $`\mathrm{\Delta }𝐑`$, the state change $`\rho (t)\rho (t+\mathrm{\Delta }t|𝐚(t)\mathrm{\Delta })`$ is called the selective state change.
We define the transformation $`\rho 𝐓\rho `$ for any trace class operator $`\rho `$ on the state space $``$ of $`𝐒`$ by the relation
$$𝐓\rho =\text{Tr}_𝒦[U(\rho \sigma )U^{}].$$
(29)
Then, $`𝐓`$ is a trace preserving completely positive linear transformation on the space $`\tau c()`$ of trace class operators on $``$ . The transformation $`𝐓`$ is determined by the apparatus preparation $`\sigma `$ and by the measuring interaction $`U`$, and is called the nonselective operation of the apparatus $`𝐀(𝐚)`$. The nonselective operation $`𝐓`$ represents the open system dynamics of the system $`𝐒`$ from $`t`$ to $`t+\mathrm{\Delta }t`$ and we have
$$𝐓\rho (t)=\rho (t+\mathrm{\Delta }t).$$
(30)
As the converse of the definition (29), it is well known that for every trace preserving completely positive linear transformation on $`\tau c()`$ there is an indirect measurement model such that $`𝐓`$ is the nonselective operation of that model .
For any bounded linear transformation $`𝐋`$ on $`\tau c()`$, its dual $`𝐋^{}`$ is defined by
$$\text{Tr}[(𝐋^{}X)\rho ]=\text{Tr}[X(𝐋\rho )]$$
for all $`\rho \tau c()`$ and $`X()`$, where $`()`$ stands for the space of bounded operators on $``$. Let $`𝐓^{}`$ be the dual of the nonselective operation $`𝐓`$. Then, $`𝐓^{}`$ is the normal unit preserving completely positive linear transformation on $`()`$ such that
$$\text{Tr}[(𝐓^{}X)\rho ]=\text{Tr}[X(𝐓\rho )]$$
(31)
for all $`\rho \tau c()`$ and $`X()`$ \[11, p. 18\]. We call $`𝐓^{}`$ the dual nonselective operation. Let $`X()`$ and $`\rho \tau c()`$. From (29) and a property of the partial trace, we have
$`\text{Tr}[X𝐓\rho ]`$ $`=`$ $`\text{Tr}[X\text{Tr}_𝒦[U(\rho \sigma )U^{}]`$
$`=`$ $`\text{Tr}[(XI)U(\rho \sigma )U^{}]`$
$`=`$ $`\text{Tr}[U^{}(XI)U(I\sigma )(\rho I)]`$
$`=`$ $`\text{Tr}[\text{Tr}_𝒦[U^{}(XI)U(I\sigma )]\rho ].`$
Hence, from (31) we have
$$\text{Tr}[(𝐓^{}X)\rho ]=\text{Tr}[\text{Tr}_𝒦[U^{}(XI)U(I\sigma )]\rho ].$$
Since $`\rho `$ is arbitrary, we have
$$𝐓^{}X=\text{Tr}_𝒦[U^{}(XI)U(I\sigma )]$$
(32)
for all $`X()`$. This characterizes the dual nonselective operation.
## VI Operation valued measures
### A Davies-Lewis postulates
Davies and Lewis postulated that given an apparatus $`𝐀(𝐚)`$ for $`𝐒`$, there is a mapping $`\mathrm{\Delta }𝐗(\mathrm{\Delta })`$ from the Borel sets to the positive linear transformations on $`\tau c()`$ satisfying the following conditions:
(DL1) For any disjoint sequence of Borel sets $`\mathrm{\Delta }_n`$ and for any $`\rho \tau c()`$,
$$𝐗(\underset{n}{}\mathrm{\Delta }_n)\rho =\underset{n}{}𝐗(\mathrm{\Delta }_n)\rho .$$
(34)
(DL2) For any $`\rho \tau c()`$,
$$\mathrm{Tr}[𝐗(𝐑)\rho ]=\mathrm{Tr}[\rho ].$$
(35)
(DL3) For any Borel set $`\mathrm{\Delta }`$,
$$\mathrm{Pr}\{𝐚(t)\mathrm{\Delta }\}=\mathrm{Tr}[𝐗(\mathrm{\Delta })\rho (t)].$$
(36)
(DL4) For any Borel set $`\mathrm{\Delta }`$ with $`\mathrm{Pr}\{𝐚(t)\mathrm{\Delta }\}>0`$,
$$\rho (t+\mathrm{\Delta }t|𝐚(t)\mathrm{\Delta })=\frac{𝐗(\mathrm{\Delta })\rho (t)}{\mathrm{Tr}[𝐗(\mathrm{\Delta })\rho (t)]}.$$
(37)
We call the above mapping $`𝐗:\mathrm{\Delta }𝐗(\mathrm{\Delta })`$ the operation valued measure or the operational distribution of the apparatus $`𝐀(𝐚)`$. In general, we call any bounded linear transformation on $`\tau c()`$ a superoperator for $``$. Any mapping $`\mathrm{\Delta }𝐗(\mathrm{\Delta })`$ from the Borel sets to the positive superoperators for $``$ is called a positive superoperator valued (PSV) measure if it satisfies condition (DL1). Moreover, it is called normalized if it satisfies condition (DL2). Accordingly, the operation valued measure of $`𝐀(𝐚)`$ is the normalized PSV measure satisfying (DL3) and (DL4).
The validity of the Davies-Lewis postulate for the apparatuses with indirect measurement models was previously demonstrated based on the joint probability formula in , where it is also shown that any normalized PSV measures which are realizable by indirect measurement models are completely positive and vice versa.
### B Determination of operation valued measures based on the projection postulate
In order to determine the operation valued measure corresponding to the given measuring apparatus, we need to describe the measuring process by an indirect measurement model. Then the measurement is divided into the two processes, the measuring interaction in the object-probe composite system and the probe measurement. Given the indirect measurement model, the current formulation has two arguments to determine the operation valued measure: one relies on the projection postulate (cf. for yes-no measurements) and the other relies on the joint probability formula .
In the first approach, the probe measurement is assumed explicitly to satisfy the projection postulate. Consequently, the operation valued measure is determined by the unitary operator of the measuring interaction and the projection operator derived by the projection postulate with partial trace over the probe.
The argument runs as follows. Assume that the apparatus $`𝐀(𝐚)`$ has the indirect measurement model $`(𝒦,\sigma ,U,M)`$, where the probe observable $`M`$ is purely discrete with eigenvalues $`a_n`$. Let us suppose that the measuring interaction between the system $`𝐒`$ and the apparatus $`𝐀(𝐚)`$ is turned on from $`t`$ to $`t+\mathrm{\Delta }t`$ and that the observer measures the probe observable $`M`$ at the time $`t+\mathrm{\Delta }t`$ using the apparatus $`𝐀(𝐦)`$. Let $`t+\mathrm{\Delta }t+\tau `$ be the time just after the measuring interaction is turned off between the probe $`𝐏`$ and the apparatus $`𝐀(𝐦)`$. The system $`𝐀(𝐦)`$ is considered as a subsystem of $`𝐀(𝐚)`$ including the later stages after the probe $`𝐏`$. Assume that the apparatus $`𝐀(𝐦)`$ satisfies the projection postulate and that the outcome is $`𝐦(t+\mathrm{\Delta }t)=a_n`$. Then, at the time $`t+\mathrm{\Delta }t+\tau `$ the composite system $`𝐒+𝐏`$ is in the state
$`\rho _{𝐒+𝐏}(t+\mathrm{\Delta }t+\tau |𝐦(t+\mathrm{\Delta }t)=a_n)`$ (38)
$`=`$ $`{\displaystyle \frac{(IE^M\{a_n\})\rho _{𝐒+𝐏}(t+\mathrm{\Delta }t)(IE^M\{a_n\})}{\text{Tr}[(IE^M\{a_n\})\rho _{𝐒+𝐏}(t+\mathrm{\Delta }t)]}}.`$ (39)
Since the outcome of the $`𝐀(𝐦)`$-measurement at the time $`t+\mathrm{\Delta }t`$ is interpreted as the outcome of the $`𝐀(𝐚)`$-measurement at the time $`t`$, the condition $`𝐦(t+\mathrm{\Delta }t)=a_n`$ is equivalent to the condition $`𝐚(t)=a_n`$. It follows that at the time $`t+\mathrm{\Delta }t+\tau `$ the system $`𝐒`$ is in the state
$`\rho (t+\mathrm{\Delta }t+\tau |𝐚(t)=a_n)`$ (40)
$`=`$ $`\text{Tr}_𝒦[\rho _{𝐒+𝐏}(t+\mathrm{\Delta }t+\tau |𝐦(t+\mathrm{\Delta }t)=a_n)].`$ (41)
From (25), (38), and (40), we have
$`\rho (t+\mathrm{\Delta }t+\tau |𝐚(t)=a_n)`$ (42)
$`=`$ $`{\displaystyle \frac{\text{Tr}_𝒦[(IE^M\{a_n\})U(\rho (t)\sigma )U^{}(IE^M\{a_n\})]}{\text{Tr}[(IE^M\{a_n\})U(\rho (t)\sigma )U^{}]}}.`$ (43)
By the well-known relation
$$\text{Tr}_𝒦[(IX)Y]=\text{Tr}_𝒦[Y(IX)]$$
(45)
for all $`X(𝒦)`$ and $`Y(𝒦)`$, we have
$`\text{Tr}_𝒦[(IE^M\{a_n\})U(\rho (t)\sigma )U^{}(IE^M\{a_n\})]`$ (46)
$`=`$ $`\text{Tr}_𝒦[(IE^M\{a_n\})U(\rho (t)\sigma )U^{}].`$ (47)
Hence, we have an important relation
$`\rho (t+\mathrm{\Delta }t+\tau |𝐚(t)=a_n)`$ (48)
$`=`$ $`{\displaystyle \frac{\text{Tr}_𝒦[(IE^M\{a_n\})U(\rho (t)\sigma )U^{}]}{\text{Tr}[(IE^M\{a_n\})U(\rho (t)\sigma )U^{}]}}.`$ (49)
Let $`\mathrm{\Delta }`$ be an arbitrary Borel set such that $`\mathrm{Pr}\{𝐚(t)\mathrm{\Delta }\}>0`$. Then, we have naturally
$`\rho (t+\mathrm{\Delta }t+\tau |𝐚(t)\mathrm{\Delta })`$ (50)
$`=`$ $`{\displaystyle \frac{_{a_n\mathrm{\Delta }}\mathrm{Pr}\{𝐚(t)=a_n\}\rho (t+\mathrm{\Delta }t+\tau |𝐚(t)=a_n)}{\mathrm{Pr}\{𝐚(t)\mathrm{\Delta })\}}}.`$ (51)
From (26), (48), and (50), we have
$`\rho (t+\mathrm{\Delta }t+\tau |𝐚(t)\mathrm{\Delta })`$ (53)
$`=`$ $`{\displaystyle \frac{\text{Tr}_𝒦[(IE^M(\mathrm{\Delta }))U(\rho (t)\sigma )U^{}]}{\text{Tr}[(IE^M(\mathrm{\Delta }))U(\rho (t)\sigma )U^{}]}}.`$ (54)
To obtain the final result, suppose that the $`𝐀(𝐦)`$-measurement is instantaneous, i.e., $`\tau 0`$, and that there is no interaction between $`𝐒`$ and the outside of $`𝐀(𝐚)`$ from $`t`$ to $`t+\mathrm{\Delta }t+\tau `$. Then, in this time interval the state changes of $`𝐒`$ are negligible due to the $`𝐀(𝐦)`$-measurement, the time evolution of $`𝐒`$, and the decoherence from the environment. Consequently, we have
$$\rho (t+\mathrm{\Delta }t|𝐚(t)\mathrm{\Delta })=\rho (t+\mathrm{\Delta }t+\tau |𝐚(t)\mathrm{\Delta }).$$
(55)
Therefore, we have reached the final form
$`\rho (t+\mathrm{\Delta }t|𝐚(t)\mathrm{\Delta })`$ (56)
$`=`$ $`{\displaystyle \frac{\text{Tr}_𝒦[(IE^M(\mathrm{\Delta }))U(\rho (t)\sigma )U^{}]}{\text{Tr}[(IE^M(\mathrm{\Delta }))U(\rho (t)\sigma )U^{}]}}.`$ (57)
From the above, the operation valued measure of $`𝐀(𝐚)`$ is determined by
$$𝐗(\mathrm{\Delta })\rho =\text{Tr}_𝒦[(IE^M(\mathrm{\Delta }))U(\rho \sigma )U^{}]$$
(58)
for all Borel sets $`\mathrm{\Delta }`$ and all trace class operators $`\rho `$. Furthermore, it follows easily from properties of partial trace that $`𝐗`$ satisfies conditions (DL1)–(DL4).
Obviously, the above approach explicitly excludes the possibility of measuring the probe by the apparatus not satisfying the projection postulate such as a photon counter. Accordingly, this approach cannot apply correctly to any measurements with continuous probe observables such as the position of the probe pointer, since no apparatuses measuring continuous observables satisfy the repeatability hypothesis . Moreover, the argument assumes that the probe measurement should be instantaneous.
### C Determination of operation valued measures based on the joint probability formula
In the second approach, generalizing von Neumann’s argument on repeated measurements of the same observable \[1, pp. 211–223\], it is assumed that the observer were to measure an arbitrary observable of the object system again at the time just after the measuring interaction and then considered is the joint probability distribution of the outcomes of the probe measurement and the second object measurement . By assuming that the above joint probability distribution satisfies the joint probability formula for the simultaneous measurement, we can determine the operation valued measure.
Since the joint probability formula is well formulated even in the case where the probe observable has continuous spectrum, the second approach can be applied to measurements of continuous observables. Moreover, in the case of the discrete probe observable, the second approach leads to the same operation valued measure as the first approach, so that the second approach is consistent with the first.
The argument in the second approach runs as follows. Let $`𝐀(𝐚)`$ be an apparatus described by the indirect measurement model $`(𝒦,\sigma ,U,M)`$. Suppose that the system $`𝐒`$ is measured at the time $`t`$ by the apparatus $`𝐀(𝐚)`$. Suppose that at the time $`t+\mathrm{\Delta }t`$ just after the measuring interaction, the observer were to measure an arbitrary observable $`B`$ of the same object $`𝐒`$ by an apparatus $`𝐀(𝐛)`$. The conditional probability of $`𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}`$ given $`𝐚(t)\mathrm{\Delta }`$ is the probability of $`𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}`$ in the state $`\rho (t+\mathrm{\Delta }|𝐚(t)\mathrm{\Delta })`$, so that the joint probability distribution of $`𝐚(t)`$ and $`𝐛(t+\mathrm{\Delta }t)`$ satisfies
$`\mathrm{Pr}\{𝐚(t)\mathrm{\Delta },𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}\}`$ (59)
$`=`$ $`\mathrm{Tr}[E^B(\mathrm{\Delta }^{})\rho (t+\mathrm{\Delta }t|𝐚(t)\mathrm{\Delta })]\mathrm{Pr}\{𝐚(t)\mathrm{\Delta }\}.`$ (60)
For any Borel set $`\mathrm{\Delta }`$, let $`𝐗(\mathrm{\Delta },\rho (t))`$ be the trace class operator defined by
$$𝐗(\mathrm{\Delta },\rho (t))=\mathrm{Pr}\{𝐚(t)\mathrm{\Delta }\}\rho (t+\mathrm{\Delta }t|𝐚(t)\mathrm{\Delta }).$$
(61)
From (59) we have
$`\mathrm{Pr}\{𝐚(t)\mathrm{\Delta },𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}\}`$ (62)
$`=`$ $`\mathrm{Tr}[E^B(\mathrm{\Delta }^{})𝐗(\mathrm{\Delta },\rho (t))].`$ (63)
On the other hand, by the indirect measurement model the output $`𝐚(t)`$ of this measurement is obtained by the measurement of the probe observable $`M`$ at the time $`t+\mathrm{\Delta }t`$. Let $`𝐀(𝐦)`$ be the apparatus measuring $`M`$ at the time $`t+\mathrm{\Delta }t`$. Then, the probabilistic event “$`𝐚(t)\mathrm{\Delta }`$” is equivalent to the probabilistic event “$`𝐦(t+\mathrm{\Delta }t)\mathrm{\Delta }`$” and hence we have
$`\mathrm{Pr}\{𝐚(t)\mathrm{\Delta },𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}\}`$ (64)
$`=`$ $`\mathrm{Pr}\{𝐦(t+\mathrm{\Delta }t)\mathrm{\Delta },𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}\}.`$ (65)
Since the observable $`M`$ of $`𝐏`$ and the observable $`B`$ of $`𝐒`$ are simultaneously measurable, if the $`𝐀(𝐦)`$-measurement and the $`𝐀(𝐛)`$-measurement can be considered as a simultaneous measurement, we have
$`\mathrm{Pr}\{𝐦(t+\mathrm{\Delta }t)\mathrm{\Delta },𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}\}`$
$`=`$ $`\text{Tr}[(E^B(\mathrm{\Delta }^{})E^M(\mathrm{\Delta }))\rho _{𝐒+𝐏}(t+\mathrm{\Delta }t)]`$
$`=`$ $`\text{Tr}[(E^B(\mathrm{\Delta }^{})E^M(\mathrm{\Delta }))U(\rho (t)\sigma )U^{}].`$
By the property of partial trace, we have
$`\mathrm{Pr}\{𝐦(t+\mathrm{\Delta }t)\mathrm{\Delta },𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}\}`$ (66)
$`=`$ $`\text{Tr}[E^B(\mathrm{\Delta }^{})\text{Tr}_𝒦[(IE^M(\mathrm{\Delta }))U(\rho (t)\sigma )U^{}].`$ (67)
Since $`B`$ and $`\mathrm{\Delta }^{}`$ are arbitrary, from (62)–(66) we have
$$𝐗(\mathrm{\Delta },\rho (t))=\text{Tr}_𝒦[(IE^M(\mathrm{\Delta }))U(\rho (t)\sigma )U^{}].$$
(68)
Suppose that $`\mathrm{Pr}\{𝐚(t)\mathrm{\Delta }\}>0`$. From (61), we have
$`\rho (t+\mathrm{\Delta }t|𝐚(t)\mathrm{\Delta })`$
$`=`$ $`{\displaystyle \frac{𝐗(\mathrm{\Delta },\rho (t))}{\text{Tr}[𝐗(\mathrm{\Delta },\rho (t))]}}`$
$`=`$ $`{\displaystyle \frac{\text{Tr}_𝒦[(IE^M(\mathrm{\Delta }))U(\rho (t)\sigma )U^{}]}{\text{Tr}[(IE^M(\mathrm{\Delta }))U(\rho (t)\sigma )U^{}]}}.`$
Hence, we have shown that relation (56) holds for the apparatus $`𝐀(𝐚)`$ given in this argument. Let $`𝐗`$ be the mapping $`\mathrm{\Delta }𝐗(\mathrm{\Delta })`$ defined by relation (58) for the present apparatus. Then, $`𝐗`$ satisfies conditions (DL1)–(DL2) by the properties of partial trace as before. From (68) we have
$$𝐗(\mathrm{\Delta })\rho (t)=𝐗(\mathrm{\Delta },\rho (t))$$
(69)
and hence $`𝐗`$ satisfies conditions (DL3)–(DL4). Thus, $`𝐗`$ satisfies the Davies-Lewis postulate for the apparatus $`𝐀(𝐚)`$ given above.
We have shown that the determination (58) of the operation valued measure holds without assuming the projection postulate for the probe measurement. Nevertheless, in order to justify the formula (58) generally we need to justify the joint probability formula without assuming the projection postulate. This put a serious constraint on the theoretical device to explore our problem. Indeed, because of the threat of the circular argument, the above arguments do not enable us to take advantage of the operation valued measures for the justification of the joint probability formula. In the conventional measurement theory, the similar kind of circular argument has been known as the infinite regress of the von Neumann chain. Despite the above difficulties, we shall show, in the following sections, an alternative approach without any fear of the circular argument.
## VII Statistical approach to the operation valued measures
### A Existence of the operation valued measures
In what follows, we shall prove the Davies-Lewis postulate without assuming the joint probability formula or the projection postulate.
Let us suppose that the system $`𝐒`$ is measured at the time $`t`$ by the apparatus $`𝐀(𝐚)`$ and at the time $`t+\mathrm{\Delta }t`$ immediately after this measurement an observable $`B`$ of $`𝐒`$ is measured using an apparatus $`𝐀(𝐛)`$. Then, the joint probability distribution of the outcomes of the $`A`$-measurement and the $`B`$-measurement satisfies (59). For any Borel set $`\mathrm{\Delta }`$, let $`𝐗(\mathrm{\Delta },\rho (t))`$ be the trace class operator defined by (61). Then, from (59), $`𝐗(\mathrm{\Delta },\rho (t))`$ satisfies (62). Since the input state $`\rho (t)`$ is assumed to be an arbitrary density operator, (61) defines the transformation $`𝐗(\mathrm{\Delta })`$ that maps $`\rho (t)`$ to $`𝐗(\mathrm{\Delta },\rho (t))`$. From (61) and (62), $`𝐗(\mathrm{\Delta })`$ satisfies the relations
$$𝐗(\mathrm{\Delta })\rho (t)=\mathrm{Pr}\{𝐚(t)\mathrm{\Delta }\}\rho (t+\mathrm{\Delta }t|𝐚(t)\mathrm{\Delta })$$
(70)
and
$$\mathrm{Pr}\{𝐚(t)\mathrm{\Delta },𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}\}=\text{Tr}[E^B(\mathrm{\Delta }^{})𝐗(\mathrm{\Delta })\rho (t)].$$
(71)
Suppose that the input state $`\rho (t)`$ is a mixture of density operators $`\rho _1`$ and $`\rho _2`$, i.e.,
$$\rho (t)=\alpha \rho _1+(1\alpha )\rho _2$$
(72)
where $`0<\alpha <1`$. This means that at the time $`t`$ the measured object $`𝐒`$ is sampled randomly from an ensemble of similar systems described by the density operator $`\rho _1`$ with probability $`\alpha `$ and from another ensemble described by the density operator $`\rho _2`$ with probability $`1\alpha `$. Thus we have naturally
$`\mathrm{Pr}\{𝐚(t)\mathrm{\Delta },𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}|\rho (t)=\alpha \rho _1+(1\alpha )\rho _2\}`$ (73)
$`=`$ $`\alpha \mathrm{Pr}\{𝐚(t)\mathrm{\Delta },𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}|\rho (t)=\rho _1\}`$ (75)
$`+(1\alpha )\mathrm{Pr}\{𝐚(t)\mathrm{\Delta },𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}|\rho (t)=\rho _2\},`$
where $`\mathrm{Pr}\{E|F\}`$ stands for the conditional probability of $`E`$ given $`F`$. From (71) and (73), we have
$`\mathrm{Tr}\left[E^B(\mathrm{\Delta }^{})𝐗(\mathrm{\Delta })\left[\alpha \rho _1+(1\alpha )\rho _2\right]\right]`$
$`=`$ $`\alpha \text{Tr}[E^B(\mathrm{\Delta }^{})𝐗(\mathrm{\Delta })\rho _1]`$
$`+(1\alpha )\text{Tr}[E^B(\mathrm{\Delta }^{})𝐗(\mathrm{\Delta })\rho _2]`$
$`=`$ $`\mathrm{Tr}[E^B(\mathrm{\Delta }^{})[\alpha 𝐗(\mathrm{\Delta })\rho _1+(1\alpha )𝐗(\mathrm{\Delta })\rho _2]].`$
Since $`B`$ and $`\mathrm{\Delta }^{}`$ are arbitrary, we have
$$𝐗(\mathrm{\Delta })\left[\alpha \rho _1+(1\alpha )\rho _2\right]=\alpha 𝐗(\mathrm{\Delta })\rho _1+(1\alpha )𝐗(\mathrm{\Delta })\rho _2.$$
(77)
It follows that $`𝐗(\mathrm{\Delta })`$ is an affine transformation from the space of density operators to the space of trace class operators, so that it can be extended to a unique positive superoperator .
We have proved that for any apparatus $`𝐀(𝐚)`$ measuring $`A`$ there is uniquely a family $`\{𝐗(\mathrm{\Delta })|\mathrm{\Delta }(𝐑)\}`$ of positive superoperators such that (70) and (71) hold, where $`(𝐑)`$ stands for the collection of all Borel sets.
By the countable additivity of probability, if $`\mathrm{\Delta }=_n\mathrm{\Delta }_n`$ for disjoint Borel sets $`\mathrm{\Delta }_n`$, we have
$`\mathrm{Pr}\{𝐚(t)\mathrm{\Delta },𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}\}`$ (78)
$`=`$ $`{\displaystyle \underset{n}{}}\mathrm{Pr}\{𝐚(t)\mathrm{\Delta }_n,𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}\}.`$ (79)
By (71) and (79), we have
$`\text{Tr}[E^B(\mathrm{\Delta }^{})𝐗(\mathrm{\Delta })\rho (t)]`$ $`=`$ $`{\displaystyle \underset{n}{}}\text{Tr}[E^B(\mathrm{\Delta }^{})𝐗(\mathrm{\Delta }_n)\rho (t)]`$
$`=`$ $`\text{Tr}[E^B(\mathrm{\Delta }^{}){\displaystyle \underset{n}{}}𝐗(\mathrm{\Delta }_n)\rho (t)].`$
Since $`B`$ and $`\mathrm{\Delta }^{}`$ are arbitrary, we have
$$𝐗(\mathrm{\Delta })\rho (t)=\underset{n}{}𝐗(\mathrm{\Delta }_n)\rho (t).$$
Since $`\rho (t)`$ is arbitrary, condition (DL1) holds for arbitrary density operator $`\rho `$ and hence by linearity, condition (DL1) holds for all $`\rho \tau c()`$. Conditions (DL3) and (DL4) are obvious from (61). From (DL3), we have
$$\text{Tr}[𝐗(𝐑)\rho (t)]=1.$$
Since $`\rho (t)`$ is arbitrary, condition (DL2) holds for arbitrary density operator $`\rho `$ and hence by linearity, condition (DL2) holds for all $`\rho \tau c()`$. Thus, the mapping $`𝐗:\mathrm{\Delta }𝐗(\mathrm{\Delta })\rho `$ satisfies the Davies-Lewis postulate.
It should be noted that the present derivation rely on neither the existence of the indirect measurement model, the joint probability formula, nor the projection postulate. The crucial assumption in the above argument is (73) which follows from the basic principle underlying the notion of the mixture of states. Thus, we can conclude that every measuring apparatus has the operation valued measure satisfying the Davies-Lewis postulate.
### B Basic properties of the operation valued measures
Let $`𝐀(𝐚)`$ be a measuring apparatus for the system $`𝐒`$ with the operation valued measure $`𝐗`$. Let us assume that the apparatus $`𝐀(𝐚)`$ measures an observable $`A`$. In this case, from (1) and (DL3) we have
$$\text{Tr}[𝐗(\mathrm{\Delta })\rho (t)]=\text{Tr}[E^A(\mathrm{\Delta })\rho (t)].$$
(80)
Let $`𝐗(\mathrm{\Delta })^{}`$ be the dual of $`𝐗(\mathrm{\Delta })`$. Then, we have
$$\text{Tr}[(𝐗(\mathrm{\Delta })^{}I)\rho (t)]=\text{Tr}[E^A(\mathrm{\Delta })\rho (t)].$$
Since $`\rho (t)`$ is arbitrary, we conclude
$$𝐗(\mathrm{\Delta })^{}I=E^A(\mathrm{\Delta })$$
(81)
for any Borel set $`\mathrm{\Delta }`$.
We say that a PSV measure $`𝐗`$ is $`A`$-compatible if $`𝐗`$ satisfies relation (81). By the above, the operation valued measure of the apparatus $`𝐀(𝐚)`$ measuring $`A`$ is an $`A`$-compatible PSV measure.
Now we are ready to state the following important relations for operation valued measures.
Theorem 3. Let $`A`$ be an observable and let $`𝐗`$ be an $`A`$-compatible PSV measure. Then, for any Borel set $`\mathrm{\Delta }`$ and any trace class operator $`\rho `$ we have
$`𝐗(\mathrm{\Delta })\rho `$ $`=`$ $`𝐗(𝐑)\left(E^A(\mathrm{\Delta })\rho \right)=𝐗(𝐑)\left(\rho E^A(\mathrm{\Delta })\right)`$ (82)
$`=`$ $`𝐗(𝐑)\left(E^A(\mathrm{\Delta })\rho E^A(\mathrm{\Delta })\right),`$ (83)
and for any bounded operator $`B`$ we have
$`𝐗(\mathrm{\Delta })^{}B`$ $`=`$ $`(𝐗(𝐑)^{}B)E^A(\mathrm{\Delta })=E^A(\mathrm{\Delta })𝐗(𝐑)^{}B`$ (84)
$`=`$ $`E^A(\mathrm{\Delta })(𝐗(𝐑)^{}B)E^A(\mathrm{\Delta }).`$ (85)
A proof of the above theorem was given in for the case where $`𝐗(\mathrm{\Delta })`$ is completely positive, and another proof was given in for the case where $`A`$ is discrete. The general proof necessary for the above theorem runs as follows .
Proof. Let $`C`$ be a bounded operator such that $`0CI`$ and let $`\mathrm{\Delta }(𝐑)`$. We define
$$\begin{array}{cc}A_{11}=𝐗(\mathrm{\Delta })^{}C,\hfill & A_{12}=𝐗(\mathrm{\Delta })^{}(IC),\hfill \\ A_{21}=𝐗(𝐑\mathrm{\Delta })^{}C,\hfill & A_{22}=𝐗(𝐑\mathrm{\Delta })^{}(IC),\hfill \\ P_1=E^A(\mathrm{\Delta }),\hfill & P_2=IE^A(\mathrm{\Delta }),\hfill \\ Q_1=𝐗(𝐑)^{}C,\hfill & Q_2=I𝐗(𝐑)^{}C.\hfill \end{array}$$
Then, for $`i,j=1,2`$ we have $`0A_{ij}P_i`$, so that $`[A_{ij},P_i]=[A_{ij},P_j]=0`$. It follows that $`Q_j=A_{1j}+A_{2j}`$ commutes with $`P_1`$ and $`P_2`$ as well. Thus,
$$A_{ij}=P_iA_{ij}P_iQ_j.$$
On the other hand, we have $`_{ij}A_{ij}=I`$ and $`_{ij}P_iQ_j=I`$, whence $`A_{ij}=P_iQ_j`$. It follows that
$$𝐗(\mathrm{\Delta })^{}C=E^A(\mathrm{\Delta })𝐗(𝐑)^{}C.$$
By taking adjoint, we also have
$$𝐗(\mathrm{\Delta })^{}C=(𝐗(𝐑)^{}C)E^A(\mathrm{\Delta }).$$
Since any bounded operator $`B`$ can be represented by $`B=_{n=0}^3\lambda _nC_n`$ with positive operators $`0C_nI`$ and complex numbers $`\lambda _n`$, we have
$$𝐗(\mathrm{\Delta })^{}B=E^A(\mathrm{\Delta })𝐗(𝐑)^{}B=(𝐗(𝐑)^{}B)E^A(\mathrm{\Delta })$$
for any $`\mathrm{\Delta }(𝐑)`$ and $`B()`$. By multiplying $`E^A(\mathrm{\Delta })`$ from the both sides, we also have
$$𝐗(\mathrm{\Delta })^{}B=E^A(\mathrm{\Delta })(𝐗(𝐑)^{}B)E^A(\mathrm{\Delta }).$$
Hence, relations (84) hold. Relations (82) follow easily by taking the duals of $`𝐗(\mathrm{\Delta })^{}`$ and $`\underset{¯}{\text{X}}(𝐑)^{}`$. $`\mathrm{}`$
By the above theorem, the operation valued measure $`𝐗`$ of an arbitrary apparatus $`𝐀(𝐚)`$ measuring $`A`$ is determined uniquely by the nonselective operation $`𝐓=𝐗(𝐑)`$ of $`𝐀(𝐚)`$.
Mathematical theory of PSV measures was introduced by Davies and Lewis based on conditions (DL1) and (DL2) as mathematical axioms; see also Davies . Their relations with measuring processes were established in and applied to analyzing various measuring processes in .
### C Operation valued measures of indirect measurement models
Suppose that the apparatus $`𝐀(𝐚)`$ measuring $`A`$ has the indirect measurement model $`(𝒦,\sigma ,U,M)`$. In this case, we can determine the operation valued measure $`𝐗`$ of the apparatus $`𝐀(𝐚)`$ without assuming the joint probability distribution or the projection postulate, as follows.
Let $`𝐗`$ be the operation valued measure of the apparatus $`𝐀(𝐚)`$. Then, $`𝐗`$ satisfies (DL1)–(DL4) and hence $`𝐗`$ is an $`A`$-compatible PSV measure. It follows from Theorem 3 that $`𝐗`$ satisfies
$$𝐗(\mathrm{\Delta })\rho =𝐗(𝐑)[E^A(\mathrm{\Delta })\rho ],$$
(86)
where $`\mathrm{\Delta }(𝐑)`$ and $`\rho \tau c()`$. Since $`𝐀(𝐚)`$ has the indirect measurement model $`(𝒦,\sigma ,U,M)`$, relation (27) holds. By (DL4), (27), and (28), we have
$$𝐗(𝐑)\rho (t)=\text{Tr}_𝒦[U(\rho (t)\sigma )U^{}].$$
Since $`\rho (t)`$ is arbitrary and $`𝐗(𝐑)`$ is linear, the above relation can be extended to trace class operators $`\rho `$ from density operators $`\rho (t)`$, so that we have
$$𝐗(𝐑)\rho =\text{Tr}_𝒦[U(\rho \sigma )U^{}]$$
(87)
for all $`\rho \tau c()`$. Now, we consider the expression
$$(\mathrm{\Delta })\rho =\text{Tr}_𝒦[(IE^M(\mathrm{\Delta }))U(\rho \sigma )U^{}],$$
(88)
where $`\mathrm{\Delta }(𝐑)`$ and $`\rho \tau c()`$. Then, we can show purely mathematically that the mapping $`:\mathrm{\Delta }(\mathrm{\Delta })`$ defined above is an $`A`$-compatible PSV measure satisfying
$$(𝐑)\rho =\text{Tr}_𝒦[U(\rho \sigma )U^{}].$$
(89)
Thus, $``$ satisfies the assumptions of Theorem 3, and hence we have
$$(\mathrm{\Delta })\rho =(𝐑)[E^A(\mathrm{\Delta })\rho ].$$
(90)
From (87) and (89), we have
$$(𝐑)=𝐗(𝐑)$$
(91)
and hence from (86) and (90) we have
$`(\mathrm{\Delta })\rho `$ $`=`$ $`(𝐑)[E^A(\mathrm{\Delta })\rho ]`$
$`=`$ $`𝐗(𝐑)[E^A(\mathrm{\Delta })\rho ]`$
$`=`$ $`𝐗(\mathrm{\Delta })\rho .`$
Therefore, we conclude that $`𝐗`$ satisfies (58).
From (58) and (82), we have the following expressions for $`𝐗`$:
$`𝐗(\mathrm{\Delta })\rho `$ $`=`$ $`\text{Tr}_𝒦[(IE^M(\mathrm{\Delta }))U(\rho \sigma )U^{}]`$ (93)
$`=`$ $`\mathrm{Tr}_𝒦[U(\rho (t)E^A(\mathrm{\Delta })\sigma )U^{}]`$ (94)
$`=`$ $`\mathrm{Tr}_𝒦[U(E^A(\mathrm{\Delta })\rho (t)\sigma )U^{}]`$ (95)
$`=`$ $`\mathrm{Tr}_𝒦[U(E^A(\mathrm{\Delta })\rho (t)E^A(\mathrm{\Delta })\sigma )U^{}].`$ (96)
Thus, if $`\mathrm{Pr}\{𝐚(t)\mathrm{\Delta }\}>0`$, we obtain the following relations.
$`\rho (t+\mathrm{\Delta }t|𝐚(t)\mathrm{\Delta })`$ (98)
$`=`$ $`{\displaystyle \frac{\mathrm{Tr}_𝒦[(IE^M(\mathrm{\Delta }))U(\rho (t)\sigma )U^{}]}{\mathrm{Tr}[E^A(\mathrm{\Delta })\rho (t)]}}`$ (99)
$`=`$ $`{\displaystyle \frac{\mathrm{Tr}_𝒦[U(\rho (t)E^A(\mathrm{\Delta })\sigma )U^{}]}{\mathrm{Tr}[E^A(\mathrm{\Delta })\rho (t)]}}`$ (100)
$`=`$ $`{\displaystyle \frac{\mathrm{Tr}_𝒦[U(E^A(\mathrm{\Delta })\rho (t)\sigma )U^{}]}{\mathrm{Tr}[E^A(\mathrm{\Delta })\rho (t)]}}`$ (101)
$`=`$ $`{\displaystyle \frac{\mathrm{Tr}_𝒦[U(E^A(\mathrm{\Delta })\rho (t)E^A(\mathrm{\Delta })\sigma )U^{}]}{\mathrm{Tr}[E^A(\mathrm{\Delta })\rho (t)]}}.`$ (102)
## VIII Disturbance in measurement
### A Disturbance and simultaneous measurability
Let $`B`$ be an arbitrary observable of $`𝐒`$. We say that the measurement using an apparatus $`𝐀(𝐚)`$ does not disturb the observable $`B`$ if the nonselective state change does not perturb the probability distribution of $`B`$, that is, we have
$$\mathrm{Tr}[E^B(\mathrm{\Delta })\rho (t+\mathrm{\Delta }t)]=\mathrm{Tr}[E^B(\mathrm{\Delta })e^{iH\mathrm{\Delta }t/\mathrm{}}\rho (t)e^{iH\mathrm{\Delta }t/\mathrm{}}]$$
(103)
for any Borel set $`\mathrm{\Delta }`$, where $`H`$ is the Hamiltonian of the system $`𝐒`$. The measurement is said to be instantaneous if the duration $`\mathrm{\Delta }t`$ of the measurement is negligible in the time scale of the time evolution of the system $`𝐒`$. Thus, the instantaneous measurement using the apparatus $`𝐀(𝐚)`$ does not disturb $`B`$ if and only if
$$\mathrm{Tr}[E^B(\mathrm{\Delta })\rho (t+\mathrm{\Delta }t)]=\mathrm{Tr}[E^B(\mathrm{\Delta })\rho (t)]$$
(104)
for any Borel set $`\mathrm{\Delta }`$.
Let $`𝐗`$ be the operation valued measure of the apparatus $`𝐀(𝐚)`$ and $`𝐓=𝐗(𝐑)`$ be the nonselective operation of $`𝐀(𝐚)`$. Then, from (DL4) we have
$$\rho (t+\mathrm{\Delta }t)=𝐓\rho (t)$$
(105)
and hence (104) is equivalent to
$$\text{Tr}[E^B(\mathrm{\Delta })𝐓\rho (t)]=\text{Tr}[E^B(\mathrm{\Delta })\rho (t)].$$
(106)
Let $`𝐓^{}`$ be the dual nonselective operation of $`𝐀(𝐚)`$. It follows from (106) that (104) is equivalent to
$$\text{Tr}[(𝐓^{}E^B(\mathrm{\Delta }))\rho (t)]=\text{Tr}[E^B(\mathrm{\Delta })\rho (t)].$$
(107)
Since $`\rho (t)`$ is arbitrary, (104) is equivalent to
$$𝐓^{}E^B(\mathrm{\Delta })=E^B(\mathrm{\Delta }).$$
(108)
Thus, we conclude that the instantaneous measurement using the apparatus $`𝐀(𝐚)`$ with nonselective operation $`𝐓`$ does not disturb the observable $`B`$ if and only if (108) holds for any Borel set $`\mathrm{\Delta }`$.
Now we are ready to state the answer to our problem.
Theorem 4. Let $`𝐀(𝐚)`$ be an apparatus measuring an observable $`A`$ instantaneously and let $`𝐀(𝐛)`$ be an arbitrary apparatus measuring an observable $`B`$. Then, the successive measurement using $`𝐀(𝐚)`$ and $`𝐀(𝐛)`$ is a simultaneous measurement of $`A`$ and $`B`$ if and only if $`𝐀(𝐚)`$ does not disturb $`B`$.
Proof. It suffices to show the equivalence between (3) and (108). From (71) and (86), we have
$`\mathrm{Pr}\{𝐚(t)\mathrm{\Delta },𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}\}`$
$`=`$ $`\mathrm{Tr}[E^B(\mathrm{\Delta }^{})𝐗(\mathrm{\Delta })\rho (t)]`$
$`=`$ $`\mathrm{Tr}[E^B(\mathrm{\Delta }^{})𝐗(𝐑)[\rho (t)E^A(\mathrm{\Delta })]]`$
$`=`$ $`\mathrm{Tr}[[𝐓^{}E^B(\mathrm{\Delta }^{})]\rho (t)E^A(\mathrm{\Delta })]`$
$`=`$ $`\mathrm{Tr}[E^A(\mathrm{\Delta })(𝐓^{}E^B(\mathrm{\Delta }^{}))\rho (t)]`$
Thus, the joint probability distribution of $`A`$ and $`B`$ is given by
$$\mathrm{Pr}\{𝐚(t)\mathrm{\Delta },𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}\}=\mathrm{Tr}[E^A(\mathrm{\Delta })(𝐓^{}E^B(\mathrm{\Delta }^{})\rho (t)].$$
(109)
If (108) holds, (3) follows immediately from (109). Conversely, suppose that (3) holds. By substituting $`\mathrm{\Delta }=𝐑`$ in (3), we have
$$\mathrm{Pr}\{𝐚(t)𝐑,𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}\}=\mathrm{Tr}[E^B(\mathrm{\Delta }^{})\rho (t)].$$
(110)
On the other hand, from (109) we have
$`\mathrm{Pr}\{𝐚(t)𝐑,𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}\}`$ (111)
$`=`$ $`\mathrm{Tr}[(𝐓^{}E^B(\mathrm{\Delta }^{}))\rho (t)].`$ (112)
Since $`\rho (t)`$ is arbitrary, from (110) and (111) we obtain (108). Therefore, (3) and (108) are equivalent. $`\mathrm{}`$
From Theorems 1 and 4, we can see that if the apparatus $`𝐀(𝐚)`$ measuring instantaneously an observable $`A`$ does not disturb an observable $`B`$, then $`A`$ and $`B`$ necessarily commute. Therefore, we can conclude the following statement.
Theorem 5. Every apparatus measuring an observable disturbs all the observables that do not commute with the measured observable.
### B Disturbance in indirect measurements
By (27) and by the property of the partial trace, we have
$`\mathrm{Tr}[E^B(\mathrm{\Delta })\rho (t+\mathrm{\Delta }t)]`$
$`=`$ $`\mathrm{Tr}\left[E^B(\mathrm{\Delta })\mathrm{Tr}_𝒦[U(\rho (t)\sigma )U^{}]\right]`$
$`=`$ $`\mathrm{Tr}[(E^B(\mathrm{\Delta })I)U(\rho (t)\sigma )U^{}]`$
$`=`$ $`\mathrm{Tr}[U^{}(E^B(\mathrm{\Delta })I)U(I\sigma )(\rho (t)I)]`$
$`=`$ $`\mathrm{Tr}\left[\mathrm{Tr}_𝒦[U^{}(E^B(\mathrm{\Delta })I)U(I\sigma )]\rho (t)\right].`$
Hence, (104) is equivalent to
$`\text{Tr}\left[\text{Tr}_𝒦[U^{}(E^B(\mathrm{\Delta })I)U(I\sigma )]\rho (t)\right]`$
$`=\mathrm{Tr}[E^B(\mathrm{\Delta })\rho (t)].`$
Since $`\rho (t)`$ is arbitrary, (104) is equivalent to
$$\mathrm{Tr}_𝒦[U^{}(E^B(\mathrm{\Delta })I)U(I\sigma )]=E^B(\mathrm{\Delta })$$
(113)
for any Borel set $`\mathrm{\Delta }`$.
Obviously from (113), if $`U`$ and $`BI`$ commute, i.e.,
$$[U,E^B(\mathrm{\Delta })I]=0$$
(114)
for any Borel set $`\mathrm{\Delta }`$, then the $`A`$-measurement does not disturb the observable $`B`$. However, (114) is not a necessary condition for nondisturbing measurement. In the case where $`\sigma `$ is a pure state $`\sigma =|\xi \xi |`$, from (113) we have the following theorem.
Theorem 6. Let $`𝐀(𝐚)`$ be an apparatus measuring an observable $`A`$ instantaneously with indirect measurement model $`(𝒦,|\mathrm{\Phi }\mathrm{\Phi }|,U,M)`$. The apparatus $`𝐀(𝐚)`$ does not disturb an observable $`B`$ if and only if
$$[U,E^B(\mathrm{\Delta })I]|\psi \mathrm{\Phi }=0$$
(115)
for any Borel set $`\mathrm{\Delta }`$ and any state vector $`\psi `$ of $`𝐒`$.
Proof. First, we note that in the case where $`\sigma =|\mathrm{\Phi }\mathrm{\Phi }|`$, relation (113) holds if and only if
$$\psi \mathrm{\Phi }|U^{}(E^B(\mathrm{\Delta })I)U|\psi \mathrm{\Phi }=\psi |E^B(\mathrm{\Delta })|\psi $$
(116)
holds for any state vector $`\psi `$. Suppose that (115) holds. We have
$$U(E^B(\mathrm{\Delta })I)|\psi \mathrm{\Phi }=(E^B(\mathrm{\Delta })I)U|\psi \mathrm{\Phi }.$$
Multiplying $`U^{}`$ from the left, we have
$$(E^B(\mathrm{\Delta })I)|\psi \mathrm{\Phi }=U^{}(E^B(\mathrm{\Delta })I)U|\psi \mathrm{\Phi },$$
and hence, we have (116). Thus, if (115) holds for any Borel set $`\mathrm{\Delta }`$ and any state vector $`\psi `$ then $`𝐀(𝐚)`$ does not disturb $`B`$. Conversely, suppose that $`𝐀(𝐚)`$ does not disturb $`B`$. Then, from (113) with $`\sigma =|\mathrm{\Phi }\mathrm{\Phi }|`$, we have
$`\varphi ^{}\mathrm{\Phi }|U^{}(E^B(\mathrm{\Delta })I)U|\varphi \mathrm{\Phi }`$
$`=`$ $`\varphi ^{}\mathrm{\Phi }|E^B(\mathrm{\Delta })I|\varphi \mathrm{\Phi }`$
for any vectors $`\varphi ,\varphi ^{}`$. Let $`\psi `$ be a state vectors. Put $`|\varphi =|\psi `$ and $`|\varphi ^{}=E^B(\mathrm{\Delta })|\psi `$, we have
$`\psi \mathrm{\Phi }|(E^B(\mathrm{\Delta })I)U^{}(E^B(\mathrm{\Delta })I)U|\psi \mathrm{\Phi }`$ (117)
$`=\psi \mathrm{\Phi }|E^B(\mathrm{\Delta })I|\psi \mathrm{\Phi }.`$ (118)
By taking complex conjugate, we have
$`\psi \mathrm{\Phi }|U^{}(E^B(\mathrm{\Delta })I)U(E^B(\mathrm{\Delta })I)|\psi \mathrm{\Phi }`$ (119)
$`=\psi \mathrm{\Phi }|E^B(\mathrm{\Delta })I|\psi \mathrm{\Phi }.`$ (120)
From (116) – (119), we have
$`[E^B(\mathrm{\Delta })IU^{}(E^B(\mathrm{\Delta })I)U]|\psi \mathrm{\Phi }^2`$
$`=`$ $`\psi \mathrm{\Phi }|E^B(\mathrm{\Delta })I|\psi \mathrm{\Phi }`$
$`\psi \mathrm{\Phi }|(E^B(\mathrm{\Delta })I)U^{}(E^B(\mathrm{\Delta })I)U|\psi \mathrm{\Phi }`$
$`\psi \mathrm{\Phi }|U^{}(E^B(\mathrm{\Delta })I)U(E^B(\mathrm{\Delta })I)|\psi \mathrm{\Phi }`$
$`+\psi \mathrm{\Phi }|U^{}(E^B(\mathrm{\Delta })I)U|\psi \mathrm{\Phi }`$
$`=`$ $`0.`$
Thus, we have
$$(E^B(\mathrm{\Delta })I)|\psi \mathrm{\Phi }=U^{}(E^B(\mathrm{\Delta })I)U|\psi \mathrm{\Phi }.$$
Multiplying $`U`$ from the left, we have
$$[U(E^B(\mathrm{\Delta })I)(E^B(\mathrm{\Delta })I)U]|\psi \mathrm{\Phi }=0,$$
and hence we have (115). Therefore, we conclude that if $`𝐀(𝐚)`$ does not disturb $`B`$, then (115) holds for any Borel set $`\mathrm{\Delta }`$ and any state vector $`\psi `$ of $`𝐒`$. $`\mathrm{}`$
## IX Local Measurements of Observables of Two Entangled System
If the two observables to be measured belong to two different subsystems respectively, then they commute each other and the measurement of one is not considered to disturb the other in general, so that the result obtained in the preceding section applies to this situation. The purpose of this section is to state this fact in the rigorous language.
Let $`C`$ be an observable of an system $`𝐒_1`$ with Hilbert space $`_1`$ and $`D`$ an observable of another system $`𝐒_2`$ with Hilbert space $`_2`$. Suppose that the composite system $`𝐒=𝐒_1+𝐒_2`$ is in the state $`\rho (t)`$ at the time $`t`$. Let us suppose that one measures the observable $`C`$ at the time $`t`$ using an apparatus $`𝐀(𝐚)`$ and that at the time, $`t+\mathrm{\Delta }t`$, just after the $`C`$-measurement one measures $`D`$ using any apparatus $`𝐀(𝐛)`$ measuring $`D`$. We assume that after the time $`t`$ there is no interaction between $`𝐒_1`$ and $`𝐒_2`$.
First, we shall consider the case where the measurement of $`C`$ satisfies the projection postulate. In this case, in the composite system $`𝐒_{12}`$ the observable $`A=CI_2`$ is measured at the time $`t`$ and the observable $`B=I_1D`$ is measured immediately after the $`A`$-measurement, where $`I_1`$ and $`I_2`$ are the identity operators on $`_1`$ and $`_2`$, respectively. From Theorem 2 the joint probability distribution satisfies
$`\mathrm{Pr}\{𝐚(t)\mathrm{\Delta },𝐛(t+\mathrm{\Delta }t)\mathrm{\Delta }^{}\}`$ (121)
$`=`$ $`\text{Tr}\left[\left(E^C(\mathrm{\Delta })E^D(\mathrm{\Delta }^{})\right)\rho (t)\right].`$ (122)
In order to compare this result with the argument given by EPR , let us consider the special case where $`C`$ and $`D`$ are nondegenerate observables in their own subsystems and the initial state $`\rho (t)`$ is a pure state. In this case, the state $`\rho (t)`$ is represented by a state vector $`\mathrm{\Psi }(t)`$ in the Hilbert space $`=_1_2`$ as
$$\rho (t)=|\mathrm{\Psi }(t)\mathrm{\Psi }(t)|.$$
Let us suppose that the observables $`C`$ and $`D`$ have the spectral decompositions
$`C`$ $`=`$ $`{\displaystyle \underset{n}{}}a_n|\varphi _n\varphi _n|,`$
$`D`$ $`=`$ $`{\displaystyle \underset{m}{}}b_m|\xi _m\xi _m|.`$
EPR expand $`\mathrm{\Psi }(t)`$ using the basis $`\{\varphi _n\}`$ of $`_1`$ as
$$\mathrm{\Psi }(t)=\underset{n}{}|\varphi _n\eta _n,$$
(123)
where $`\eta _n`$ are uniquely determined vectors in $`_2`$ not necessarily orthogonal and, according to EPR, are to be regarded merely as the coefficients of the expansion of $`\mathrm{\Psi }(t)`$ into a series of orthogonal vectors $`\varphi _n`$. Then, EPR considered the process of “reduction of the wave packet”
$$\underset{n}{}|\varphi _n\eta _nN|\varphi _n\eta _n,$$
(124)
where $`N`$ is the normalization constant determined up to a phase factor by
$$N=\varphi _n\eta _n^1,$$
(125)
and stated that the state after the measurement conditional upon the outcome $`𝐚(t)=a_n`$ is determined as
$$\mathrm{\Psi }(t+\mathrm{\Delta }t|𝐚(t)=a_n)=N|\varphi _n\eta _n,$$
(126)
where
$`|\mathrm{\Psi }(t+\mathrm{\Delta }t|𝐚(t)=a_n)\mathrm{\Psi }(t+\mathrm{\Delta }t|𝐚(t)=a_n)|`$ (127)
$`=`$ $`\rho (t+\mathrm{\Delta }t|𝐚(t)=a_n).`$ (128)
From this, we have the joint probability formula
$`\mathrm{Pr}\{𝐚(t)=a_n,𝐛(t+\mathrm{\Delta }t)=b_m\}`$ (129)
$`=`$ $`|\varphi _n\xi _m|\mathrm{\Psi }(t)|^2,`$ (130)
which is a special case of (121).
Now, let us show that the EPR argument is equivalent with the argument based on the projection postulate for the $`A`$-measurement. From the projection postulate, if the outcome of the $`𝐀(𝐚)`$-measurement is $`a_n`$, the state of the composite system at the time just after the measurement is
$$\mathrm{\Psi }(t+\mathrm{\Delta }t|𝐚(t)=a_n)=\frac{(|\varphi _n\varphi _n|I_2)\mathrm{\Psi }(t)}{(|\varphi _n\varphi _n|I_2)\mathrm{\Psi }(t)}.$$
(131)
Then, from (123) we have
$$(|\varphi _n\varphi _n|I_2)\mathrm{\Psi }(t)=|\varphi _n\eta _n.$$
(132)
Thus, we have shown that (126) is the consequence from the projection postulate (131).
In the following, we shall consider the general case. For instance, consider the case where the $`A`$-measurement leaves the system $`𝐒_1`$ in a fixed state $`\varphi _1`$ independent of the outcome such as the vacuum state after photon counting. Does (121) hold even in this case? The answer to this question might depend on the way of measuring $`A`$. However, if the measurement of $`A`$ is carried out so as not to affect the system $`𝐒_2`$, then from the result in the preceding section we will be able to conclude relation (121). In order to ensure that the measurement of $`A`$ does not affect the system $`𝐒_2`$, we introduce the following condition.
We will say that the apparatus $`𝐀(𝐚)`$ measuring $`A`$ is local in the system $`𝐒_1`$ if the measuring interaction is confined in the system $`𝐒_1`$ and the apparatus $`𝐀(𝐚)`$, as formulated precisely as follows. Let $`𝒦`$ be the Hilbert space of the probe $`𝐏`$ in the apparatus $`𝐀(𝐚)`$ and suppose that $`𝐏`$ is prepared in the state $`\sigma `$ at the time $`t`$ of measurement and let $`U`$ be the unitary operator of $`𝒦_1_2`$ representing the time evolution of the composite system $`𝐒+𝐏`$. Then, the apparatus $`𝐀(𝐚)`$ is said to be local in the system $`𝐒_1`$ if we have
$$[U,I_1XI_𝒦]=0,$$
(133)
for any bounded operator $`X`$ on $`_2`$, where $`I_𝒦`$ is the identity on $`𝒦`$.
Theorem 7. Suppose that the composite system $`𝐒=𝐒_1+𝐒_2`$ is in the state $`\rho (t)`$ at the time $`t`$ of measurement. Let $`C`$ and $`D`$ be observables of $`𝐒_1`$ and $`𝐒_2`$, respectively. If the apparatus $`𝐀(𝐚)`$ measuring $`A=CI_2`$ instantaneously is local in the system $`𝐒_1`$ then (121) holds.
Proof. Let $`\sigma `$ be the state of the probe at $`t`$. From Theorem 5 it suffices to show that $`𝐀(𝐚)`$ does not disturb the observable $`B=ID`$. By assumption, we have
$`[U,E^B(\mathrm{\Delta })I_𝒦]`$ $`=`$ $`[U,I_1E^D(\mathrm{\Delta })I_𝒦]`$
$`=`$ $`0`$
for all Borel set $`\mathrm{\Delta }`$. Thus, relation (114) holds, so that $`𝐀(𝐚)`$ does not disturb the observable $`B=I_1D`$. Therefore, equation (121) follows from Theorem 4. $`\mathrm{}`$
From the above theorem, we have also the following statement: Any pair of local instantaneous measuring apparatuses of $`A=CI_2`$ and $`B=I_1D`$ satisfies the joint probability formula
$$\mathrm{Pr}\{𝐚(t)\mathrm{\Delta },𝐛(t)\mathrm{\Delta }^{}\}=\mathrm{Tr}[(E^C(\mathrm{\Delta })E^D(\mathrm{\Delta }^{}))\rho (t)],$$
(134)
regardless of the order of the measurement, where we identify $`t`$ with $`t+\mathrm{\Delta }t`$.
In the EPR paper , the so called EPR correlation is derived theoretically under the assumption that the pair of measurements satisfies the projection postulate, but the present result concludes that the EPR correlation holds for any pair of local instantaneous measurements as experiments have already suggested.
## X Minimum disturbing measurements
Classical measurements are usually considered to disturb no measured systems. This does not mean, however, that no classical measurement disturbs the system but means that among all the possible measurement the minimum disturbing measurement does not disturb the system in principle. In this section, we shall introduce the notion of the minimum disturbing measurement in quantum mechanics and show that this is equivalent to the measurement satisfying the projection postulate.
For an apparatus $`𝐀(𝐱)`$, we denote by $`𝒟(𝐱)`$ the set of observables that are disturbed by $`𝐀(𝐱)`$, i.e., $`𝒟(𝐱)`$ is the set of observables $`B`$ such that $`𝐓^{}E^B(\mathrm{\Delta })E^B(\mathrm{\Delta })`$ for some Borel set $`\mathrm{\Delta }`$, where $`𝐓`$ is the nonselective operation of $`𝐀(𝐱)`$. Let $`A`$ be an observable of the system $`𝐒`$ and let $`𝐀(𝐚)`$ be an apparatus measuring $`A`$ instantaneously. The apparatus $`𝐀(𝐚)`$ is called minimum disturbing if $`𝒟(𝐚)𝒟(𝐱)`$ for any apparatus $`𝐀(𝐱)`$ measuring $`A`$ instantaneously. Then, we have the following statement.
Theorem 8. Let $`𝐀(𝐚)`$ be an apparatus measuring a discrete observable $`A`$ instantaneously. The apparatus $`𝐀(𝐚)`$ is minimum disturbing if and only if $`𝐀(𝐚)`$ satisfies the projection postulate.
Proof. Let $`𝒞(A)`$ be the set of observables that do not commute with $`A`$. From Theorem 5, we have
$$𝒞(A)^c𝒟(𝐱)$$
(135)
for any apparatus $`𝐀(𝐱)`$ measuring $`A`$ instantaneously, where <sup>c</sup> stands for the complement in the set of observables. Let $`𝐀(𝐚)`$ be an apparatus measuring $`A`$ instantaneously. Suppose that $`𝐀(𝐚)`$ satisfies the projection postulate. Then, from Theorem 2 we have
$$𝒟(𝐚)𝒞(A)^c,$$
(136)
and hence from (135) we conclude that $`𝐀(𝐚)`$ is minimum disturbing and
$$𝒟(𝐚)=𝒞(A)^c.$$
(137)
Conversely, suppose that $`𝐀(𝐚)`$ is minimum disturbing. We have an indirect measurement model that measures $`A`$ instantaneously and satisfying the projection postulate . Hence, there is an apparatus $`𝐀(𝐱)`$ measuring $`A`$ instantaneously such that $`𝒟(𝐱)=𝒞(A)^c`$. By assumption, $`𝐀(𝐚)`$ is minimum disturbing, so that $`𝒟(𝐚)=𝒞(A)^c`$. Then, the operation valued measure $`𝐗`$ of $`𝐀(𝐚)`$ is such that $`𝐗(𝐑)^{}E^B(\mathrm{\Delta }^{})=E^B(\mathrm{\Delta }^{})`$ for all $`B𝒞(A)`$ and $`\mathrm{\Delta }^{}(𝐑)`$. Thus, we have
$`\text{Tr}[E^B(\mathrm{\Delta }^{})𝐗(\mathrm{\Delta })\rho (t)]`$
$`=`$ $`\text{Tr}[[𝐗(\mathrm{\Delta })^{}E^B(\mathrm{\Delta }^{})]\rho (t)]`$
$`=`$ $`{\displaystyle \underset{a\mathrm{\Delta }}{}}\text{Tr}[[𝐗\{a\}^{}E^B(\mathrm{\Delta }^{})]\rho (t)]`$
$`=`$ $`{\displaystyle \underset{a\mathrm{\Delta }}{}}\text{Tr}[(E^A\{a\}[𝐗(𝐑)^{}E^B(\mathrm{\Delta }^{})]E^A\{a\})\rho (t)]`$
$`=`$ $`{\displaystyle \underset{a\mathrm{\Delta }}{}}\text{Tr}[E^A\{a\}E^B(\mathrm{\Delta }^{})E^A\{a\}\rho (t)]`$
$`=`$ $`\text{Tr}[E^B(\mathrm{\Delta }^{}){\displaystyle \underset{a\mathrm{\Delta }}{}}E^A\{a\}\rho (t)E^A\{a\}].`$
Since $`B`$ and $`\mathrm{\Delta }^{}`$ are arbitrary, we have
$$𝐗(\mathrm{\Delta })\rho (t)=\underset{a\mathrm{\Delta }}{}E^A\{a\}\rho (t)E^A\{a\}$$
and hence
$$\rho (t+\mathrm{\Delta }t|𝐚\mathrm{\Delta })=\frac{_{a\mathrm{\Delta }}E^A\{a\}\rho (t)E^A\{a\}}{\text{Tr}[E^A(\mathrm{\Delta })\rho ]}.$$
Thus, $`𝐀(𝐚)`$ satisfies the projection postulate. $`\mathrm{}`$
We refer to for different approaches to the minimum disturbance condition. The present approach leads to the simplest characterization of the measurements satisfying the projection postulate, which can be called eventually as the minimum disturbing measurements.
## XI Concluding remarks
As anticipated from the ordinary interpretation of the uncertainty principle or the principle of complementarity formulated by the noncommutativity of observables, every measurement of an observable disturbs every observable that does not commute with the measured observable. It should be noticed, however, that this does not imply a prevailing interpretation of the Heisenberg uncertainty principle that the measurement of the position with accuracy $`ϵ`$ must bring about an indeterminacy $`\eta =\mathrm{}/2ϵ`$ in the value of the position \[1, p. 239\]. In fact, we can construct an indirect measurement model of the postion measurement that counters the above statement ; this model has the complete accuracy $`ϵ=0`$ but disturbs the momentun arbitrarily small if the input state is arbitarily close to the momentum eigenstate. This example suggests that the relation between the accuracy and the disturbance is more complicated than the relation $`ϵ\eta \mathrm{}/2`$ suggested by the Robertson uncertainty relation . The detailed investigation will be presented in the forthcomming article.
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# Untitled Document
Cylindrical sources in full Einstein and Brans-Dicke gravity
Andrés Arazi<sup>1</sup><sup>1</sup>1Electronic mail: arazi@tandar.cnea.gov.ar and Claudio Simeone<sup>2</sup><sup>2</sup>2Electronic mail: simeone@tandar.cnea.gov.ar
Departamento de Física, Comisión Nacional de Energía Atómica
Av. del Libertador 8250, 1429 Buenos Aires, Argentina
ABSTRACT
It was shown by Hiscock that the energy-momentum tensor commonly used to model local cosmic strings in linearized Einstein gravity can be extended and used in the full theory, obtaining a metric in the exterior of the source with the same deficit angle. Here we show that this tensor is an exception within a family for which a static solution does not exist in full Einstein nor in Brans-Dicke gravity.
KEY WORDS: Static cylindrical solutions; cosmic string; Brans-Dicke.
PACS numbers: 98.80.Cq 04.50.+h 11.27.+d
1. INTRODUCTION
In a gauge theory, spontaneous symmetry breaking of a complex scalar field leads to cylindrical topological defects known as local cosmic strings . The gravitational effects of such objects are of particular interest since they are considered as possible “seeds” for galaxy formation and gravitational lenses. The metric around a local string was first calculated by Vilenkin in the linear approximation of general relativity. Local strings are characterized by having an energy-momentum tensor whose only non null components are $`T_t^t=T_z^z`$. As linearized Einstein equations are formally analogous to Maxwell equations, the exterior solution does not depend on the radial distribution of the source. Hence, a Dirac $`\delta `$ was used to approximate the radial distribution of the energy-momentum tensor for a cosmic string along the $`z`$ axis:
$$\stackrel{~}{T}_\mu ^\nu \delta (x)\delta (y)T_\mu ^\nu (x,y)𝑑x𝑑y=\delta (x)\delta (y)\text{diag}(\mu ,0,0,\mu ),$$
(1)
where $`\mu `$ is the linear mass density. Under this assumptions, Vilenkin obtained a spacetime metric which is flat but with a deficit angle $`\mathrm{\Delta }\phi =8\pi G\mu `$, up to first order in $`G\mu `$ (in GUT strings $`G\mu 10^6`$). Since this metric has $`g_{tt}=1`$, i.e. the Newtonian potential is null, rest particles are not affected by the string.
Some years later, Hiscock , motivated by the possibility of theories which may lead to values of $`G\mu `$ closer to one, showed that Vilenkin’s results are actually valid to all orders in $`G\mu `$. As a source, he considered a thick cylinder of radius $`a`$ with uniform tension and linear mass density, whose tensor is
$$T_\mu ^\nu (x,y)=\text{diag}(\mu ,0,0,\mu )\frac{\theta (ra)}{a^2}.$$
(2)
He solved full Einstein equations in the interior and matched the resulting static metric with the vacuum solution for the exterior.
On the other hand, from the point of view of structure formation it is important to determine whether an object interacts with rest particles. Vachaspati and Vilenkin obtained a metric with non-null Newtonian potential considering a source whose tensor has $`T_z^z=𝒯`$ (effective tension) different from $`T_t^t=`$ (energy per unit length). For this, they again considered the approximation of an infinitesimally thin ($`\delta `$type) source and worked within linearized Einstein gravity. Similar results were obtained for such a source in linearized Brans-Dicke gravity .
In the present work we show that the case $`T_t^t=T_z^z`$ solved by Hiscock is an exception: thick sources with energy–momentum tensor
$$T_\mu ^\nu =\text{diag}(,0,0,𝒯)F(r)$$
(3)
do not admit a static solution in full Einstein nor in Brans–Dicke theories of gravitation. In (3) $`F(r)`$ is any distribution function whose integral over the string transverse section is equal to unity. In this general case, we can obtain the static metrics for the exterior by solving full Einstein and Brans–Dicke vacuum equations for the most general metric with cylindrical symmetry. However, we find that static interior solutions do not exist in either theories.
2. GENERAL RELATIVITY
A. Weak field, $`\delta `$ source
Vachaspati and Vilenkin solved the linearized Einstein equations to obtain the metric in the exterior of an infinitesimally thin source described by the energy-momentum tensor
$$\stackrel{~}{T}_\mu ^\nu =\text{diag}(,0,0,𝒯)\delta (x)\delta (y).$$
(4)
They found a solution which in cylindrical coordinates has the form
$`ds^2`$ $`=`$ $`[1+4G(𝒯)\mathrm{ln}(r/r_0)]dt^2`$ (5)
$`[14G(+𝒯)\mathrm{ln}(r/r_0)](dr^2+r^2d\phi ^2)`$
$`[14G(𝒯)\mathrm{ln}(r/r_0)]dz^2`$
where $`r_0`$ is a constant of integration. As $`𝒯`$ we have $`g_{00}1`$ and, differing from the case $`T_t^t=T_z^z`$, there is an interaction with rest particles.
B. Full equations, finite cylindrical source
We shall start from the most general static metric with cylindrical symmetry :
$$ds^2=e^{2(KU)}(dt^2dr^2)e^{2U}W^2d\phi ^2e^{2U}dz^2,$$
(6)
where $`K`$, $`U`$ and $`W`$ are $`r`$–dependent functions. In terms of these functions, the full Einstein equations for the energy-momentum tensor of equation (3) take the form:
$$\frac{W^{\prime \prime }}{W}+K^{}\frac{W^{}}{W}U_{}^{}{}_{}{}^{2}=8\pi GF(r)e^{2(KU)},$$
(7)
$$K^{}\frac{W^{}}{W}U_{}^{}{}_{}{}^{2}=0,$$
(8)
$$K^{\prime \prime }+U_{}^{}{}_{}{}^{2}=0,$$
(9)
$$\frac{W^{\prime \prime }}{W}+2U^{\prime \prime }+2U^{}\frac{W^{}}{W}K^{\prime \prime }U_{}^{}{}_{}{}^{2}=8\pi G𝒯F(r)e^{2(KU)},$$
(10)
where primes denote derivatives with respect to $`r`$. In the exterior of the source ( $`F(r)=0`$ ) these equations lead to the Weyl vacuum metric which, with our coordinates choice, has the form
$$ds^2=r^{2d(d1)}(dt^2dr^2)W_0^2r^{2d}r^2d\phi ^2r^{2d}dz^2.$$
(11)
Since in a cylindrically symmetric problem the exterior solution is not independent of the interior metric, as it happens in a spherical problem, the integration constants $`W_0`$ and $`d`$ should be determined by matching both metrics in the boundary. With $`d=0`$ or $`d=2`$ the metric (11) becomes Lorentz invariant in the $`z`$ direction; the case $`d=0`$ is the one solved by Hiscock.
We shall show, however, that a static interior solution does not exist. From (9) and (10) we have
$$\frac{W^{\prime \prime }}{W}+2U^{\prime \prime }+2U^{}\frac{W^{}}{W}=8\pi G𝒯F(r)e^{2(KU)},$$
(12)
and from (7) and (8)
$$\frac{W^{\prime \prime }}{W}=8\pi GF(r)e^{2(KU)},$$
(13)
so that
$$U^{\prime \prime }+U^{}\frac{W^{}}{W}=\frac{1}{2}\frac{W^{\prime \prime }}{W}\left(\frac{𝒯}{}\right).$$
(14)
The conservation equation
$$T_{\rho }^{\sigma }{}_{;\sigma }{}^{}=\frac{}{x^\sigma }\left(T_\rho ^\sigma \sqrt{g}\right)\frac{1}{2}\sqrt{g}\frac{g_{\sigma \tau }}{x^\rho }T^{\sigma \tau }=0$$
(15)
yields
$$(K^{}U^{})+U^{}𝒯=0.$$
(16)
Using this equation we can write $`K^{}=U^{}(𝒯)/`$ and $`K^{\prime \prime }=U^{\prime \prime }(𝒯)/`$, and then from (8) and (9) we obtain
$$\left(U^{\prime \prime }+U^{}\frac{W^{}}{W}\right)\left(\frac{𝒯}{}\right)=0.$$
(17)
In the particular case $`𝒯=`$ these equations are compatible and yield the interior metric found by Hiscock ($`g_{\phi \phi }=(a^2/8\pi G)\mathrm{sin}^2(\sqrt{8\pi G}r/a)`$ ). However, for $`𝒯`$ equations (14) and (17) yield $`W^{\prime \prime }=0`$. If so, equation (13) gives $`=0`$, which means that there is no string. Hence an interior static solution cannot exist in the full theory.
3. BRANS-DICKE GRAVITY
In the framework of present unified theories a scalar field should exist besides the metric of the spacetime. Scalar-tensor theories of gravitation would be important when studying the early universe, where it is supposed the coupling of the matter to the scalar field could be nonnegligible. Topological defects are produced in phase transitions in the early universe, so that it seems natural to study them in a scalar-tensor theory of gravitation as that of Brans and Dicke .
In Brans-Dicke theory matter and nongravitational fields generate a long-range scalar field $`\varphi `$, which, together with them, acts as a source of gravitational field. The field $`\varphi `$ is a solution of the equation
$$\varphi _{;\sigma }^{;\sigma }=\frac{1}{\sqrt{g}}\frac{}{x^\sigma }\left(\sqrt{g}g^{\sigma \tau }\frac{\varphi }{x^\tau }\right)=\frac{8\pi T}{2\omega +3}$$
(18)
where $`T=\delta _\mu ^\nu T_\nu ^\mu `$ and $`\omega `$ is a dimensionless constant; the metric equations replacing those of General Relativity are
$$R_{\mu \nu }\frac{1}{2}g_{\mu \nu }R=8\pi \frac{T_{\mu \nu }}{\varphi }+\frac{\omega }{\varphi ^2}\varphi _{,\mu }\varphi _{,\nu }\frac{\omega }{2\varphi ^2}g_{\mu \nu }\varphi _{,\alpha }\varphi ^{,\alpha }+\frac{1}{\varphi }\varphi _{,\mu ;\nu }\frac{1}{\varphi }g_{\mu \nu }\varphi _{;\sigma }^{;\sigma }.$$
(19)
A. Weak field, $`\delta `$ source
In the linearized approximation the $`\varphi `$ field is expanded as $`\varphi \varphi _0+\xi =G^1+\xi `$ so that the equations for the metric and $`\varphi `$ are
$$R_{\mu \nu }^{(1)}=8\pi G\left(T_{\mu \nu }\frac{\omega +1}{2\omega +3}\eta _{\mu \nu }T\right)+G\xi _{,\mu ,\nu },\mathrm{}\varphi =\mathrm{}\xi =\frac{8\pi T}{2\omega +3}.$$
(20)
In the Brans-Dicke gauge $`(h_\mu ^\nu \delta _\mu ^\nu h)_{,\mu }=G\xi _{,\nu }`$ the perturbation $`h_{\mu \nu }`$ decouples from $`\varphi `$ and the equations for the metric take the simple form
$$^2h_{\mu \nu }=16\pi G\left(T_{\mu \nu }\frac{\omega +1}{2\omega +3}\eta _{\mu \nu }T\right).$$
(21)
Solving this equations for the energy-momentum tensor (4) we obtain
$`ds^2`$ $`=`$ $`\left[1+{\displaystyle \frac{8G}{2\omega +3}}[(\omega +2)𝒯(\omega +1)]\mathrm{ln}(r/r_0)\right]dt^2`$ (22)
$`\left[18G(+𝒯)\left({\displaystyle \frac{\omega +1}{2\omega +3}}\right)\mathrm{ln}(r/r_0)\right](dr^2+r^2d\phi ^2)`$
$`\left[1{\displaystyle \frac{8G}{2\omega +3}}[(\omega +1)𝒯(\omega +2)]\mathrm{ln}(r/r_0)\right]dz^2.`$
In the limit $`\omega \mathrm{}`$ the metric (5) is recovered. If we write
$$ds^2=g_{00}\left(dt^2+\underset{i=1}{\overset{3}{}}g^{00}g_{ii}(dx^i)^2\right)$$
and redefine the radial coordinate by
$$(18G\mathrm{ln}(r/r_0))r^2=(18G)\rho ^2,(18G\mathrm{ln}(r/r_0))dr^2d\rho ^2,$$
we can put the metric in the form
$`ds^2`$ $`=`$ $`(1+{\displaystyle \frac{8G}{2\omega +3}}[(\omega +2)𝒯(\omega +1)]\mathrm{ln}(\rho /\rho _0))\times `$ (23)
$`\times \left(dt^2d\rho ^2(18G)\rho ^2d\phi ^2[18G(𝒯)\mathrm{ln}(\rho /\rho _0)]dz^2\right).`$
In a plane perpendicular to the $`z`$ axis the metric is conformal to one with a deficit angle $`\mathrm{\Delta }=8\pi G`$ which does not depend on the Brans-Dicke constant $`\omega `$.
B. Full equations, finite cylindrical source
For the source (3), the Brans-Dicke equations (19) read
$$\frac{W^{\prime \prime }}{W}+K^{}\frac{W^{}}{W}U_{}^{}{}_{}{}^{2}=8\pi \frac{}{\varphi }F(r)e^{2(KU)}+\frac{\omega \varphi _{}^{}{}_{}{}^{2}}{2\varphi ^2}(K^{}U^{})\frac{\varphi ^{}}{\varphi }+\frac{1}{\varphi }\left(\varphi ^{\prime \prime }+\varphi ^{}\frac{W^{}}{W}\right)$$
(24)
$$K^{}\frac{W^{}}{W}U_{}^{}{}_{}{}^{2}=\frac{\omega \varphi _{}^{}{}_{}{}^{2}}{2\varphi ^2}(K^{}U^{})\frac{\varphi ^{}}{\varphi }\frac{W^{}\varphi ^{}}{W\varphi }$$
(25)
$$K^{\prime \prime }+U_{}^{}{}_{}{}^{2}=\frac{\omega \varphi _{}^{}{}_{}{}^{2}}{2\varphi ^2}U^{}\frac{\varphi ^{}}{\varphi }\frac{\varphi ^{\prime \prime }}{\varphi }$$
(26)
$$\frac{W^{\prime \prime }}{W}+2U^{\prime \prime }+2U^{}\frac{W^{}}{W}K^{\prime \prime }U_{}^{}{}_{}{}^{2}=8\pi \frac{𝒯}{\varphi }F(r)e^{2(KU)}+\frac{\omega \varphi _{}^{}{}_{}{}^{2}}{2\varphi ^2}+\frac{1}{\varphi }\left(\varphi ^{\prime \prime }+\varphi ^{}\frac{W^{}}{W}\right)U^{}\frac{\varphi ^{}}{\varphi }$$
(27)
and the equation (18) for the $`\varphi `$ field takes the form
$$\varphi ^{\prime \prime }+\varphi ^{}\frac{W^{}}{W}=8\pi \left(\frac{+𝒯}{2\omega +3}\right)F(r)e^{2(KU)}.$$
(28)
We shall first find the metric in the exterior of the source by solving these equations for vacuum, that is, with $`F(r)=0`$. For this case, we inmediately see that $`\varphi =1/G`$ is a particular solution of (28) which leads to the equations of general relativity. To find the general solution we shall subtract equation (24) from (25) to get
$$W\varphi =ar+b.$$
(29)
Adding (25) and (26) and using (28) we obtain
$$K^{}=c\frac{a}{W\varphi }$$
while adding (26) and (27), with the use of (28) and (29), we get
$$U^{}=d\frac{a}{W\varphi }.$$
This yields
$$U=d\mathrm{ln}\left|\frac{ar+b}{e}\right|,K=c\mathrm{ln}\left|\frac{ar+b}{f}\right|,$$
$$W=g(ar+b)^n,\varphi =\frac{1}{g}(ar+b)^{1n},$$
(30)
where $`a\mathrm{}g,n`$ are integration constants; from (24) or (25) the relation $`c=d(d+1n)+\frac{1}{2}\omega (1n)^2+n(n1)`$ can be obtained. The resulting metric can therefore be put in the form
$$ds^2=r^{2d(dn)+(\omega +2n)(n1)}(dt^2dr^2)W_0^2r^{2(nd)}d\phi ^2r^{2d}dz^2.$$
(31)
Choosing $`n=1`$ (which corresponds to $`\varphi =constant`$) the Weyl metric of equation (11) is recovered.
Now let us study the possibility of obtaining a static solution for the interior of the source. From (25), (26) and (28) we obtain
$$K^{\prime \prime }+K^{}\left(\frac{W^{}}{W}+\frac{\varphi ^{}}{\varphi }\right)=\frac{8\pi }{\varphi }\left(\frac{+𝒯}{2\omega +3}\right)F(r)e^{2(KU)},$$
(32)
and from (24) and (27)
$$K^{\prime \prime }2U^{\prime \prime }+(K^{}2U^{})\left(\frac{W^{}}{W}+\frac{\varphi ^{}}{\varphi }\right)=8\pi \left(\frac{𝒯}{\varphi }\right)F(r)e^{2(KU)}.$$
(33)
Now, using the conservation equation (16), equation (33) can be put as
$$\left(\frac{+𝒯}{𝒯}\right)\left[K^{\prime \prime }+K^{}\left(\frac{W^{}}{W}+\frac{\varphi ^{}}{\varphi }\right)\right]=8\pi \left(\frac{𝒯}{\varphi }\right)F(r)e^{2(KU)}.$$
(34)
Then, comparing this with equation (32) we find that it should be
$$\left(\frac{+𝒯}{𝒯}\right)^2=2\omega +3.$$
(35)
Hence, for $`𝒯`$ a static solution would be possible only if
$$\omega <\frac{3}{2},$$
(36)
which corresponds (see reference 14) to $`G<0`$, that is, to a theory in which gravitation is repulsive. Hence a static solution for the interior metric cannot exist in Brans-Dicke full theory.
4. DISCUSSION
Hiscock showed that the deficit angle obtained by Vilenkin within the linear approximation of general relativity is correct to all orders in $`G\mu `$. For this, he used a thick cylinder as a source and found an exact interior static solution which he matched with the exterior metric. Here we have shown that this procedure cannot be carried out with a more general tensor with $`T_t^tT_z^z`$: in this case there is no static interior solution.
Tensors of this kind were considered by Vachaspati and Vilenkin when they studied the effect of wiggles propagating along a string. They used $`\stackrel{~}{T}_\mu ^\nu =\text{diag}(,0,0,𝒯)\delta (x)\delta (y)`$ for calculating the exterior metric within linearized Einstein gravity. The energy per unit length $``$ and the effective tension $`𝒯`$ (fulfilling $`𝒯=\mu ^2`$) are obtained by averaging over a distance and a time much greater than the typical wavelength and oscillating period of the wiggles. It may be thought, following Hiscock’s idea, that a natural extension would be to use $`T_\mu ^\nu =\text{diag}(,0,0,𝒯)F(r)`$ and to find a stationary interior solution; this solution would be the time average of the actual time-dependent metric (this implies neglecting gravitational radiation). In this picture, $`F(r)`$ would play the role of a spatial distribution obtained as a time average of the radial position of the string; the interior of the source would be the region defined by $`r`$ less than the maximum amplitude of the wiggles. Regardless of the validity of this approximation, it is clear from our analysis that the inmediate extension valid for local straight strings cannot be applied in the case of wiggly strings. To obtain a solution valid to all orders it may be necessary either to improve the approximations made in the energy-momentum tensor or to consider the possibility of a time-dependent solution.
ACKNOWLEDGMENT
We wish to thank F. D. Mazzitelli for reading the manuscript and making helpful comments.
REFERENCES
1. A. Vilenkin and E. P. S. Shellard, Cosmic Strings and Other Topological Deffects, Cambridge University Press, Cambridge (1994).
2. A. Vilenkin, Phys. Rep. 121, 263 (1985).
3. A. Vilenkin, Phys. Rev. D 23, 852 (1981).
4. W. A. Hiscock, Phys. Rev. D 31, 3288 (1985).
5. T. Vachaspati and A. Vilenkin, Phys. Rev. Lett. 67, 1057 (1991).
6. T. Vachaspati, Phys. Rev. D 45, 3487 (1992).
7. A. Arazi and C. Simeone, submitted to Phys. Rev. D.
8. B. Linet, Gen. Rel. Grav. 17, 1109 (1985).
9. K. S. Thorne, Phys. Rev. 138, 251 (1965).
10. L. D. Landau and E. M. Lifshitz, The Classical Theory of Fields, Pergamon Press, Oxford (1975).
11. A. A. Sen, N. Banerjee and A. Banerjee, Phys. Rev. D 56, 3706 (1997).
12. C. Gundlach and M. Ortiz, Phys. Rev. D 42, 2521 (1990).
13. B. Boisseau and B. Linet, Gen. Rel. Grav. 30, 963 (1998).
14. S. Weinberg, Gravitation and Cosmology, John Wiley and sons, New York (1972).
15. C. W. Misner, K. S. Thorne and J. Wheeler, Gravitation, W. H. Freeman and company, New York (1997).
16. A. Barros and C. Romero, J. Math. Phys. 36, 5800 (1995).
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# Valleyless Sequences
## 1 Introduction
Throughout this article we use the notation $`𝒮_n`$ to denote finite sequences of length $`n`$. A length $`n`$ sequence $`s=s_1,s_2,\mathrm{},s_n`$ of positive integers is a permutation whenever each $`s_i`$ is a distinct member of the set $`\{1,2,\mathrm{},n\}`$. In this case we write the sequence $`s`$ as a word $`\pi =\pi _1\pi _2\mathrm{}\pi _n`$ and denote the set of all permutations of length $`n`$ as $`𝒫_n`$ . In the general case, the sequence of positive integers $`s`$ can be associated with a composition of the integer $`s_1+s_2+\mathrm{}+s_n`$. The composition has $`n`$ parts where $`s_i`$ is the size of the $`i^{th}`$ part. We let $`𝒞_n`$ be the set of all compositions of the positive integer $`n`$.
A sequence is defined herein to be *valleyless* provided $`1i<j<kn`$ implies $`s_j\mathrm{min}\{s_i,s_k\}`$. A graph of the pairs $`(i,s_i)`$ reveals that a valleyless sequence never has a valley. As a sequence, it is either nondecreasing, nonincreasing, or is nondecreasing to a point and thereafter nonincreasing. We let $`𝒱_n`$ be the set of all length $`n`$ valleyless sequences of positive integers and $`𝒱_{n,k}`$ the set of all length $`n`$ valleyless sequences of positive integers with maximum part $`k`$.
Valleyless sequences occur in tree enumeration problems and provide an interesting correspondence between permutations and compositions. We explore this correspondence first and then enumerate valleyless sequences in three different ways.
## 2 Valleyless permutations
A permutation $`\pi =\pi _1\pi _2\mathrm{}\pi _n𝒫_n`$ is uniquely determined by its inversion table $`I(\pi )=(a_1,a_2,\mathrm{},a_n)`$. The entry $`a_k`$ is the number of symbols in the word $`\pi `$ to the left of $`k`$ that are greater than $`k`$. Thus, for all $`k`$, $`0a_knk`$. The total number of inversions of $`\pi `$ is $`i(\pi )=a_1+a_2+\mathrm{}+a_n`$ and
$$\underset{\pi 𝒮_n}{}q^{i(\pi )}=(1+q)(1+q+q^2)\mathrm{}(1+q+q^2+\mathrm{}+q^{n1}).$$
As an example we have $`\pi =2731546`$, $`a_1=3`$, $`I(\pi )=(3,0,1,2,1,1,0)`$, and $`i(\pi )=8`$.
A finite sequence $`s=s_1,s_2,\mathrm{},s_n`$ viewed as a composition $`s_1+s_2+\mathrm{}+s_n=N`$ can be encoded as a unique subset of $`\{1,2,\mathrm{},N1\}`$ by $`\mathrm{\Theta }(s)=\{s_1,s_1+s_2,\mathrm{},s_1+s_2+\mathrm{}+s_{n1}\}`$. For example, the sequence $`s=2,3,2,1,1`$ is the composition $`2+3+2+1+1=9`$ with $`\mathrm{\Theta }(s)=\{2,5,7,8\}`$. The process is reversible so that there are $`2^{N1}`$ different compositions of $`N`$.
Next, we present a result relating *valleyless* permutations and compositions.
###### Theorem 1.
The subset of valleyless permutations in $`𝒫_n`$ is in one-to-one correspondence with the compositions of $`n`$. That is, $`|𝒱_n𝒫_n|=|𝒞_n|`$.
###### Proof.
Let $`\pi =\pi _1\pi _2\mathrm{}\pi _n`$ be valleyless, then either $`\pi _1=1`$ or $`\pi _n=1`$. Add one to every entry of $`\pi `$ and then append 1 to the beginning or end. The new permutation is valleyless of length $`n+1`$. The process is reversible so there are twice as many valleyless permutations of length $`n+1`$ as there are of length $`n`$. Since $`|𝒫_1𝒮_1|=1`$ we have $`|𝒱_n𝒫_n|=2^{n1}=|𝒞_n|`$ and we are done. ∎
The statement of the theorem implies a stronger correspondence than the cardinality result. There is indeed such a correspondence which we present now. Our goal in presenting the above proof was to introduce the reader to the technique of altering a sequence by adding one to every entry and then appending a 1 at the beginning or end. We will use this technique in the future. The stronger result is obtained from the following.
###### Theorem 2.
A permutation $`\pi =\pi _1\pi _2\mathrm{}\pi _n`$ with inversion table $`I(\pi )=(a_1,a_2,\mathrm{},a_n)`$ is valleyless if and only if $`a_k=0`$ or $`nk`$, $`k=1,2,\mathrm{},n`$.
###### Proof.
Suppose that a permutation $`\pi =\pi _1\pi _2\mathrm{}\pi _n`$ with inversion table $`(a_1,a_2,\mathrm{},a_n)`$ is valleyless. If $`a_k>0`$ for some $`k`$, $`1kn`$, then there is at least one symbol in $`\pi `$ to the left of $`k`$ that is greater than $`k`$. If there is any symbol in $`\pi `$ greater than $`k`$ to the right of $`k`$, then $`k`$ would be a valley contradicting our assumption. Hence all numbers greater than $`k`$ must be to the left of $`k`$ and $`a_k=nk`$. Conversely, suppose that $`I(\pi )=(a_1,a_2,\mathrm{},a_n)`$ and $`a_k=0`$ or $`nk`$, $`k=1,2,\mathrm{},n`$. We want to show that $`\pi `$ belongs to $`𝒱_n𝒫_n`$. Suppose not, then there is at least one symbol $`\pi _k`$ in $`\pi `$ such that
$$\pi _k<\pi _i,\text{ for some }i<k\text{ and }\pi _k<\pi _j,\text{ for some }j>k.$$
Then $`a_k0`$ and $`a_k<nk`$ contradicting our assumption. Hence $`\pi `$ belongs to $`𝒱_n𝒫_n`$ and we are done. ∎
Note that there are two choices for each $`a_k`$ except for $`k=n`$ in which case there is only one. Thus, there are $`2^{n1}`$ valleyless permutations of length $`n`$. If we again use $`q`$ to record the number of inversions in a permutation, then we have
###### Corollary 3.
$$\underset{\pi 𝒱_n𝒫_n}{}q^{i(\pi )}=(1+q)(1+q^2)\mathrm{}(1+q^{n1}).$$
## 3 Permutations with a valley
One of the fundamental statistics associated with a permutation $`\pi =\pi _1\mathrm{}\pi _n`$ is the descent set $`D(\pi )=\{i|\pi _i>\pi _{i+1}\}.`$ See for the enumeration of permutations by their descent set. We say that a length $`n`$ sequence $`s=s_1,s_2,\mathrm{},s_n`$ has exactly $`k`$ valleys if
$$|\{j:s_j<min\{s_{j1},s_{j+1}\}\}|=k.$$
In this section we obtain a recursive formula for the generating function of permutations with exactly $`k`$ valleys which we denote by $`g_k(x)`$ and use the formula to reproduce the table of permutations with peaks in .
We have shown in the previous section that $`g_0(x)`$ satisfies the relation
$`g_0(x)={\displaystyle \frac{x}{12x}}.`$ (1)
One can easily see that length $`n+1`$ permutations with exactly $`k+1`$ valleys can be obtained from:
1. length $`n`$ permutations with exactly $`k+1`$ valleys by adding $`1`$ to every entry and then appending a $`1`$ either at the beginning, end, to the left or to the right of the $`k+1`$ valleys.
2. length $`n`$ permutations with exactly $`k`$ valleys by adding $`1`$ to every entry and then inserting a $`1`$ in one of the $`(n1)2k=n(2k+1)`$ middle positions.
Hence
$`g_{k+1}(x)=(2(k+1)+2)xg_{k+1}(x)+x^{2k+3}D_x\left({\displaystyle \frac{g_k(x)}{x^{2k+1}}}\right),\text{ for }k0.`$ (2)
Substituting $`f_k(x)=\frac{g_k(x)}{x^{2k+1}}`$ in $`(2)`$ we obtain a system of linear differential equations
$`f_{k+1}(x)={\displaystyle \frac{1}{1(2(k+1)+2)x}}D_x(f_k(x))`$ (3)
which can also be written in matrix form as
$`\left[\begin{array}{c}f_0\\ f_1\\ f_2\\ f_3\\ f_4\\ \mathrm{}\end{array}\right]=[{\displaystyle \frac{1}{12x}},{\displaystyle \frac{1}{14x}},{\displaystyle \frac{1}{16x}},{\displaystyle \frac{1}{18x}},\mathrm{}]D_x\left[\begin{array}{c}x\\ f_0\\ f_1\\ f_2\\ f_3\\ \mathrm{}\end{array}\right]`$
Even though $`(3)`$ is an elegant recurrence relation just like many known simple relations for the number of partitions of an integer, we do not know how to solve it at this time.
| k | $`g_k(x)`$ |
| --- | --- |
| 0 | $`\frac{x}{12x}`$ |
| 1 | $`\frac{2x^3}{(12x)^2(14x)}`$ |
| 2 | $`\frac{16x^5(13x)}{(12x)^3(14x)^2(16x)}`$ |
| 3 | $`\frac{16x^7(17184x+636x^2720x^3)}{(12x)^4(14x)^3(16x)^2(18x)}`$ |
| 4 | $`\frac{256x^9(31788x+8096x^243132x^3+126072x^4192672x^5+120960x^6)}{(12x)^5(14x)^4(16x)^3(18x)^2(110x)}`$ |
Table 1: Generating functions of permutations with valleys
The formal Taylor series expansion of the functions in in Table 1 reproduces the table of permutations with peaks in .
| k/n | 1 | 2 | 3 | 4 | 5 | 6 | 7 | 8 | 9 | 10 |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| 0 | 1 | 2 | 4 | 8 | 16 | 32 | 64 | 128 | 256 | 512 |
| 1 | | | 2 | 16 | 88 | 416 | 1824 | 7680 | 31616 | 128512 |
| 2 | | | | | 16 | 272 | 2880 | 24576 | 185856 | 1304832 |
| 3 | | | | | | | 272 | 7936 | 137216 | 1841152 |
| 4 | | | | | | | | | 7936 | 353792 |
Table 2: The number of permutations of length $`n`$ with exactly $`k`$ valleys
If we denote the number of permutations of $`n`$ numbers with $`k`$ valleys by $`P(n,k)`$, the functional recursion (2) can also be written as
$`P(n,k)=2(k+1)P(n1,k)+(n2k)P(n1,k1)`$ (5)
and notice the remarkable similarity of the recurrence relation (4) with that of Eulerian numbers
$$E(n,k)=(k+1)E(n1,k)+(nk)E(n1,k1).$$
## 4 Valleyless sequences
In this section we enumerate valleyless sequences of length $`n`$ and maximum entry $`k`$ which we have denoted by $`𝒱_{n,k}`$ using the method of generating functions. Table 3 shows the cardinalities of $`𝒱_{n,k}`$ for $`1n6`$ and $`1k5`$.
| n/k | 1 | 2 | 3 | 4 | 5 |
| --- | --- | --- | --- | --- | --- |
| 1 | 1 | 1 | 1 | 1 | 1 |
| 2 | 1 | 3 | 5 | 7 | 9 |
| 3 | 1 | 6 | 15 | 28 | 45 |
| 4 | 1 | 10 | 35 | 84 | 165 |
| 5 | 1 | 15 | 70 | 210 | 495 |
| 6 | 1 | 21 | 126 | 462 | 1287 |
Table 3: Valleyless Sequences of Length $`n`$ with maximum entry $`k`$
Let $`V(x,y)`$ be a generating function which enumerates valleyless sequences of length $`n`$ and maximum entry $`k`$. Then
$$V(x,y)=\underset{n,k}{}V_{n,k}x^ny^k=1+xy+xy^2+\mathrm{},$$
where $`V_{n,k}=|𝒱_{n,k}|`$. As pointed out in the proof of Theorem 1, any valleyless sequence of length $`n`$ and maximum entry $`k`$ can be obtained from the base case:
$$1,11,111,1111,\mathrm{}$$
by adding one to every entry and then appending a $`1`$ at the beginning or end. Hence $`V(x,y)`$ satisfies the recursion
$`V(x,y)`$ $`=`$ $`\underset{\text{base case}}{\underset{}{\{x+x^2y+x^3y+\mathrm{}\}}}+\stackrel{\text{adding 1}}{\stackrel{}{\{yV(x,y)\}}}\times \underset{\text{appending a 1}}{\underset{}{\{1+2x+3x^2+\mathrm{}\}}}`$
$`=`$ $`{\displaystyle \frac{xy}{1x}}+yV(x,y){\displaystyle \frac{1}{(1x)^2}}.`$
Therefore,
$`V(x,y)\left(1{\displaystyle \frac{y}{(1x)^2}}\right)={\displaystyle \frac{xy}{1x}}\text{ and }V(x,y)={\displaystyle \frac{xy}{1x\frac{y}{1x}}}.`$ (6)
On the other hand, it is well known that
$$\underset{k=0}{\overset{\mathrm{}}{}}\left(\genfrac{}{}{0pt}{}{n1+k}{k}\right)y^k=\frac{1}{(1y)^n}.$$
Using this closed form and the ’snake oil’ method, we observe that
$`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n1+2(k1)}{2(k1)}}\right)y^k\right)x^n`$
$`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{y}{2}}\left({\displaystyle \frac{1}{(1\sqrt{y})^n}}+{\displaystyle \frac{1}{(1+\sqrt{y})^n}}\right)x^n`$
$`=`$ $`{\displaystyle \frac{y}{2}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{x}{1\sqrt{y}}}\right)^n+{\displaystyle \frac{y}{2}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{x}{1+\sqrt{y}}}\right)^n`$
$`=`$ $`{\displaystyle \frac{y}{2}}\left({\displaystyle \frac{x}{1\sqrt{y}x}}+{\displaystyle \frac{x}{1+\sqrt{y}x}}\right)`$
$`=`$ $`{\displaystyle \frac{xy}{(1x)\frac{y}{1x}}}.`$ (8)
Hence from $`(5)`$ and $`(6)`$ one can see that
###### Theorem 4.
The number of valleyless sequences of length $`n`$ with maximum entry $`k`$ is
$$V_{n,k}=\left(\genfrac{}{}{0pt}{}{n1+2(k1)}{2(k1)}\right).$$
## 5 q-Analog
If we use $`q`$ to record the sum of the entries of members of $`𝒱_{n,k}`$ and let $`V(x,q,y)`$ its generating function, then
$$V(x,q,y)=\underset{n,p,k}{}V_{n,p,k}x^nq^py^k.$$
Using the same argument as in the previous case, we see that $`V(x,q,y)`$ satisfies the functional recursion
$`V(x,q,y)`$ $`=`$ $`{\displaystyle \frac{xqy}{1xq}}+yV(xq,q,y){\displaystyle \frac{1}{(1xq)^2}}`$
$`(1xq)^2V(x,q,y)`$ $`=`$ $`xqyx^2q^2y+yV(xq,q,y).`$ (9)
Equation $`(7)`$ is a linear $`q`$difference equation and we seek series solution of the form
$`V(x,q,y)={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}b_n(x,q)y^n.`$ (10)
Substituting $`(8)`$ into $`(7)`$ and comparing coefficients of $`y^n`$ we obtain:
$`b_n(x,q)={\displaystyle \frac{b_{n1}(xq,q)}{(1xq)^2}}\text{ for }n>1,`$ (11)
where $`b_0=0\text{ and }b_1=\frac{xq}{1xq}.`$
Repeated application of equation (9) gives the explicit solution
$`b_n={\displaystyle \frac{xq^n(1xq^n)}{(xq)_n^2}},`$ (12)
where the common $`q`$symbol
$$(xq)_n=\underset{i=1}{\overset{n}{}}(1xq^i).$$
Therefore,
$`V(x,q,y)`$ $`=`$ $`{\displaystyle \frac{xq}{1xq}}y+{\displaystyle \underset{n=2}{\overset{\mathrm{}}{}}}{\displaystyle \frac{xq^n(1xq^n)}{(xq)_n^2}}y^n`$ (13)
$`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{xq^n(1xq^n)}{(xq)_n^2}}y^n`$
An implementation of $`(11)`$ using Maple<sup>1</sup><sup>1</sup>1Send e-mail to zelekem@wpunj.edu to obtain the Maple Package bxq accompanying this paper. shows that the number of valleyless sequences of length $`10`$, sum of entries $`20`$, and maximum part $`5`$, for instance, is 325.
## 6 A nonlinear relation
Expanding $`V(x,q,y)`$ in $`x`$ as $`V(x,q,y)=a_n(q,y)x^n`$, substituting it into $`(5)`$, and comparing coefficients of $`x^n`$, we obtain a three term recurrence relation
$`a_n={\displaystyle \frac{2qa_{n1}q^2a_{n2}}{1yq^n}}\text{ for }n>2`$ (14)
where $`a_0=0,a_1=\frac{yq}{1yq},\text{ and }a_2=\frac{q^2y+q^3y^2}{(yq;q)_2}`$.
This nonlinear recurrence relation provides yet another way of enumerating valleyless sequences. Indeed, one can obtain valleyless sequences of length $`n`$ from sequences of length $`n1`$ by adding a 1 at the beginning or end, and then appending 1s to every entry. This process contributes $`\frac{2q}{1yq^n}a_{n1}(q,y)`$ to $`a_n(q,y)`$. However, sequences with the same first and last entry, for example $`23522`$, are counted twice in this algorithm and we need to subtract $`\frac{q^2}{1yq^n}a_{n2}(q,y)`$, the number of valleyless sequences of length $`n`$ with the same first and last entry. Hence,
$$a_n(q,y)=\frac{2q}{1yq^n}a_{n1}(q,y)\frac{q^2}{1yq^n}a_{n2}(q,y).$$
Unlike (9) equation (12) is nonlinear and we do not know its explicit solution.
## 7 Conclusion
We came across valleyless sequences while enumerating ordered trees by their number of leaves, total path length and number of vertices. Given a positive integer $`n`$, consider its composition or ordered partition and draw a composition tree corresponding to $`n`$. Then try to construct all possible ordered trees of total path length $`n`$ by identifying non-terminal nodes or vertices of the composition tree.
A sequence obtained from an ordered partition of a positive integer $`n`$ by taking the minimum of two consecutive entries gives a sequence without isolated valley and we hope that the ideas developed in this article will provide an alternative means to solve the ordered tree enumeration problem.
## 8 Acknowledgements
The second author acknowledges research release time from William Paterson University and summer support from The Center for Research of the College of Science and Health.
jrieper@cybernex.net
zelekem@wpunj.edu
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# 𝐷^∗𝐷𝜋 and 𝐵^∗𝐵𝜋 form factors from QCD Sum Rules
## Abstract
The $`H^{}H\pi `$ form factor for $`H=B`$ and $`D`$ mesons is evaluated in a QCD sum rule calculation. We study the Borel sum rule for the three point function of two pseudoscalar and one vector meson currents up to order four in the operator product expansion. The double Borel transform is performed with respect to the heavy meson momenta. We discuss the momentum dependence of the form factors and two different approaches to extract the $`H^{}H\pi `$ coupling constant.
PACS numbers 14.40.Lb, 14.40.Nd, 12.38.Lg, 11.55.Hx
The coupling of the pion to the heavy mesons ($`g_{B^{}B\pi }`$ and $`g_{D^{}D\pi }`$) is related to the form factor at zero pionic momentum and its precise value has been often needed in phenomenology. In particular, the $`g_{D^{}D\pi }`$ coupling is needed in the context of quark gluon plasma (QGP) physics. Suppression of charmonium production in heavy ion collisions is one of the signatures of QGP formation . Therefore a precise evaluation of the background, i.e., conventional $`J/\psi `$ absorption by co-moving pions and $`\rho `$ mesons , is of fundamental importance. Since pions are so abundant in a dense nuclear environment, the reactions $`\pi +J/\psi D+\overline{D^{}}`$ (and consequently the coupling $`g_{D^{}D\pi }`$) are of special relevance .
In the case of $`g_{D^{}D\pi }`$, the $`D_{}^{}{}_{}{}^{+}D^0\pi ^+`$ decay is observed experimentally. However, present data provide only an upper bound: $`g_{D^{}D\pi }21`$ . For $`g_{B^{}B\pi }`$, there cannot be a direct experimental indication because there is no phase space for the $`B^{}B\pi `$ decay. Recently, a direct preliminary determination of $`g_{B^{}B\pi }`$ on the lattice has been attempted .
The $`D^{}D\pi `$ and $`B^{}B\pi `$ couplings have been studied by several authors using different approaches of the QCD sum rules (QCDSR): two point function combined with soft pion techniques , light cone sum rules , light cone sum rules including perturbative corrections , sum rules in a external field , double momentum sum rules . Unfortunately, the numerical results from these calculations may differ by almost a factor two.
In this work we use the three-point function approach to evaluate the $`D^{}D\pi `$ and $`B^{}B\pi `$ form factors and coupling constants. The advantage of using the three-point function approach with a double Borel transformation compared with the two-point function with a single Borel transformation is the elimination of the terms associated with the pole-continuum transitions .
The three-point function associated with a $`H^{}H\pi `$ vertex, where $`H`$ and $`H^{}`$ are respectively the lowest pseudoscalar and vector heavy mesons, is given by
$$\mathrm{\Gamma }_\mu (p,p^{})=d^4xd^4y0|T\{j(x)j_5(y)j_\mu ^{}(0)\}|0e^{ip^{}.x}e^{i(p^{}p).y},$$
(1)
where $`j=i\overline{Q}\gamma _5u`$, $`j_5=i\overline{u}\gamma _5d`$ and $`j_\mu ^{}=\overline{d}\gamma _\mu Q`$ are the interpolating fields for $`H`$, $`\pi ^{}`$ and $`H^{}`$ respectively with $`u`$, $`d`$ and $`Q`$ being the up, down, and heavy quark fields.
The phenomenological side of the vertex function, $`\mathrm{\Gamma }_\mu (p,p^{})`$, is obtained by the consideration of $`H`$ and $`H^{}`$ state contribution to the matrix element in Eq. (1):
$`\mathrm{\Gamma }_\mu ^{(phen)}(p,p^{})`$ $`=`$ $`{\displaystyle \frac{1}{p^2m_H^{}^2}}{\displaystyle \frac{1}{p_{}^{}{}_{}{}^{2}m_H^2}}0|j|H(p^{})\times `$ (3)
$`H(p^{})|j_5|H^{}(p,ϵ)H^{}(p,ϵ)|j_\mu ^{}|0+\text{higher resonances}.`$
The matrix element of the pseudoscalar element, $`j_5`$, defines the vertex form factor $`g_{H^{}H\pi }(q^2)`$:
$$H(p^{})|j_5|H^{}(p,ϵ)=\frac{f_\pi m_\pi ^2}{m_u+m_d}\frac{g_{H^{}H\pi }(q^2)}{q^2m_\pi ^2}q_\nu ϵ^\nu ,$$
(4)
where $`q=p^{}p`$, $`f_\pi `$ is the pion decay constant and $`ϵ^\nu `$ is the polarization of the vector meson. The vacuum to meson transition amplitudes appearing in Eq. (3) are given in terms of the corresponding meson decay constants $`f_H`$ and $`f_H^{}`$ by
$$0|j|H(p^{})=\frac{m_H^2f_H}{m_Q},$$
(5)
and
$$H^{}(p,ϵ)|j_\mu ^{}|0=m_H^{}f_H^{}ϵ_\mu ^{}.$$
(6)
Therefore, using Eqs. (4), (5) and (6) in Eq. (3) we get
$`\mathrm{\Gamma }_\mu ^{(phen)}(p,p^{})`$ $`=`$ $`C_{HH^{}}{\displaystyle \frac{g_{H^{}H\pi }(q^2)}{q^2m_\pi ^2}}{\displaystyle \frac{1}{p^2m_H^{}^2}}{\displaystyle \frac{1}{p_{}^{}{}_{}{}^{2}m_H^2}}\times `$ (8)
$`\left(p_\mu ^{}+{\displaystyle \frac{m_H^{}^2+m_H^2q^2}{2m_H^{}^2}}p_\mu \right)+\text{higher resonances},`$
where
$$C_{HH^{}}=\frac{m_H^2m_H^{}m_\pi ^2f_Hf_H^{}f_\pi }{(m_u+m_d)m_Q}.$$
(9)
The contribution of higher resonances and continuum in Eq. (8) will be taken into account as usual in the standard form of ref. .
The QCD side, or theoretical side, of the vertex function is evaluated by performing Wilson’s operator product expansion (OPE) of the operator in Eq. (1). Writing $`\mathrm{\Gamma }_\mu `$ in terms of the invariant amplitudes:
$$\mathrm{\Gamma }_\mu (p,p^{})=\mathrm{\Gamma }_1(p^2,p_{}^{}{}_{}{}^{2},q^2)p_\mu +\mathrm{\Gamma }_2(p^2,p_{}^{}{}_{}{}^{2},q^2)p_\mu ^{},$$
(10)
we can write a double dispersion relation for each one of the invariant amplitudes $`\mathrm{\Gamma }_i(i=1,2)`$, over the virtualities $`p^2`$ and $`p_{}^{}{}_{}{}^{2}`$ holding $`Q^2=q^2`$ fixed:
$$\mathrm{\Gamma }_i(p^2,p_{}^{}{}_{}{}^{2},Q^2)=\frac{1}{4\pi ^2}_{m_Q^2}^{\mathrm{}}𝑑s_{m_Q^2}^{\mathrm{}}𝑑u\frac{\rho _i(s,u,Q^2)}{(sp^2)(up_{}^{}{}_{}{}^{2})},$$
(11)
where $`\rho _i(s,u,Q^2)`$ equals the double discontinuity of the amplitude $`\mathrm{\Gamma }_i(p^2,p_{}^{}{}_{}{}^{2},Q^2)`$ on the cuts $`m_Q^2s\mathrm{}`$, $`m_Q^2u\mathrm{}`$, which can be evaluated using Cutkosky’s rules .
Finally we perform a double Borel transformation in both variables $`P^2=p^2`$ and $`P_{}^{}{}_{}{}^{2}=p_{}^{}{}_{}{}^{2}`$ and equate the two representations described above. We get one sum rule for each invariant function. In the $`p_\mu `$ structure:
$`C_{HH^{}}{\displaystyle \frac{m_H^{}^2+m_H^2+Q^2}{2m_H^{}^2}}{\displaystyle \frac{g_{H^{}H\pi }(q^2)}{Q^2+m_\pi ^2}}e^{m_H^{}^2/M^2}e^{m_H^2/M_{}^{}{}_{}{}^{2}}`$ $`=`$ $`{\displaystyle \frac{1}{4\pi ^2}}{\displaystyle _{m_Q^2}^{s_0}}ds{\displaystyle _{m_Q^2}^{u_0}}du[`$ (13)
$`\rho _1(s,u,Q^2)e^{s/M^2}e^{u/M_{}^{}{}_{}{}^{2}}],`$
and in the $`p_\mu ^{}`$ structure:
$$C_{HH^{}}\frac{g_{H^{}H\pi }(q^2)}{Q^2+m_\pi ^2}e^{m_H^{}^2/M^2}e^{m_H^2/M_{}^{}{}_{}{}^{2}}=\frac{1}{4\pi ^2}_{m_Q^2}^{s_0}𝑑s_{m_Q^2}^{u_0}𝑑u\left[\rho _2(s,u,Q^2)e^{s/M^2}e^{u/M_{}^{}{}_{}{}^{2}}\right],$$
(14)
where $`s_0`$ and $`u_0`$ are the continuum thresholds for the $`H^{}`$ and $`H`$ mesons respectively, which are, in general, taken from the mass sum rules. The two Borel masses $`M^2`$ and $`M_{}^{}{}_{}{}^{2}`$ are, in principle, independent and they should vary in the vicinity of the corresponding meson masses: $`m_H^{}^2`$ and $`m_H^2`$ respectively. Since for heavy mesons $`m_H`$ and $`m_H^{}`$ are very close, many authors use $`M^2=M_{}^{}{}_{}{}^{2}`$ . To allow for different values of $`M^2`$ and $`M_{}^{}{}_{}{}^{2}`$ we take them proportional to the respective meson masses, which leads us to study the sum rule as a function of $`M^2`$ at a fixed ratio
$$\frac{M^2}{M_{}^{}{}_{}{}^{2}}=\frac{m_H^{}^2}{m_H^2}.$$
(15)
We will consider diagrams up to dimension four which include the perturbative diagram and the gluon condensate. The quark condensate term does not contribute since it depends only on one external momentum and, therefore, it is eliminated by the double Borel transformation. Higher dimension condensates are strongly suppressed in the case of heavy quarks . The double discontinuity of the perturbative contribution reads:
$$\rho _1^{(pert)}(s,u,Q^2)=\frac{3Q^2u(2m_Q^2suQ^2)}{[(s+u+Q^2)^24su]^{3/2}},$$
(16)
$$\rho _2^{(pert)}(s,u,Q^2)=\frac{3Q^2[m_Q^2(s+u+Q^2)2su]}{[(s+u+Q^2)^24su]^{3/2}},$$
(17)
and the integration limit condition is
$$(sm_Q^2)(um_Q^2)Q^2m_Q^2.$$
(18)
In this paper we focus on the structure $`p_\mu `$ which we found to be the more stable one. For consistency we use in our analysis the QCDSR expressions for the decay constants up to dimension four in lowest order of $`\alpha _s`$:
$$f_H^2=\frac{3m_Q^2}{8\pi ^2m_H^4}_{m_Q^2}^{u_0}𝑑u\frac{(um_Q^2)^2}{u}e^{(m_H^2u)/M_{}^{}{}_{}{}^{2}}\frac{m_Q^3}{m_H^4}\overline{q}qe^{(m_H^2m_Q^2)/M_{}^{}{}_{}{}^{2}},$$
(19)
$$f_H^{}^2=\frac{1}{8\pi ^2m_H^{}^2}_{m_Q^2}^{s_0}𝑑s\frac{(sm_Q^2)^2}{s}\left(2+\frac{m_Q^2}{s}\right)e^{(m_H^{}^2s)/M^2}\frac{m_Q}{m_H^{}^2}\overline{q}qe^{(m_H^{}^2m_Q^2)/M^2},$$
(20)
where we have omitted the numerically insignificant contribution of the gluon condensate.
The parameter values used in all calculations are $`m_u+m_d=14\text{MeV}`$, $`m_c=1.5\text{GeV}`$, $`m_b=4.7\text{GeV}`$, $`m_\pi =140\text{MeV}`$, $`m_D=1.87\text{GeV}`$, $`m_D^{}=2.01\text{GeV}`$, $`m_B=5.28\text{GeV}`$, $`m_B^{}=5.33\text{GeV}`$, $`f_\pi =131.5\text{MeV}`$, $`\overline{q}q=(0.23)^3\text{GeV}^3`$, $`g^2G^2=0.5\text{GeV}^4`$. We parametrize the continuum thresholds as
$$s_0=(m_H^{}+\mathrm{\Delta }_s)^2,$$
(21)
and
$$u_0=(m_H+\mathrm{\Delta }_u)^2.$$
(22)
The values of $`u_0`$ and $`s_0`$ are, in general, extracted from the two-point function sum rules for $`f_H`$ and $`f_H`$ in Eqs. (19) and (20). Using the Borel region $`2M^25\text{GeV}^2`$ (for the $`D^{}`$ and $`D`$ mesons) and $`10M^225\text{GeV}^2`$ (for the $`B^{}`$, and $`B`$ mesons) we found a good stability for $`f_H`$ and $`f_H`$ with $`\mathrm{\Delta }_s=\mathrm{\Delta }_u0.5\text{GeV}`$, in agreement with the results in ref. . We have checked that bigger values for $`\mathrm{\Delta }_{s(u)}`$, of order of 1 GeV, lead to unstable results for $`f_H`$ and $`f_H`$, in the case of the sum rules Eqs. (19) and (20). In our study we will allow for a small variation in $`\mathrm{\Delta }_s`$ and $`\mathrm{\Delta }_u`$ to test the sensitivity of our results to the continuum contribution.
We first discuss the $`D^{}D\pi `$ form factor. In Fig. 1 we show the behavior of the perturbative and gluon condensate contributions to the form factor $`g_{D^{}D\pi }(Q^2)`$ at $`Q^2=1\text{GeV}`$ as a function of the Borel mass $`M^2`$ using $`\mathrm{\Delta }_s`$ and $`\mathrm{\Delta }_u`$ given in Eqs. (21) and (22) equal to $`0.5\text{GeV}`$. We can see that, in the case of the form factor, the gluon condensate is not negligible and it helps the stability of the curve, as a function of $`M^2`$, providing a rather stable plateau for $`M^23\text{GeV}^2`$. The behavior of the curve for other $`Q^2`$ and continuum treshold values is similar. Fixing $`M^2=3.5\text{GeV}^2`$ we show, in Fig. 2, the momentum dependence of the form factor (dots). Since the present approach cannot be used at $`Q^2=0`$, to extract the $`g_{D^{}D\pi }`$ coupling from the form factor we need to extrapolate the curve to $`Q^2=0`$ (in the approximation $`m_\pi ^2=0`$). In order to do this extrapolation we fit the QCD sum rule results (dots) with an analytical expression. We tried to fit our results with a monopole form, since this is very often used for form factors, but the fit is very poor. We obtained good fits using both the gaussian form
$$g_{H^{}H\pi }(Q^2)=g_{H^{}H\pi }e^{(Q^2+m_\pi ^2)^2/\mathrm{\Gamma }^4}$$
(23)
and a curve of the form
$$g_{H^{}H\pi }(Q^2)=g_{H^{}H\pi }\frac{1+(a/\mathrm{\Lambda })^4}{1+(a/\mathrm{\Lambda })^4e^{(Q^2+m_\pi ^2)^2/\mathrm{\Lambda }^4}}.$$
(24)
In Fig. 2 we show that the $`Q^2`$ dependence of the form factor, represented by the dots, can be well reproduced by the parametrization in Eqs. (23) (dashed line) and (24) (solid line). The value of the parameters in Eqs. (23) and (24) are given in Table I for two different values of the continuum threshold.
| $`\mathrm{\Delta }_s=\mathrm{\Delta }_u(\text{GeV})`$ | $`g_{D^{}D\pi }`$ | $`\mathrm{\Lambda }(\text{GeV})`$ | $`a(\text{GeV})`$ | $`\mathrm{\Gamma }(\text{GeV})`$ |
| --- | --- | --- | --- | --- |
| 0.5 | 5.3 | 1.66 | 1.90 | - |
| 0.6 | 6.0 | 1.89 | 3.05 | - |
| 0.5 | 5.7 | - | - | 1.74 |
| 0.6 | 6.1 | - | - | 1.92 |
TABLE I: Values of the parameters in Eqs. (23) and (24) which reproduce the QCDSR results for $`g_{D^{}D\pi }(Q^2)`$, for two different values of the continuum thresholds in Eqs. (21) and (22).
In view of the uncertainties involved, the results obtained with the two parametrizations are consistent with each other, the systematic error being of the order of $`10\%`$.
In refs. it was found that the form factor in the semileptonic decay $`H\pi l\overline{\nu }`$, which is also normalized by the $`H^{}H\pi `$ coupling constant, can be well approximated by a monopole form factor. In the case of the $`H\pi l\overline{\nu }`$ form factor, a vector dominance approximation gives a phenomenological explanation for a pole fit at $`q^2=m_H^{}^2`$, which is not the case of the form factor studied here. It is important to notice that here the dispersion relation is written in terms of the two heavy meson momenta, while in the case of semileptonic decay the dispersion relation is a function of the $`H`$ and $`\pi `$ momenta. Therefore, our form factor is a function of the pion momentum, exhibiting a peak at the pion pole $`Q^2=0`$.
To test if our fit gives a good extrapolation to $`Q^2=0`$ we can write a sum rule, based on the three-point function Eq. (1), but valid only at $`Q^2=0`$, as suggested in for the pion-nucleon coupling constant. This method was also applied to the nucleon-hyperon-kaon coupling constant and to the nucleon$`\mathrm{\Lambda }_cD`$ coupling constant . It consists in neglecting the pion mass in the denominator of Eq. (8) and working at $`Q^2=0`$, making a single Borel transformation to both $`P^2=P_{}^{}{}_{}{}^{2}M^2`$.
As discussed in the introduction, the problem of doing a single Borel transformation is the fact that the single pole contribution, associated with the $`NN^{}`$ transition, is not suppressed . In ref. it was explicitly shown that the pole-continuum transition has a different behavior as a function of the Borel mass as compared with the double pole contribution and continuum contribution: it grows with $`M^2`$ as compared with the double pole contribution. Therefore, the single pole contribution can be taken into account through the introduction of a parameter $`A`$, in the phenomenological side of the sum rule . Thus, neglecting $`m_\pi ^2`$ in the denominator of Eq. (8) and doing a single Borel transform in $`P^2=P_{}^{}{}_{}{}^{2}`$, we get for the structure $`p_\mu `$
$$\stackrel{~}{\mathrm{\Gamma }}_1^{(phen)}(M^2,Q^2)=\frac{C_{H^{}H}}{2m_H^2Q^2}\frac{m_H^2+m_H^{}^2+Q^2}{m_H^{}^2m_H^2}\left(e^{m_H^2/M^2}e^{m_H^{}^2/M^2}\right)(g_{H^{}H\pi }+AM^2),$$
(25)
where $`C_{H^{}H}`$ in given in Eq. (9) with $`f_H`$ and $`f_H^{}`$ given by Eqs. (19) and (20).
On the OPE side only terms proportional to $`1/Q^2`$ will contribute to the sum rule. Therefore, up to dimension four the only diagram that contributes is the quark condensate given by
$$\stackrel{~}{\mathrm{\Gamma }}_1^{<\overline{q}q>}(M^2,Q^2)=\frac{2m_Q\overline{q}q}{Q^2}e^{m_Q^2/M^2}.$$
(26)
Equating Eqs. (25) and (26) and taking $`Q^2=0`$ we obtain the sum rule for $`g_{H^{}H\pi }+AM^2`$, where $`A`$ denotes the contribution from the unknown single poles terms. It is interesting to point out that in the limit $`m_H^2+m_H^{}^2=2m_H^{}^2`$, the sum rule obtained in the $`p_\mu ^{}`$ structure coincides with the sum rule in the $`p_\mu `$ structure. In Fig. 3 we show, for $`\mathrm{\Delta }_s=\mathrm{\Delta }_u=0.5\text{GeV}`$, the QCDSR results for $`g_{D^{}D\pi }+AM^2`$ as a function of $`M^2`$ (dots) from where we see that, in the Borel region $`2M^25\text{GeV}^2`$, they follow a straight line. The value of the coupling constant is obtained by the extrapolation of the line to $`M^2=0`$ . Fitting the QCDSR results to a straight line we get
$$g_{D^{}D\pi }5.4,$$
(27)
in excellent agreement with the values obtained with the extrapolation of the form factor to $`Q^2=0`$, given in Table I.
It is reassuring that both methods, with completely different OPE sides and Borel transformation approaches, give the same value for the coupling constant.
In the case of $`B^{}B\pi `$ vertex, we show in Fig. 4, for $`\mathrm{\Delta }_s=\mathrm{\Delta }_u=0.5\text{GeV}`$, the $`Q^2=0`$ sum rule results for $`g_{B^{}B\pi }+AM^2`$ (dots) as a function of $`M^2`$. It also follows a straight line in the Borel region $`10M^225\text{GeV}^2`$, and the extrapolation to $`M^2=0`$ gives
$$g_{B^{}B\pi }10.6.$$
(28)
In Fig. 6 we show the QCDSR result for the perturbative and gluon condensate contributions to the form factor $`g_{B^{}B\pi }(Q^2)`$ at $`Q^2=2\text{GeV}^2`$ as a function of $`M^2`$ using $`\mathrm{\Delta }_s=\mathrm{\Delta }_u=0.5\text{GeV}`$. In this case the gluon condensate is very small but it still goes in the right direction of providing a stable plateau for $`M^215\text{GeV}^2`$. Fixing $`M^2=17\text{GeV}^2`$ we show, in Fig. 6, the $`Q^2`$ behavior of the form factor (dots). The dots can still be well fitted by Eq. (24) (solid line). However, the fit with Eq. (23) is not so good, as can be seen by the dashed line in Fig. 6. In Table II we give the value of the parameters in Eqs. (23) and (24) that reproduce our results for two different choices of the continuum thresholds.
In this case the agreement of the two different approaches to extract the coupling constant is not so good, but the numbers are still compatible. One possible reason for that is the fact that for heavier quarks the perturbative contribution (or hard physics) becomes more important, as can be observed by the decrease of the importance of the gluon condensate in Fig. 5 as compared with Fig. 1. Since in the sum rule given by Eqs. (25) and (26) there is only soft physics information, we expect $`\alpha _s`$ corrections to the sum rule to be more important in the case of $`g_{B^{}B\pi }(Q^2)`$ than for $`g_{D^{}D\pi }(Q^2)`$.
| $`\mathrm{\Delta }_s=\mathrm{\Delta }_u(\text{GeV})`$ | $`g_{B^{}B\pi }`$ | $`\mathrm{\Lambda }(\text{GeV})`$ | $`a(\text{GeV})`$ | $`\mathrm{\Gamma }(\text{GeV})`$ |
| --- | --- | --- | --- | --- |
| 0.5 | 14.7 | 1.62 | 1.37 | - |
| 0.6 | 16.3 | 1.81 | 1.67 | - |
| 0.5 | 17.2 | - | - | 1.79 |
| 0.6 | 18.4 | - | - | 1.97 |
TABLE II: Values of the parameters in Eqs. (23) and (24) which reproduce the QCDSR results for $`g_{B^{}B\pi }(Q^2)`$, for two different values of the continuum thresholds in Eqs. (21) and (22).
Comparing Table I with Table II we see that the cut-offs are of the same order in the two vertices and are very hard. Concerning the parameter $`a`$, it is smaller in the case of the $`B^{}B\pi `$ vertex. This is because of the fact that the form factor $`g_{B^{}B\pi }(Q^2)`$ has a flatter peak around $`Q^2=0`$ than $`g_{D^{}D\pi }(Q^2)`$. This can be interpreted as an indication that the spatial extension of the vertex is smaller for $`B^{}B\pi `$ than for $`D^{}D\pi `$. This is also the reason why the gaussian fit is not so good in the case of the $`B^{}B\pi `$ vertex, and leads to bigger values for the coupling. It is interesting to notice that our results for the coupling constants are completely consistent with the QCDSR calculation of ref. .
As a final exercise, we use our result for $`g_{B^{}B\pi }`$ to extract the coupling constant $`g`$ which controls the interaction of the pion with infinitely heavy fields in effective lagrangian approaches . They are related by
$$g_{B^{}B\pi }=\frac{2m_B}{f_\pi }g.$$
(29)
The knowledge of $`g`$ is of great phenomenological value, since its strenght is required in the analyzes of many electroweak processes . Therefore, during the last years, a large number of theoretical papers has been devoted to the calculation of $`g`$. However, the variation of the value obtained for $`g`$, even within a single class of models, turns out to be quite large. For instance, using different quark models one obtains $`1/3g1`$ while QCDSR calculations points in the direction of small $`g`$, with a typical value in the range $`g0.130.35`$ .
Using the values for $`g_{B^{}B\pi }`$ given in Table II we get, at order $`\alpha _s=0`$
$$g=0.130.23,$$
(30)
therefore, we corroborate the overall conclusion drawn from different QCDSR calculations, that the coupling $`g`$ is small.
In conclusion, we extracted the $`H^{}H\pi `$ coupling constant using two different approaches of the QCDSR based on the three-point function. We have obtained for the coupling constants:
$$g_{D^{}D\pi }=5.7\pm 0.4,$$
(31)
$$g_{B^{}B\pi }=14.5\pm 3.9,$$
(32)
where the errors reflect variations in the continuum thresholds, different parametrizations of the form factors and the use of two different sum rules. There are still sources of errors in the values of the condensates and in the choice of the Borel mass to extract the form factor, which were not considered here. Therefore, the errors quoted are probably underestimated.
In Table III we present a compilation of the estimates of the coupling constants $`g_{D^{}D\pi }`$ and $`g_{B^{}B\pi }`$ from distinct QCDSR calculations.
| approach | $`g_{D^{}D\pi }`$ | $`g_{B^{}B\pi }`$ |
| --- | --- | --- |
| this work | $`5.7\pm 0.4`$ | $`14.5\pm 3.9`$ |
| two-point function + soft pion techniques (2PFSP) | $`9\pm 2`$ | $`20\pm 4`$ |
| 2PFSP + perturbative corrections | $`7\pm 2`$ | $`15\pm 4`$ |
| light cone sum rules (LCSR) | $`11\pm 2`$ | $`28\pm 6`$ |
| LCSR + perturbative corrections | $`10.5\pm 3`$ | $`22\pm 9`$ |
| double momentum sum rule | $`6.3\pm 1.9`$ | $`14\pm 4`$ |
TABLE III: Summary of QCDSR estimates for $`g_{D^{}D\pi }`$ and $`g_{B^{}B\pi }`$.
From this Table we see that our result is in a fair agreement with the calculations in refs. , while LCSR calculations point to bigger values for the coupling constants. This discrepancy has still to be solved.
The $`D^{}D\pi `$ coupling is directly related with the $`D^{}D\pi `$ decay width through
$$\mathrm{\Gamma }(D_{}^{}{}_{}{}^{}\overline{D^0}\pi ^{})=\frac{g_{D^{}D\pi }^2|\stackrel{}{q}_\pi |^3}{24\pi m_D^{}^2}.$$
(33)
Using Eq. (31) we get
$$\mathrm{\Gamma }(D_{}^{}{}_{}{}^{}\overline{D^0}\pi ^{})=6.3\pm 0.9\text{keV},$$
(34)
which is much smaller then the current upper limit $`\mathrm{\Gamma }(D_{}^{}{}_{}{}^{}\overline{D^0}\pi ^{})<89`$ keV.
Acknowledgements: This work has been supported by CNPq and FAPESP (under project number 1998/2249-4). C.L.S. thanks FAPERJ for financial support.
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# Monopole Black Hole Skyrmions
## I Introduction
It is known that baryons can emerge as solitons in an effective meson field theory. In these models the baryon number is identified with a topological charge. Since this identification was originally made by T.H.R.Skyrme $`40`$ years ago , the soliton solution is called the skyrmion.
Skyrmion models provide a useful way of analysing the monopole catalysis of proton decay, as originally pointed out by Callan and Witten . The decay process can be allowed because, in the presence of a monopole, the baryon number is no longer equal to the topological charge of the meson field. In fact, there exist non-topological solitons with non-zero baryon number which can decay without topological problems.
Skyrmions have also been used to examine the absorbtion of baryons by microscopic black holes . Solutions representing skyrmions partially absorbed by Schwarzschild black holes provide a semiclassical framework to study the absorption rate of a proton by a black hole of comparable size. However, this process is rather insignificant because black holes of the size of a proton have large fluxes of Hawking radiation which swamp the proton decay.
In this paper we analyse black hole solutions with skyrmion hair and magnetic charge, providing the semiclassical framework in which to study monopole black hole catalysis of proton decay. The flux of Hawking radiation from these holes is less than ordinary holes and even vanishes in extremal cases . The black hole mass for abelian monopoles has the lower bound
$$M=\frac{p}{\sqrt{G}}=pm_{\mathrm{pl}}2.54\times 10^7Kg,$$
(1)
where $`\mathrm{e}`$ is a rationalised electric charge in gaussian units determined from the Dirac quantisation condition $`(\mathrm{e}p=1)`$ as $`p11.7`$ and $`m_{\mathrm{pl}}`$ is the planck mass which is approximately $`2.1768\times 10^8Kg`$. This kind of small black hole may be considered to have been created in the early universe and remain as a stable relic today. More complicated monopole solutions are also possible in Grand Unified Theories, depending on the details of the Higgs sector , but we will restrict attention to the simplest case.
In the extreme case, $`M=pm_{\mathrm{pl}}`$, we find the remarkable situation that multiple black hole solutions are possible in which the gravitational, electromagnetic, and strong forces between the monopoles are all in balance. In this respect the solitons behave in an analagous way to BPS monopole solutions in the Yang-Mills-Higgs system .
We find that the non-topological skyrmion solutions are stable within the confines of our model. This is, in part, due to the fact that we have not included the electron. The black hole cannot swallow the proton whole because this would tip it over the extremal limit. Light, charged particles are therefore needed to carry away the proton’s charge when it decays. We discuss a way of including this effect in the conclusion.
## II The lagrangian
We shall consider models with a charged $`SU(2)`$ meson field $`U`$. The Lagrangian is based on a gauged version of the original skyrmion Lagrangian . It should be regarded as a minimal form of the meson effective action, since extra terms could also be included. However, some of the terms in the Lagrangian which couple the chiral field to electromagnetism can be deduced from current algebra techniques and were constructed by Callan and Witten .
The Lagrangian can be divided into into four parts,
$$=_1+_2+_3+_4$$
(2)
where the meson parts are
$`_1`$ $`=`$ $`{\displaystyle \frac{F_\pi ^2}{16}}\mathrm{tr}\left(U^1D_\mu UU^1D_\mu U\right)+{\displaystyle \frac{1}{32a^2}}\mathrm{tr}\left([U^1D_\mu U,U^1D_\mu U]\right)`$
$`_2`$ $`=`$ $`{\displaystyle \frac{\mathrm{e}}{16\pi ^2}}ϵ^{\mu \nu \rho \sigma }A_\mu \mathrm{tr}\left\{Q\left(_\nu UU^1_\rho UU^1_\sigma UU^1+U^1_\nu UU^1_\rho UU^1_\sigma U\right)\right\}`$
$`+`$ $`{\displaystyle \frac{i\mathrm{e}^2}{8\pi ^2}}ϵ^{\mu \nu \rho \sigma }(_\mu A_\nu )A_\rho \mathrm{tr}\left(Q^2_\sigma UU^1+Q^2U^1_\sigma U+{\displaystyle \frac{1}{2}}Q_\sigma UQU^1{\displaystyle \frac{1}{2}}QUQ_\sigma U^1\right)`$
and the free actions are
$$_3=\frac{1}{16\pi }F_{\mu \nu }F^{\mu \nu },_4=\frac{1}{16\pi G}R.$$
(3)
The quantity $`Q`$ is the charge matrix of quarks and $`D`$ is a covariant derivative defined by
$$Q=\left(\begin{array}{cc}\frac{2}{3}& 0\\ 0& \frac{1}{3}\end{array}\right)$$
(4)
$$D_\mu U=_\mu Ui\mathrm{e}A_\mu [Q,U].$$
(5)
The abelian gauge field $`A_\mu `$ and electric charge $`\mathrm{e}`$ are in unrationalised units, and $`\mathrm{}=c=1`$.
In the spherically symmetric case with a magnetic charge the gauge field has the form
$$A=p(1\mathrm{cos}\theta )d\varphi +\mathrm{\Phi }dt$$
(6)
where $`p`$ is a magnetic charge and the Dirac quantisation condition is $`p\mathrm{e}=1`$. The usual Skyrmion has a magnetic moment which would interact with a magnetic monopole and break the spherical symmetry. We use instead a non-topological ansatz for the chiral field,
$$U=e^{if(r,t)\sigma _3}.$$
(7)
Despite the fact that this field is made up of neutral pions and commutes with the charge matrix $`Q`$, it has a non-zero total electric and baryonic charge due to the effects of anomalies, as we shall see in the next section.
## III Baryon Number and Electric Charge
The gauge invariant baryon current was constructed by Callan and Witten
$`j_B^\mu `$ $`=`$ $`{\displaystyle \frac{ϵ^{\mu \nu \alpha \beta }}{48\pi ^2}}[\mathrm{tr}U^1_\nu UU^1_\alpha UU^1_\beta U`$ (10)
$`+3i\mathrm{e}A_\nu \mathrm{tr}Q(U^1_\alpha UU^1_\beta U_\alpha UU^1_\beta UU^1)`$
$`+3i\mathrm{e}_\nu A_\alpha \mathrm{tr}Q(U^1_\beta U+_\beta UU^1)].`$
Only the third term survives for the chiral field anzatz (7), with a resulting baryon charge
$$n_B=\frac{ep}{2\pi }\left[f(\mathrm{})f(0)\right].$$
(11)
The solution with the boundary conditions $`f(0)=0`$ and $`f(\mathrm{})=2\pi `$ posesses unit baryon number and hence it can be interpreted as a baryon surrounding the monopole.
If the field anzatz (7) is substituted into the meson and electromagnetic interaction terms in the Lagrangian, they become
$`_1`$ $`=`$ $`{\displaystyle \frac{F_\pi ^2}{8}}(f)^2,`$ (12)
$`_2`$ $`=`$ $`{\displaystyle \frac{e^2}{8\pi ^2}}E_iB^if,`$ (13)
$`_3`$ $`=`$ $`{\displaystyle \frac{1}{8\pi }}\left(E^2B^2\right),`$ (14)
where the index $`i=1,2,3`$ and the electromagnetic fields $`E_i`$ and $`B_i`$ are defined by $`F_{0i}=(g_{tt})^{1/2}E_i`$ and $`F_{ij}=ϵ_{ijk}B^k`$. When combined,
$$_2+_3=\frac{1}{8\pi }\left(E\frac{e^2}{2\pi }Bf\right)^2\frac{1}{8\pi }B^2\left(1\frac{e^4}{4\pi ^2}f^2\right)$$
(15)
The extrema of the action occur when the electric field is given by
$$E=\frac{e^2}{2\pi }Bf$$
(16)
This situation is reminisent of the factorisation of the Lagrangian that occurs for a BPS monopole .
The electric field implies a total charge
$$q=\frac{e^2p}{2\pi }f$$
(17)
or asymptotically
$$q_{\mathrm{}}=n_Be$$
(18)
and the $`n_B=1`$ solution can therefore be interpreted as a proton.
If a black hole appears in the background, the inner boundary condition for the field $`f`$ should be imposed not at the origin but at the event horizon $`r=r_+`$. Thus the baryon number in the presence of an event horizon will be defined as
$$n_B=\frac{ep}{2\pi }\left[f(\mathrm{})f(r_+)\right].$$
(19)
If $`f(r_+)=0`$ the baryon number is still an integer and conserved. This configuration represents a proton tightly bound to the black hole. On the other hand if $`f`$ takes some positive value at the horizon the baryon number is not an integer and the skyrmion carries fractional baryon number and electric charge. This configuration will be interpreted as a proton partially swallowed by the black hole. In particular, $`f(r_+)=2\pi `$ means that the black hole has swallowed a whole proton leaving nothing outside the horizon.
It is interesting to observe that, while the baryon number disappears inside the horizon, the electric charge of the black hole can still be measured outside, turning the monopole black hole into a dyon black hole. Therefore, while the baryon number conservation is violated, charge conservation is not violated.
## IV Extremal black hole solutions
In the extremal case we can obtain a general solution based on the Papapetrou-Majumbar metrics . We begin with the background metric fixed and later generalise to solve the full Einstein equations with chiral matter. The Papapetrou Majumbar metrics have the form
$$ds^2=U^2dt^2+U^2(dx^2+dy^2+dz^2)$$
(20)
where
$$U=1+\underset{n=1}{\overset{n_M}{}}\frac{GM_n}{R_n}$$
(21)
and $`R_n`$ is the ordinary Euclidean distance from the point mass $`M_n`$ located in three dimensional space. We also associate these point masses with magnetic charges $`P_n=G^{1/2}M_n`$, and the magnetic field
$$B=G^{1/2}U^1U.$$
(22)
The matter lagrangian obtained earlier (12) has the form
$$=\frac{1}{8}F_\pi ^2(f)^2\frac{1}{8\pi }B^2\left(1+\alpha ^2f^2\right)$$
(23)
where we will set
$$\alpha =\frac{e^2}{2\pi },\mu ^2=\pi GF_\pi ^2$$
(24)
The skyrme field equation obtained from the Lagrangian on this background becomes
$$\mu ^2^2f+\alpha ^2U^2(U)^2f=0$$
(25)
We loose no generality by taking equal charges $`P_n=p`$. The solution with baryon number $`n_B=n_M`$ is then
$$f=2\pi U^s$$
(26)
where
$$s=\frac{1}{2}+\sqrt{\frac{1}{4}+\frac{\alpha ^2}{\mu ^2}}$$
(27)
Since $`F_\pi m_{pl}`$, we can use $`s\alpha /\mu `$ for the pion model.
For a single black hole, the Reissner-Nordstrom coordinate $`r`$ is related to $`R`$ by $`R=rr_+`$ and we obtain
$$f=2\pi \left(1\frac{r_+}{r}\right)^s.$$
(28)
The field is effectively expelled by the black hole and vanishes on the horizon $`r=r_+`$.
The mass of the chiral field configuration can be obtained by integrating the Lagrangian (23),
$$m_f=\frac{1}{8}F_\pi ^2_\mathrm{\Sigma }f_ifdS^i$$
(29)
where $`\mathrm{\Sigma }`$ is a large surface containing all of the masses. This gives
$$m_f=2\pi ^3sF_\pi ^2G\underset{n}{}M_n\pi ^{3/2}n_BeF_\pi $$
(30)
The total mass in the chiral field is much less than one baryon mass per mass point. We can see how it is energetically favourable for a free skyrmion to change its internal configuration from the original Skyrme form to the simpler form used here when it comes into contact with a black hole monopole. The topological description of this transformation for a single monopole is exactly as described in reference .
It is interesting to see that the electrostatic energy cancells due to the factorisation occuring in the lagrangian (15). The chiral field mass is independent of the separation of the holes and therefore there are no forces between then. This is similar to situation for BPS monopole solutions , and suggests that there is a solution of the full Einstein-matter system. This existence of this solution will now be demonstrated.
The spatial part of the Einstein tensor for the matric (20) is
$$G_{ij}=2U^2(_iU)(_jU)+U^2(U)^2g_{ij}$$
(31)
and the Ricci scalar is
$$R=2U^1^2U$$
(32)
Substituting the Einstein tensor for the Lagrangian (23) into the Einstein equations gives
$`U^1^2U`$ $`=`$ $`\mu ^2(f)^2`$ (33)
$`U^2(_iU)(_jU)`$ $`=`$ $`\mu ^2(_if)(_jf)+GB_iB_j(1+\alpha ^2f^2)`$ (34)
These make up a complete system of equations when we include the Maxwell equation
$$(UB)=0$$
(35)
The second Einstein equation implies that $`f`$, $`U`$ and $`B`$ are all parallel. We therefore impose a condition $`ff(u)`$, $`B=b(u)U^1U`$, where
$$u=\mu ^1\mathrm{log}U$$
(36)
The system of equations becomes equivalent to an ordinary differential equation with independent parameter $`u`$,
$$f^{\prime \prime }+\mu \left(1+f^2\right)f^{}\alpha ^2b^2f=0$$
(37)
where
$$b^2=\frac{1+f^2}{1+\alpha ^2f^2}.$$
(38)
The horizon corresponds to $`u\mathrm{}`$ and the far region to $`u0`$. The horizon must therefore be at a critical point of the first order system corresponding to (37). There is only one critical point, $`(f,f^{})=(0,0)`$, hence
$$(f,f^{})(0,0)\text{ as }u\mathrm{}.$$
(39)
Since the critical point is a saddle, the solution is unique and exists for all values of $`\mu `$. Having obtained the unique solution to (37), we then define
$$V(u)=1+\mu _u^0b(x)e^{\mu x}𝑑x.$$
(40)
It is easily seen from (36) that
$$_iV=V^{}_iu=UB$$
(41)
Hence the Maxwell equation implies $`^2V=0`$ and we can write
$$V=1+\underset{n}{}\frac{GM_n}{R_n}$$
(42)
Inverting (40) gives $`u(V)`$.
## V Spherically Symmetric Solutions
In the non-extremal case we shall consider spherically symmetric metrics which can be parameterised in the form
$$ds^2=\frac{\mathrm{\Delta }}{r^2}e^{2\delta }dt^2+\frac{r^2}{\mathrm{\Delta }}dr^2+r^2d\mathrm{\Omega }^2$$
(43)
where $`\mathrm{\Delta }`$ and $`\delta `$ are functions of $`r`$ and $`t`$. After inserting the metric and the other field ansatzë (6)-(7) into the Einstein field equations, one obtains
$`\left(\mathrm{\Delta }e^\delta f^{}\right)^{}{\displaystyle \frac{\lambda ^2}{r^2}}e^\delta f`$ $`=`$ $`{\displaystyle \frac{2r^4}{\mathrm{\Delta }^3}}e^\delta \dot{\mathrm{\Delta }}\dot{f}{\displaystyle \frac{r^4}{\mathrm{\Delta }}}e^\delta \dot{\delta }\dot{f}+{\displaystyle \frac{r^4}{\mathrm{\Delta }^2}}e^\delta \ddot{f}`$ (44)
$`\delta ^{}`$ $`=`$ $`\mu ^2r\left({\displaystyle \frac{r^4}{\mathrm{\Delta }}}e^{2\delta }\dot{f}^2+f^2\right)`$ (45)
$`e^\delta \left({\displaystyle \frac{\mathrm{\Delta }e^\delta }{r}}\right)^{}`$ $`=`$ $`1{\displaystyle \frac{\mu ^2\lambda ^2}{r^2}}f^2{\displaystyle \frac{Gp^2}{r^2}}`$ (46)
where $`\mu `$ and $`\lambda `$ are constants,
$$\mu ^2=\pi F_\pi ^2G,\lambda ^2=\frac{e^4p^2}{4\pi ^3F_\pi ^2}$$
(47)
The electric charge within a sphere of radius $`r`$ is given by equation (17).
For very small $`\mu `$, which is the case for pions, the chiral field has little effect on the background metric and we may take $`\delta =0`$ and express $`\mathrm{\Delta }`$ in terms of the mass $`M`$, electric charge $`Q`$ and magnetic charge $`p`$ of the black hole as
$$\mathrm{\Delta }=r^22GMr+G(Q^2+p^2).$$
(48)
The skyrme field equation (44) on this background is therefore
$$\left(\mathrm{\Delta }f^{}\right)^{}\frac{\lambda ^2}{r^2}f=0,$$
(49)
This should be solved subject to the boundary condition on $`f_{\mathrm{}}`$ which fixes the total charge,
$$q_{\mathrm{}}=\frac{\mathrm{e}^2p}{2\pi }f_{\mathrm{}}=n_B\mathrm{e}.$$
(50)
The non-extremal black hole posseses two horizons at $`r=r_{}`$ and $`r=r_+`$ $`(r_+>r_{})`$, related to the mass and charge by
$$GM=\frac{1}{2}(r_{}+r_+),GQ^2=r_{}r_+Gp^2.$$
(51)
The solution to equation (49) can be obtained analytically,
$$f=2\pi n_B\frac{P_q\left(\frac{r_++r_{}}{r_+r_{}}\frac{2r_+r_{}}{r_+r_{}}\frac{1}{r}\right)}{P_s\left(\frac{r_++r_{}}{r_+r_{}}\right)}$$
(52)
where $`P_s(z)`$ is a Legendre function and
$$s=\frac{1}{2}+\sqrt{\frac{1}{4}+\frac{\lambda ^2}{r_+r_{}}}$$
(53)
The black hole becomes a dyon with electric charge related to the value of $`f`$ at the event horizon,
$$Q=\frac{n_Be}{P_s\left(\frac{r_++r_{}}{r_+r_{}}\right)}.$$
(54)
This relation can be solved, together with (51), to obtain $`QQ(M)`$, showing the existence of a one parameter family of solutions ($`n_B`$ and $`p`$ being regarded as fixed). In particular, $`Q0`$ as $`M`$ approaches the extremal limit $`pm_{pl}`$ and the meson field is expelled from the hole.
Larger values of $`\mu `$ may be realised for a hypothetical model where $`U`$ is unrelated to pions and this is discussed below. We will consider a static solution. We can replace $`\mathrm{\Delta }`$ by a mass function $`m(r)`$, defined implicitly by the relation
$$\mathrm{\Delta }=r^22Gmr+G(p^2+q^2)$$
(55)
where the charge is given by equation (17). The equations become
$`m^{}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }}{2Gr}}\delta ^{}+\mu ^2\lambda ^2ff^{}`$ (56)
$`\delta ^{}`$ $`=`$ $`\mu ^2r(f^{})^2`$ (57)
$`f^{\prime \prime }`$ $`+`$ $`\left({\displaystyle \frac{\mathrm{\Delta }^{}}{\mathrm{\Delta }}}+\delta ^{}\right)f^{}{\displaystyle \frac{\lambda ^2}{r^2\mathrm{\Delta }}}f=0`$ (58)
Suitable boundary conditions are $`f2\pi `$ and $`\delta 0`$ as $`r\mathrm{}`$.
In the numerical results the fields are scaled to the horizon size,
$$\widehat{r}=\frac{r}{r_+},\widehat{m}=\frac{Gm}{r_+}$$
(59)
The solutions are parameterised by a parameter $`\widehat{p}`$, defined by
$$\widehat{p}^2=\frac{Gp^2}{r_+^2}$$
(60)
which is restricted to $`\widehat{p}1`$.
The extremal black hole solutions have $`\mathrm{\Delta }_+=\mathrm{\Delta }_+^{}=0`$. The regular solution to equation (58) has,
$$f_+=0,Q=\frac{\mathrm{e}}{2\pi }f_+=0,\widehat{p}=1.$$
(61)
Hence the proton lies fully outside the black hole, as we saw before. The numerical solution for $`f`$ is shown in fig.(1). This agrees well with the result (28) of the previuos section, because the value of $`\mu `$ used here is still quite small. The results are still qualitatively similar for chiral models with $`\mu `$ of order one.
For the non-extremal solution, we begin the integration of the field equations close to the horizon, with
$`\widehat{m}`$ $`=`$ $`\widehat{m}_0+\widehat{m}_1(\widehat{r}1)+\widehat{m}_2(\widehat{r}1)^2+\mathrm{}`$ (62)
$`\delta `$ $`=`$ $`\delta _++\delta _1(\widehat{r}1)+\mathrm{}`$ (63)
$`f`$ $`=`$ $`f_++f_1(\widehat{r}1)+\mathrm{}`$ (64)
where $`\delta _+`$ and $`f_+`$ are shooting parameters determined so as to satisfy the boundary conditions at infinity.
As can be seen from the above expansion, the skyrme field must have a nonzero value at the horizon otherwise the only allowed solution is the trivial one. Consequently the nonextremal black hole acquires an electric charge
$$Q=\frac{\mathrm{e}}{2\pi }f_+,$$
(65)
and allows the skyrmion to have fractional electric charge. The numerical results for this solution are shown in fig.(2)-(3). Again, these agree well with the fixed background for small values of $`\mu `$.
We have a single one parameter family of solutions with $`\widehat{p}1`$. In fig.(4)-(5), we plot the horizon radius $`r_+`$ and skyrmion mass $`m_f`$ as functions of black hole mass $`M`$. Figure (4) is related to the entropy of the black hole ($`4\pi r_+^2`$). The other figure shows how the proportion of the skyrmion which is swallowed by the black hole increases with the black hole mass.
The last figure (6) shows how the horizon value of $`f`$ changes as the coupling constant $`\mu `$ changes. Since $`\mu `$ characterises the mass scale of the chiral model, this is amounts to a comparison of different models. As can be seen from the figure, for small $`\mu `$ the electromagnetic interaction is dominant so that the skyrmion is absorbed by the black hole to a lesser extent. On the other hand for large $`\mu `$ the gravitational interaction is dominant so that most of the skyrmion is absorbed by the black hole.
## VI Stability Analysis
In this section we show that the skyrmion solutions which we have obtained are stable under spherically symmetric linear perturbations. We shall begin with the analysis of a skyrmion on the fixed background.
In the fixed background case the skyrme field is the only perturbed field and can be expanded about the skyrmion solution $`f_0`$ by writing
$$f(r,t)=f_0(r)+e^{i\omega t}\xi (r).$$
(66)
Equation (66) is inserted into equation (49) to obtain the eigenvalue equation
$$\left(\mathrm{\Delta }_0\xi ^{}\right)^{}+\frac{\lambda ^2}{r^2}\xi =\frac{r^4}{\mathrm{\Delta }_0}\omega ^2\xi ,$$
(67)
where the background equation has been used.
If $`\omega `$ is real and $`\omega ^2>0`$ the solution is stable, and if $`\omega `$ is imaginary and $`\omega ^2<0`$ it is unstable since the mode can grow or decay exponentially under the small perturbation. To show which is the case we multiply both sides of equation (67) by $`\xi `$ and integrate in $`r`$ from the horizon to infinity
$$_{r_+}^{\mathrm{}}\left[\frac{\mathrm{\Delta }_0}{2}\xi ^2+\frac{\lambda ^2}{r^2}\xi ^2\right]𝑑r=\omega ^2_{r_+}^{\mathrm{}}\frac{r^4}{\mathrm{\Delta }_0}\xi ^2𝑑r,$$
(68)
where integration by parts and boundary conditions have been used. It can be seen that the integrands of both sides are positive definite, which means that $`\omega ^2>0`$. Hence the skyrmion on the fixed background is linearly stable .
Next we analyse the stability of the skyrmion with backreaction. In this case we have to expand the metric as well as the skyrme field around the classical solutions $`f_0`$, $`\delta _0`$ and $`\mathrm{\Delta }_0`$
$`f(r,t)`$ $`=`$ $`f_0(r)+f_1(r,t)`$
$`\delta (r,t)`$ $`=`$ $`\delta _0+\delta _1(r,t)`$
$`\mathrm{\Delta }(r,t)`$ $`=`$ $`\mathrm{\Delta }_0+\mathrm{\Delta }_1(r,t).`$
These are substituted into eq.(44)-(46) to obtain the following coupled equations up to first order
$`\left[\left(\mathrm{\Delta }_0\delta _1f_0^{}+\mathrm{\Delta }_0f_1^{}+\mathrm{\Delta }_1f_0^{}\right)e^{\delta _0}\right]^{}{\displaystyle \frac{\lambda ^2}{r^2}}\left(\delta _1f_0+f_1\right)e^{\delta _0}`$ $`=`$ $`{\displaystyle \frac{r^4}{\mathrm{\Delta }_0}}e^{\delta _0}\ddot{f}_1`$ (69)
$`\delta _1^{}`$ $`=`$ $`2\mu ^2rf_0^{}f_1^{}`$ (70)
$`\left({\displaystyle \frac{2\mu ^2\lambda ^2}{r^2}}f_0f_1+{\displaystyle \frac{\mathrm{\Delta }_0}{r}}\delta _1^{}\right)e^{\delta _0}`$ $`=`$ $`\left({\displaystyle \frac{\mathrm{\Delta }_1}{r}}e^{\delta _0}\right)^{}.`$ (71)
Equation (71) can be integrated with the help of the static field equation,
$$\mathrm{\Delta }_1=2\mu ^2r\mathrm{\Delta }_0f_0^{}f_1.$$
(72)
Substituting equation (70) and equation (72) into equation (69) one obtains the first order equation for $`f_1`$
$$\left(\mathrm{\Delta }_0e^{\delta _0}f_1^{}\right)^{}\left[2\mu ^2\left(r\mathrm{\Delta }_0e^{\delta _0}f_0^2\right)^{}+\frac{\lambda ^2(M_c)^2}{r^2}e^{\delta _0}\right]f_1=\frac{r^4}{\mathrm{\Delta }_0}e^{\delta _0}\ddot{f}_1.$$
(73)
Setting $`f_1(r,t)=\xi (r)e^{i\omega t}`$ one obtains an eigenvalue equationfor $`\xi `$,
$$\left(\mathrm{\Delta }_0e^{\delta _0}\xi ^{}\right)^{}+\left[2\mu ^2\left(r\mathrm{\Delta }_0e^{\delta _0}f_0^2\right)^{}+\frac{\lambda ^2(M_c)^2}{r^2}e^{\delta _0}\right]\xi =\omega ^2\frac{r^4}{\mathrm{\Delta }_0}e^{\delta _0}\xi .$$
(74)
We introduce the tortoise coordinate $`r^{}`$ such that
$$\frac{dr^{}}{dr}=\frac{1}{\mathrm{\Delta }_0e^{\delta _0}}$$
(75)
and $`r^{}`$ runs from $`\mathrm{}`$ to $`+\mathrm{}`$ as $`r`$ runs from $`r_+`$ to $`+\mathrm{}`$. Then eq.(74) is reduced to the Sturm-Liouville equation
$$\frac{d^2\xi }{dr^2}+U\xi =\omega ^2r^4\xi ,$$
(76)
where
$$U=\left[2\mu ^2\left(r\mathrm{\Delta }_0e^{\delta _0}f_0^2\right)^{}+\frac{\lambda ^2}{r^2}e^{\delta _0}\right]\mathrm{\Delta }_0e^{\delta _0}.$$
(77)
On the left-hand side we have $`r^4`$, which makes the equation different from the previous eigenvalue equation. However, since $`r^4`$ remains positive through the whole space, the same conditions for stability as the ordinary eigenvalue equation can be applied. As can be seen by examining $`U`$, $`U0`$ as $`rr_+`$, (i.e. $`r^{}\mathrm{}`$), $`UU_{\mathrm{}}>0`$ as $`r\mathrm{}`$, and $`U>0`$ in between. In addition the solution does not change its shape for any value of the coupling constant $`\mu `$. Therefore we can safely conclude that a skyrmion with backreaction is also linearly stable.
## VII Conclusion
We are now able to describe some of the features of the interaction between a slowly moving proton and a black hole monopole. The free skyrmion has a magnetic moment and, if it has the correct orientation, it will be attracted to the monopole. When the proton approaches the black hole monopole, the fields rearrange themselves into the energetically prefered configuration of skyrmion hair solution described in this paper.
The skyrmion solution on the fixed black hole background was obtained analytically in both extremal and nonextremal cases. In the extremal case, the baryon number and electric charge are expelled by the horizon and the system represents protons bound to monopole black holes. We found a general solution on the Eintein equations in the extremal case with many similarities to the BPS monopole system.
In the nonextremal case, the monopole black hole partially swallows the proton and transforms to a dyon black hole, leaving fractional electric charge and baryon number outside the horizon. However, Hawking radiation cannot be ignored in this case.
We have also obtained numerical skyrmion solutions in the nonextremal case with gravitational backreaction. The effect of the backreaction is important only for very massive skyrmions i.e. very large coupling constant. The results are qualitatively similar to the solutions without backreaction.
Since the solutions are stable, baryon decay can only take place when extra particle fields are included in the model. If we introduce a charged field $`\varphi `$ on the fixed background, then following the proceedure described in reference , the equations become (approximately)
$`\left(\mathrm{\Delta }f^{}\right)^{}{\displaystyle \frac{\lambda ^2}{r^2}}(f\varphi )={\displaystyle \frac{r^4}{\mathrm{\Delta }}}\ddot{f},`$ (78)
$`\left(\mathrm{\Delta }\varphi ^{}\right)^{}+{\displaystyle \frac{\lambda ^2}{r^2}}(f\varphi )={\displaystyle \frac{r^4}{\mathrm{\Delta }}}\ddot{\varphi },`$ (79)
The stability arguments no longer apply. The dynamical process of a black hole swallowing a proton can be examined by solving these time-dependent field equations numerically.
###### Acknowledgements.
We are grateful for discussions with K. Maeda. NS was partially supported by a British Council ORS scholarship.
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# 28SiO , 29SiO and 30SiO excitation: effects of infrared line overlaps
## 1 Introduction
SiO masers are excited in the inner layers of circumstellar envelopes of late-type stars, before the dust formation zone, but beyond the photosphere where outflow and inflow of matter or shocks coexist and give rise to complex physical conditions.
Since the discovery of SiO maser emission by Snyder & Buhl (1974), many papers have been published to explain the SiO maser phenomenon, and several improvements involving collisional and radiative pumping have been proposed to model the emission from late-type stars. However, it is still difficult to explain the disparities between some rotational line intensities observed inside a same vibrational state or, more generally, in adjacent transitions of <sup>30</sup>SiO , <sup>29</sup>SiO and $`v3`$ <sup>28</sup>SiO . Olofsson et al.(1981, 1985) introduced the hypothesis of a line overlap between two ro-vibrational lines of SiO and water to explain the weakness of the <sup>28</sup>SiO $`v=2,J=21`$ emission. This emission was observed by Bujarrabal et al.(1996) in stars with different O-abundances, and the overlap between two infrared lines of H<sub>2</sub>O and SiO was confirmed. In addition, several overlaps between the infrared lines of the main (<sup>28</sup>SiO ) and rare (<sup>29</sup>SiO and <sup>30</sup>SiO ) isotopic species of silicon monoxyde have been suggested to play a role in the excitation of these species. Using the LVG approximation Cernicharo et al.(1991) modeled the pumping of <sup>29</sup>SiO by <sup>28</sup>SiO . The importance of overlaps was further investigated by Cernicharo & Bujarrabal (1992) and Cernicharo et al.(1993). Recently, González-Alfonso et al.(1996) and González-Alfonso & Cernicharo (1997) used a non-local radiative transfer code to study the question of line overlaps. They successfully explained most of the <sup>29</sup>SiO , <sup>30</sup>SiO and high-v ($`v3`$) <sup>28</sup>SiO masers; however, several lines remain unexplained.
Our main goal in this work is to develop a simple and modular calculation tool in order to investigate the impact of various line overlap effects among the three main isotopic forms of SiO. We do not intend here to realistically model the SiO maser line profiles. We rather are interested in predicting and discussing relative peak line intensity ratios. In Sect. 2 we give details on spectroscopy, the used radiative transfer code and physical parameters, and on validation of our code. In Sect. 3, we include several line overlaps and compare our results to the observations. Some major results of this work are summarized in Sect. 4.
## 2 The Model
The LVG code used here must be regarded as a simple but powerful computational tool for maser emission prediction. Bujarrabal (1994) has clearly demonstrated that the LVG code applied to the SiO maser zone is a fairly good approximation in comparison with the exact radiative transfer solution. Other uncertainties do not play a minor role in comparison with the differences between the LVG and the exact radiative transfer results; in particular uncertainties in the collision laws and the exact geometry of the maser region result in weaker predictions. Thus, we believe that the physical laws and the spherical symmetry adopted in our model are sufficient for the LVG approximation.
### 2.1 Spectroscopy, collisions and radiative transfer
Modeling the SiO emission requires the use of accurate spectroscopic data and collisional rates. Energy levels of SiO have been calculated using the most recent calculations of the Dunham coefficients (Campbell et al. 1995). For the Einstein A-coefficients we used the recent work of Drira et al.(1997).
Basic data concerning the rotational and rovibrational rate coefficients of SiO with H<sub>2</sub> and their dependence on temperature were taken from Bieniek & Green (1983). We also used the new rate coefficients derived by Lockett & Elitzur (1992) for $`v3`$ and $`\mathrm{\Delta }v=2`$ transitions and we applied corrections to their rates to compensate for the effect of the limited number of rotational and vibrational levels taken in the calculations. In addition, Langer & Watson (1984) made a first attempt to estimate collision rates of SiO with atomic hydrogen from a comparison with the expected H-CO and H<sub>2</sub>-CO rates. The main features of these H-SiO rates used in our work are: the vibrational collision rates are 10 times larger than those for H<sub>2</sub>, and the rotational collision rates are decreased by a factor of 7. H can be formed in two different ways: (i) the stellar radiation may photodissociate H<sub>2</sub> close to the star, thus forming an atomic hydrogen layer; (ii) shocks may dissociate $`H_2`$ to produce a uniform distribution of H throughout the envelope. We assume here that $`\chi [H_I/H_2]`$ is constant and equal to $`10^2`$ as in Bowers & Knapp (1987). Adding atomic hydrogen to H<sub>2</sub> increases the role of collisions. Note however that its impact on the inversions is weak (a few % for $`\chi [H/H_2]<10^1`$). Beyond this value, inversions for the lowest rotational levels tend to decrease contrary to what is observed for high $`J`$ transitions. The higher the vibrational level, the larger the effect of the atomic hydrogen is; this is expected because of larger vibrational collision rates for H than for H<sub>2</sub>. Nevertheless, we stress that these collision rates and the relative abundance of the atomic hydrogen are poorly known.
Because the purpose of this work is to explain in a simple way the main characteristics of <sup>28</sup>SiO , <sup>29</sup>SiO and <sup>30</sup>SiO emissions, we have used the LVG (Sobolev 1958) approximation, based on the model of a homologously expanding envelope, in which the escape probability does not depend on the angle. If the velocity gradient (logarithmic velocity gradient $`\epsilon _r=d\mathrm{ln}V/d\mathrm{ln}r`$) in the circumstellar envelope is large enough, each cell of the discretized medium will be decoupled from all other cells, thus allowing a local treatment of the radiative transfer. The interaction area is limited to a small zone where photons can be absorbed or emitted around a resonance point. We have limited the maximum size of the cells to the Sobolev length, $`L`$, defined by
$$L=r\frac{\mathrm{\Delta }V_D/V}{\epsilon _r}$$
with
$$\mathrm{\Delta }V_D=\sqrt{\mathrm{\Delta }V_{th}^{}{}_{}{}^{2}+\mathrm{\Delta }V_{turb}^{}{}_{}{}^{2}}$$
where $`r`$ is the radial distance, $`\mathrm{\Delta }V_{th}`$ the thermal velocity, $`\mathrm{\Delta }V_{turb}`$ the turbulence, and $`V`$ the expansion velocity. $`\mathrm{\Delta }V_D`$ characterizes the local absorption coefficient. The formula giving the Sobolev length applies to any form of expansion velocity and to the general case in which the photon escape probability depends on the direction. For a spherically symmetric envelope and for the special case where $`V`$ is proportional to $`r`$, the angle-dependent terms vanish, the logarithmic gradient $`\epsilon _r=1`$, and the escape probability is isotropic and proportional to the inverse of the peak opacity (as soon as the opacity becomes large). The Sobolev length is just the length over which the velocity varies by an amount corresponding to the width of the local absorption coefficient. One essential advantage of the LVG code is that it requires less computer time than more exact treatments of the radiative transfer problem based on integral solutions (Bujarrabal 1994) or on the Monte Carlo method (González-Alfonso & Cernicharo 1997). In the frame of the LVG formalism and for physical conditions relevant to the circumstellar environment of late-type stars we have derived the solutions of the statistical equilibrium equations and the opacities of various SiO lines. More complex effects related to non spherical geometry, polarization or saturation and beaming angles are beyond the scope of this paper. Concerning saturation, we emphasize that our rate equations are solved for both positive and negative opacities, but that beaming related to saturation cannot be investigated here because we assume the escape probability to be isotropic. Note that beaming angles for saturated masers in the case of the LVG approximation remain an unsolved problem (Elitzur 1992). Saturation and the effect of competitive gains in SiO lines were investigated in Doel et al.(1995).
Our calculations include 40 rotational levels for each of the 5 vibrational states $`v=0`$ to 4. To correct for the limited number of rotational and vibrational levels (this overestimates the opacities), we have applied a correction similar to that used by Bujarrabal & Rieu (1981) for the rotational populations. The required final precision on the populations is better than $`10^4`$. Calculations are done simultaneously for <sup>28</sup>SiO , <sup>29</sup>SiO and <sup>30</sup>SiO . We derive the populations for each level ($`v,J`$) and derive the opacities for all allowed transitions.
### 2.2 Circumstellar parameters
We have used physical parameters appropriate to the circumstellar environment of evolved stars (see details in Herpin 1998).
#### $``$ Kinetic temperature <br>
The gas temperature follows the law adopted by Langer & Watson (1984):
$$T_c=T_{}(\frac{r}{R_{}})^{0.5}$$
where $`R_{}`$ is the stellar radius ($`R_{}=\mathrm{7.7\; 10}^{13}`$cm), $`T_{}`$ is the stellar temperature ($`T_{}=2500`$ K) and r is the radial distance. This law is in rough agreement with the observations of SiO thermal emission of Lucas et al.(1992) and Bujarrabal et al.(1986), and also with the calculations of Willson (1987) who finds 1500 K in the maser region, for a typical Mira. Turbulence is a free parameter in our code and contributes to net line broadening.
#### $``$ Expansion velocity <br>
The expansion velocity of the gas adopted in our model is:
$$V=V_{\mathrm{}}\frac{1}{1+(\frac{R_{acc}}{r})^{e_v}}$$
where $`R_{acc}=\mathrm{1.8\; 10}^{15}`$ cm, $`e_v[1;3]`$ and $`V_{\mathrm{}}`$ is the asymptotic expansion velocity of the gas (9.5 km s<sup>-1</sup>). A logarrithmic velocity gradient of 3 was used in the calculations.
#### $``$ SiO abundance <br>
The SiO$`/`$H<sub>2</sub> abundance ratio is uncertain, and we adopt the solar abundance by default ($`\chi _0[`$SiO$`/`$H$`{}_{2}{}^{}]=\mathrm{5.\; 10}^5`$). This abundance decreases with the distance to the star, the SiO molecules being progressively condensed onto the grains. We follow the law adopted by Bujarrabal et al.(1989), and we adopt for the <sup>29</sup>SiO and <sup>30</sup>SiO isotopic abundances ($`A_i`$) relative to <sup>28</sup>SiO , $`1/19.5`$ and $`1/29.5`$, respectively (see fig.1):
$$\chi (r)=\chi _0[1\frac{\kappa (r)}{\kappa _0}+\chi ^{}]/A_i$$
with $`\chi _0`$ the initial ratio (here the solar abundance), and $`\chi ^{}`$ the remaining fraction of SiO after the grain formation ($`10^2`$). The terms $`\kappa (r)`$ and $`\kappa _0`$ depend respectively on the radiative pressure and the total quantity of grains formed. We take:
$$\kappa (r)=\kappa _0\frac{(rR_1)^2}{(rR_1)^2+(R_dR_1)^2}$$
where $`R_d`$ is the radius where half the total quantity of grains is formed ($`\mathrm{1.5\; 10}^{15}`$ cm), and $`R_1`$ is the initial point ($`\mathrm{2.\; 10}^{14}`$ cm).
#### $``$ H<sub>2</sub>density <br>
We adopt (Elitzur 1992):
$$n_{H_2}=\frac{\dot{M}}{8\pi r^2Vm_p}$$
with $`V`$ the expansion velocity, $`\dot{M}`$ the steallar mass loss rate and $`m_p`$ the proton mass (see Fig.1). Maser effects will require densities of a few $`10^9`$ particles/cm<sup>3</sup>, whereas thermalization is clearly reached beyond $`10^{10}10^{11}`$ particles/cm<sup>3</sup>.
#### $``$ Grain radiation <br>
IR spectra from the envelopes of late-type stars show that circumstellar grains have various chemical compositions. Here we only consider silicate grains and we assume that most circumstellar grains lie in a shell outside the SiO maser zone (with $`R_{in}=10^{15}`$ cm, and $`R_{out}=10^{16}`$ cm). The extinction coefficient of the dust is $`Q_\nu Q_0(\frac{\lambda }{\lambda _0})^p`$ with $`p[1;2]`$ for wavelengths $`\lambda 1\mu `$m (Mathis 1990). For silicate grains, we take $`p=1.1`$, $`\lambda _0=80\mu `$m, and $`Q_0=2.10^3`$ according to the work of Ivezic & Elitzur (1995). We consider that the silicate grains behave as black-bodies with temperature $`T_d`$. We assume that the dust is optically thin, and that $`T_d`$ is constant throughout the envelope ($`T_d=`$600 K). We take, following Netzer & Elitzur (1993), $`\chi _d=\frac{M_{dust}}{M_{H_2}}=10^2`$, $`3.0`$g.cm<sup>-3</sup> for the grain volumic density, and $`5.10^6`$cm for the grain radius.
### 2.3 Line Overlaps
In this work, we only deal with local line overlaps due to thermal line broadening and local turbulence in successive circumstellar layers. For simplicity, two nearby spectral lines are considered to overlap when $`\mathrm{\Delta }V\mathrm{\Delta }V_0`$ where $`\mathrm{\Delta }V`$ is the difference in velocity between the line centers, and $`\mathrm{\Delta }V_0`$ is an a priori fixed value. In the examples discussed in Sect. 3 we have adopted $`\mathrm{\Delta }V_0=5`$ km s<sup>-1</sup>corresponding to turbulent inner layers ($`\mathrm{\Delta }V_{thermal}1.2`$ km s<sup>-1</sup>). Line overlaps are treated in a very simple manner. For two nearby lines $`\nu _{12}`$ and $`\nu _{34}`$, we derive the opacities $`\tau _{12}`$, $`\tau _{34}`$, and the individual source functions $`S_{12}`$, $`S_{34}`$. Adding the absorption and emissivity coefficients of each line, the source function becomes
$$S_{12}^{}{}_{}{}^{}=S_{12}(\frac{\tau _{12}}{\tau _{12}+\tau _{34}})+S_{34}(\frac{\tau _{34}}{\tau _{12}+\tau _{34}})$$
and the total opacity is now $`\tau _{12}^{}{}_{}{}^{}=\tau _{12}+\tau _{34}`$. The average intensity is then given by:
$$\overline{J}_{12}=[1\beta _{12}^{}{}_{}{}^{}(\tau ^{})]S_{12}^{}{}_{}{}^{}+\beta _{12bb}^{}{}_{}{}^{}(\tau ^{})I_{bb}$$
$$+\beta _{12bg}^{}{}_{}{}^{}(\tau ^{})(I_{bg}+I_{dust})$$
where $`\beta `$ is the usual escape probability. The cosmic background $`I_{bg}`$ is described by a 3 K blackbody and the dust contribution is given by $`I_{dust}`$. The emission of the central star $`I_{bb}`$ is described by a blackbody spectrum at temperature $`T_{}`$; a geometrical dilution factor is applied. In these formulae, the populations of the 4 levels are involved. If $`\tau _{12}<\tau _{34}`$, after overlap we have $`S_{12}^{}{}_{}{}^{}=S_{34}^{}{}_{}{}^{}S_{34}`$ and $`\tau _{34}^{}\tau _{34}`$; this implies minor changes for $`\tau _{34}`$ and the $`n_3`$ and $`n_4`$ populations. On the other hand, photons from the stronger line will be absorbed by the more optically thin line $`12`$, thus enhancing the non-equilibrium distribution of the populations in the levels 1 and 2.
Such a treatment of the overlaps is oversimplified, but allows us to test quickly the effects of this mechanism. Our code incorporates three overlaps simultaneously; this number can of course be increased. To account for line overlaps, the code first derives the populations and opacities separately for each isotope without overlap, and then solves again for the radiative transfer, combining the overlapping lines. We describe in Sect. 3 how we proceed to test the effects of line overlaps.
In fact, non local line overlaps have long been recognized as important in the excitation of the OH radical, and, more recently, in SiO models (cf. Introduction and Sect. 3). Line overlaps may occur even if two lines are not excited in immediately adjacent gas layers; e.g. a Doppler-shifted line emitted from a first cloud may interact with another line in a second cloud provided that the two clouds have the appropriate relative velocities. In a somewhat different process, Olofsson et al.(1981, 1985) suggested that the anomalously weak $`v=2,J=21`$ emission line in stars could be explained by non local overlap with the H<sub>2</sub>O $`\nu _212_{75}\nu _111_{66}`$ line. Photons emitted by the water line increase the radiation field at the frequency of the vibrational transition of SiO creating an excess of absorption for this SiO line, and thus destroying the inversion in $`v=2,J=21`$. This mechanism was recently investigated in detail by Bujarrabal et al.(1996).
### 2.4 Model validation
To test the validity of our code, we have used the same initial conditions as in Alcolea et al.(1989). (For this comparison, we have not applied corrections on the collision rates for the limited number of levels.) The main differences with Alcolea et al. are the number of rotational levels (40 here compared to 12 in Alcolea et al.), the collision rates (Langer & Watson 1984 in Alcolea et al.), and the more recent spectroscopic coefficients used here. In Figure 2, we plot the $`J=10`$ negative opacities for $`v=0`$ to 4 versus the ratio $`n(SiO)LV^1\epsilon ^1`$ (where $`L`$ is the Sobolev length, $`V`$ the expansion velocity, $`\epsilon `$ the logarithmic velocity gradient). Our results are much similar to those of Alcolea et al. and show that the $`J=10`$ inversions are well differenciated according to the vibrational state. However, some minor discrepancies remain. In particular, our work shows that $`v=1`$ inversions occur over a broader range of parameters. We do not obtain increasing opacities with increasing $`v`$; but in Alcolea et al. this trend was probably the consequence of their smaller number of rotational levels. Finally, we note that results in Fig. 2 are similar to those of Lockett & Elitzur (1992).
## 3 Discussion of line overlap effects
Standard radiative and collisional SiO pumping models are able to predict several of the observed features of SiO maser emission from late-type stars (e.g., Bujarrabal 1994). These models predict similar line profiles and smooth line strength variations when one goes up the rotational energy ladder. They cannot explain, however, the anomalous line intensity ratios observed among several adjacent rotational transitions in high $`v`$ states of <sup>28</sup>SiO and in $`v=0`$ to 3 states of <sup>29</sup>SiO and <sup>30</sup>SiO . On the other hand, line overlap effects may strongly modify the SiO populations and thus explain apparent line anomalies. The most complete work in this domain (González-Alfonso et al.1996, González-Alfonso & Cernicharo 1997) includes a non-local treatment of line overlap effects. Their model correctly predicts several observed features. It also predicts excitation of the <sup>29</sup>SiO $`v=3,J=87`$ transition which was subsequently detected (González-Alfonso & Cernicharo 1997) thus demonstrating the importance of line overlaps. Nevertheless, excitation of several transitions remains unexplained in the work of González-Alfonso et al. In particular, emission from <sup>29</sup>SiO $`v=0,J=54`$ and $`v=1,2,J=65`$ and <sup>30</sup>SiO $`v=0,J=21`$ and $`43`$ must be understood (see discussion below and Table 1).
The present model is conceived as a simple tool to quickly investigate the impact of various line overlaps between ro-vibrational lines of <sup>28</sup>SiO , <sup>29</sup>SiO and <sup>30</sup>SiO . The model incorporates up to 3 simultaneous different overlaps, namely six lines, according to the simple prescription outlined in Sect. 2.3. Once the maximum line shift allowed between two lines is fixed (this depends on $`\mathrm{\Delta }V_0`$ defined in Sect. 2.3), a specific module of our code searches for all possible overlaps. Within the 10 first rotational transitions we find a total of 27, 48 and 79 overlaps for $`\mathrm{\Delta }V_0=5,10`$ and 15 km s<sup>-1</sup>, respectively. There are of course more overlaps involving higher J rotational levels, and we have also explored the influence of high $`J`$-level transitions overlapping low $`J`$-level transitions. As briefly explained in Sect. 2.3, the effect of a line overlap between two lines strongly depends on the relative strength of these two lines. In general, when an overlap occurs between <sup>28</sup>SiO and <sup>29</sup>SiO or <sup>30</sup>SiO transitions, the rare isotopic transitions which are optically thin tend to be more affected than the <sup>28</sup>SiO transitions; and the impact on the more optically thin lines is higher for lower $`v`$ states. Our code derives the net optical depths compared to the non-overlapping cases. Taking $`\mathrm{\Delta }V_0=5`$ km s<sup>-1</sup>results in 31 overlapping line pairs. These overlaps affect <sup>28</sup>SiO , <sup>29</sup>SiO or <sup>30</sup>SiO rotational transitions with $`v4`$ and $`J10`$. There is however the interesting case of <sup>28</sup>SiO $`v=21,J=34`$ which requires $`\mathrm{\Delta }V_010`$ km s<sup>-1</sup>to overlap with <sup>30</sup>SiO $`v=10,J=12`$.
Because of radiative pumping through high $`v`$ levels, we stress that very different ($`v,J`$) levels are connected. That is why some infrared overlaps have sometimes a strong influence on ro-vibrational levels not involved in these overlaps: because of this excitation mechanism numerous very different levels are connected together.
We have grouped our results into four categories: (i) the overlaps invoked by González-Alfonso et al.(1996) and González-Alfonso & Cernicharo (1997) to explain some important newly discovered rotational transitions; (ii) the overlaps which may account for the behaviour of maser rotational transitions explained elsewhere by other models and other overlaps; (iii) the overlaps providing explanations for the maser lines not yet explained; and (iv) the overlaps suggesting a search for new maser lines.
$``$ In the first category it is worth mentioning that the overlap of <sup>28</sup>SiO $`v=21,J=2021`$ with <sup>29</sup>SiO $`v=32,J=89`$, although it involves relatively high $`v`$ states for <sup>29</sup>SiO , results in strong enhancement of the <sup>29</sup>SiO $`v=3,J=87`$ transition. This was also predicted by González-Alfonso et al.(1996). We have also verified that two line overlaps, <sup>28</sup>SiO $`v=21,J=23`$ with <sup>29</sup>SiO $`v=32,J=1110`$, and <sup>29</sup>SiO $`v=32,J=1110`$ with <sup>30</sup>SiO $`v=10,J=01`$, result in a clear enhancement of the <sup>29</sup>SiO $`v=3,J=1110`$ line intensity. Another example is the enhancement of the <sup>30</sup>SiO $`v=2,J=87`$ transition resulting, as in González-Alfonso et al.(1996), from the effects of 3 simultaneous line overlaps.
$``$ We have successfully explained (case (ii)) the enhancement of the <sup>28</sup>SiO $`v=3,J=10`$ and $`v=4,J=54`$ transitions, and of the <sup>30</sup>SiO $`v=0,J=54`$ transition. Standard modelisation does not easily invert the <sup>28</sup>SiO $`v=3,J=10`$ line although it is observed in a variety of sources (Scalise & Lépine 1978, Alcolea et al.1989). One overlap may lead to the excitation of this line: <sup>28</sup>SiO $`v=32,J=23`$ with <sup>29</sup>SiO $`v=21,J=56`$. On the other hand, observations indicate that the <sup>28</sup>SiO $`v=3,J=21`$ line should be quenched. This could result from four different overlaps including the low $`J`$ rotation level overlap of <sup>28</sup>SiO $`v=21,J=34`$ with <sup>30</sup>SiO $`v=10,J=12`$. This overlap is also consistent with the extinction of the $`v=3,J=21`$ line and with the enhancement of the $`v=4,J=54`$ line as observed by Cernicharo et al.(1993) in VY CMa. Our model easily leads to maser emission of the $`v=0,J=10`$ transition. This is observed in several stars for <sup>29</sup>SiO (Alcolea & Bujarrabal 1992) and <sup>30</sup>SiO (Cernicharo & Bujarrabal 1992).
$``$ Table 1 (case (iii)) and Fig.3 suggest explanations of the <sup>29</sup>SiO and <sup>30</sup>SiO maser transitions not explained in the work of González-Alfonso et al.(1996). Only overlaps can explain inversions of the <sup>30</sup>SiO $`v=0,J=21,43`$ (not explained in previous works) or $`J=54`$ transitions. Line overlaps are also required to explain several <sup>29</sup>SiO $`v=1`$ and 2 rotational transitions. For the $`v=1,J=65`$ line observed by Cernicharo & Bujarrabal (1992), Table 1 suggests that it may be explained by several overlaps. Of special interest is the simultaneous overlap of <sup>28</sup>SiO $`v=21,J=43`$ with <sup>29</sup>SiO $`v=10,J=10`$ and of <sup>28</sup>SiO $`v=21,J=54`$ with <sup>29</sup>SiO $`v=10,J=21`$. The same arguments apply to $`v=2,J=65`$ which has also been observed by Cernicharo & Bujarrabal. The overlap of <sup>29</sup>SiO $`v=10,J=34`$ with <sup>30</sup>SiO $`v=10,J=10`$ is clearly important because it enhances simultaneously the observed $`J=43`$ and $`21`$ transitions. It is interesting to note that we have succeeded in producing weak maser emission in <sup>29</sup>SiO $`v=0,J=54`$ (never explained before). Another important result concerns the <sup>30</sup>SiO $`v=1,J=10`$ maser emission produced by our model with different overlaps (see Table 1 and Fig. 4); this line was recently discovered by Cho & Ukita (1998) in TX Cam.
$``$ Table 2 lists some relatively strong maser transitions as predicted by our model (case (iv)). Note that the <sup>29</sup>SiO and <sup>30</sup>SiO $`v=1,J=21`$ maser lines have also been predicted by Rausch et al.(1996) on the basis of an optical pumping model. Four transitions (<sup>28</sup>SiO $`v=4,J=32`$, <sup>29</sup>SiO $`v=1,J=21`$ and $`v=3,J=43`$, <sup>30</sup>SiO $`v=2,J=21`$) are easily excited even without line overlaps; however, all other transitions require line overlaps for being detectable. For <sup>29</sup>SiO $`v=1`$ we find that 6 different overlaps of ro-vibrational lines of <sup>28</sup>SiO with <sup>29</sup>SiO tend to enhance the $`J=54`$ transition. Six transitions (given in bold face in Table 2: <sup>28</sup>SiO $`v=4,J=32`$, <sup>29</sup>SiO $`v=2,J=10,32`$ and $`v=3,J=21,43`$, and <sup>30</sup>SiO $`v=2,J=21`$) have not yet been searched for maser emission as far as we are aware. However, all other transitions in Table 2 have already been observed without success by several authors (e.g., Cernicharo & Bujarrabal 1992). Nevertheless, further investigations of these lines would be useful, especially as <sup>29</sup>SiO $`v=1,J=21`$ was recently detected in one star (J.R. Pardo, private communication).
Finally, we stress that other types of line overlaps may simultaneously play a role in the anomalous excitation of <sup>28</sup>SiO , <sup>29</sup>SiO and <sup>30</sup>SiO . The infrared lines of water are especially important in this context because water is an abundant species in the late-type stars where SiO is present. Apart from the H<sub>2</sub>O-SiO line overlap discussed by Olofsson et al.(1981, 1985) and Bujarrabal et al.(1996), there are several other infrared transitions of ortho- or para-H<sub>2</sub>O close to ro-vibrational transitions of SiO and isotopes which would deserve modelisation.
## 4 Conclusion
We have presented results for a new and simple model of SiO masers which incorporates both radiative and collisional pumpings as well as line overlap effects. Radiative pumping includes stellar radiation and circumstellar grains although their extinction coefficient is not well known. Collisional pumping involves molecular hydrogen as well as atomic hydrogen (under the form of simply guessed H-SiO collision rates). Our model involves 40 rotational levels for each of the first five vibrational states and derives the populations and opacities for the three species <sup>28</sup>SiO , <sup>29</sup>SiO and <sup>30</sup>SiO , and for a variety of circumstellar conditions. Several low and high $`J`$ rotational transitions are easily inverted in the vibrational states $`v=1,2,3`$ and 4 as observed in many late-type stars. Prediction of relative line intensities must account for line overlap effects which have been shown to play an important role in the excitation of SiO masers.
Our model incorporates in a very simple way the effects of line overlaps between ro-vibrational lines of <sup>28</sup>SiO , <sup>29</sup>SiO and <sup>30</sup>SiO . Three different overlaps may be treated simultaneously, and we have examined the impact on the level populations and opacities of 31 overlapping line pairs. We have found that line overlaps are important to enhance <sup>28</sup>SiO $`v=3,J=10`$ and $`v=4,J=54`$ maser emission or to turn off $`v=3,J=21`$ line emission. Various overlaps may lead to weak <sup>29</sup>SiO $`v=0,J=54`$ line emission and strong <sup>29</sup>SiO $`v=1`$ and 2, $`J=65`$ emission; these lines had not been explained in other works. More generally, several <sup>29</sup>SiO $`v=1`$ and 2 maser lines are well explained with line overlaps. We have also suggested possible explanations to several <sup>30</sup>SiO maser lines including the newly discovered transition $`v=1,J=10`$. In addition, we have predicted that several new maser transitions for <sup>28</sup>SiO ($`v=4`$), <sup>29</sup>SiO ($`v=0`$ to 3) and <sup>30</sup>SiO ($`v=1`$ and 2) could result from line overlap effects.
###### Acknowledgements.
We thank the referee, Dr. W.H. Kegel, for useful comments on the manuscript.
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# Quintessence and CMB
## Abstract
A particular kind of quintessence is considered, with equation of motion $`p_Q/\rho _Q=1`$, corresponding to a cosmological term with time-dependence $`\mathrm{\Lambda }(t)=\mathrm{\Lambda }(t_0)(R(t_0)/R(t))^P`$ which we examine initially for $`0P<3`$. Energy conservation is imposed, as is consistency with big-bang nucleosynthesis, and the range of allowed $`P`$ is thereby much restricted to $`0P<0.2`$. The position of the first Doppler peak is computed analytically and the result combined with analysis of high-Z supernovae to find how values of $`\mathrm{\Omega }_m`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ depend on $`P`$.
preprint: IFP-784-UNC astro-ph/0005406 May 2000
Our knowledge of the universe has changed dramatically even in the last five years. Five years ago the best guess, inspired partially by inflation, for the makeup of the present cosmological energy density was $`\mathrm{\Omega }_m=1`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$. However, the recent experimental data on the cosmic background radiation and the high - $`Z`$ ($`Z`$ = red shift) supernovae strongly suggest that both guesses were wrong. Firstly $`\mathrm{\Omega }_m0.3\pm 0.1`$. Second, and more surprisingly, $`\mathrm{\Omega }_\mathrm{\Lambda }0.7\pm 0.2`$. The value of $`\mathrm{\Omega }_\mathrm{\Lambda }`$ is especially unexpected for two reasons: it is non-zero and it is $`120`$ orders of magnitude below its “natural” value.
The fact that the present values of $`\mathrm{\Omega }_m`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ are of comparable order of magnitude is a “cosmic coincidence” if $`\mathrm{\Lambda }`$ in the Einstein equation
$`R_{\mu \nu }{\displaystyle \frac{1}{2}}g_{\mu \nu }R=8\pi G_NT_{\mu \nu }+\mathrm{\Lambda }g_{\mu \nu }`$
is constant. Extrapolate the present values of $`\mathrm{\Omega }_m`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ back, say, to redshift $`Z=100`$. Suppose for simplicity that the universe is flat $`\mathrm{\Omega }_C=0`$ and that the present cosmic parameter values are $`\mathrm{\Omega }_m=0.300\mathrm{}`$ exactly and $`\mathrm{\Omega }_\mathrm{\Lambda }=0.700\mathrm{}`$ exactly. Then since $`\rho _mR(t)^3`$ (we can safely neglect radiation), we find that $`\mathrm{\Omega }_m0.9999..`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }0.0000..`$ at $`Z=100`$. At earlier times the ratio $`\mathrm{\Omega }_\mathrm{\Lambda }/\mathrm{\Omega }_m`$ becomes infinitesimal. There is nothing to exclude these values but it does introduce a second “flatness” problem because, although we can argue for $`\mathrm{\Omega }_m+\mathrm{\Omega }_\mathrm{\Lambda }=1`$ from inflation, the comparability of the present values of $`\mathrm{\Omega }_m`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ cries out for explanation.
In the present paper we shall consider a specific model of quintessence. In its context we shall investigate the position of the first Doppler peak in the Cosmic Microwave Background (CMB) analysis using results published by two of us with Rohm earlier. Other works on the study of CMB include. We shall explain some subtleties of the derivation given in that have been raised since its publication mainly because the formula works far better than its expected order-of-magnitude accuracy. Data on the CMB have been provided recently in and especially in .
The combination of the information about the first Doppler peak and the complementary analysis of the deceleration parameter derived from observations of the high-red-shift supernovae leads to fairly precise values for the cosmic parameters $`\mathrm{\Omega }_m`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$. We shall therefore also investigate the effect of quintessence on the values of these parameters.
In , by studying the geodesics in the post-recombination period a formula was arrived at for the position of the first Doppler peak, $`l_1`$. For example, in the case of a flat universe with $`\mathrm{\Omega }_C=0`$ and $`\mathrm{\Omega }_M+\mathrm{\Omega }_\mathrm{\Lambda }=1`$ and for a conventional cosmological constant:
$$l_1=\pi \left(\frac{R_t}{R_0}\right)\left[\mathrm{\Omega }_M\left(\frac{R_0}{R_t}\right)^3+\mathrm{\Omega }_\mathrm{\Lambda }\right]^{1/2}_1^{\frac{R_0}{R_t}}\frac{dw}{\sqrt{\mathrm{\Omega }_Mw^3+\mathrm{\Omega }_\mathrm{\Lambda }}}$$
(1)
If $`\mathrm{\Omega }_C<0`$ the formula becomes
$$l_1=\frac{\pi }{\sqrt{\mathrm{\Omega }_C}}\left(\frac{R_t}{R_0}\right)\left[\mathrm{\Omega }_M\left(\frac{R_0}{R_t}\right)^3+\mathrm{\Omega }_\mathrm{\Lambda }+\mathrm{\Omega }_C\left(\frac{R_0}{R_t}\right)^2\right]^{1/2}\mathrm{sin}\left(\sqrt{\mathrm{\Omega }_C}_1^{\frac{R_0}{R_t}}\frac{dw}{\sqrt{\mathrm{\Omega }_Mw^3+\mathrm{\Omega }_\mathrm{\Lambda }}}\right)$$
(2)
For the third possibility of a closed universe with $`\mathrm{\Omega }_C>0`$ the formula is:
$$l_1=\frac{\pi }{\sqrt{\mathrm{\Omega }_C}}\left(\frac{R_t}{R_0}\right)\left[\mathrm{\Omega }_M\left(\frac{R_0}{R_t}\right)^3+\mathrm{\Omega }_\mathrm{\Lambda }+5\mathrm{\Omega }_C\left(\frac{R_0}{R_t}\right)^2\right]^{1/2}\mathrm{sinh}\left(\sqrt{\mathrm{\Omega }_C}_1^{\frac{R_0}{R_t}}\frac{dw}{\sqrt{\mathrm{\Omega }_Mw^3+\mathrm{\Omega }_\mathrm{\Lambda }}}\right)$$
(3)
The use of these formulas gives iso-$`l_1`$ lines on a $`\mathrm{\Omega }_M\mathrm{\Omega }_\mathrm{\Lambda }`$ plot in $`2550`$% agreement with the corresponding results found from computer code. On the insensitivity of $`l_1`$ to other variables, see. The derivation of these formulas was given in . Here we add some more details.
The formula for $`l_1`$ was derived from the relation $`l_1=\pi /\mathrm{\Delta }\theta `$ where $`\mathrm{\Delta }\theta `$ is the angle subtended by the horizon at the end of the recombination transition. Let us consider the Legendre integral transform which has as integrand a product of two factors, one is the temperature autocorrelation function of the cosmic background radiation and the other factor is a Legendre polynomial of degree $`l`$. The issue is what is the lowest integer $`l`$ for which the two factors reinforce to create the doppler peak? For small $`l`$ there is no reinforcement because the horizon at recombination subtends a small angle about one degree and the CBR fluctuations average to zero in the integral of the Legendre transform. At large $`l`$ the Legendre polynomial itself fluctuates with almost equispaced nodes and antinodes. The node-antinode spacing over which the Legendre polynomial varies from zero to a local maximum in magnitude is, in terms of angle, on average $`\pi `$ divided by $`l`$. When this angle coincides with the angle subtended by the last-scattering horizon, the fluctuations of the two integrand factors are, for the first time with increasing $`l`$, synchronized and reinforce (constructive interference) and the corresponding partial wave coefficient is larger than for slightly smaller or slightly larger $`l`$. This explains the occurrence of $`\pi `$ in the equation for the $`l_1`$ value of the first doppler peak written as $`l_1=\pi /\mathrm{\Delta }\theta `$.
Another detail concerns the use of the photon horizon as opposed to the acoustic horizon. If we examine the evolution of the recombination transition given in the degree of ionization is 99% at $`5,000^0`$K (redshift $`Z=1,850`$) falling to 1% at $`3,000^0`$K ($`Z=1,100`$). These times represent the beginning and end of the recombination transition. For matter domination $`Rt^{2/3}`$ so in cosmic time the start of recombination is at $`t=1.7\times 10^5`$y ($`Z=1,850`$) and the end is at $`t=3.8\times 10^5`$y ($`Z=1,100`$). The photons we detect are from $`Z=1,100`$. The baryon-photon plasma has fluctuations at $`Z=1,850`$ with size the acoustic horizon $`1.7\times 10^5y/\sqrt{3}=1.0\times 10^5y`$. Between $`Z=1,850`$ and $`Z=1,100`$ the transparency increases as does the size of the fluctuation. The sound speed is gradually replaced by the light speed in the fluctuation evoluton. One can see quantitatively that during the cosmic time $`2\times 10^5`$y of the recombination transition the fluctuation can grow by the required amount if one uses a speed between sound and light. The agreement of the formula for $`l_1`$ with experiment is at the 10% level which shows phenomenologcally that the fluctuation does grow (approximately by $`\sqrt{3}`$) during the recombination transition and that is why there is no $`\sqrt{3}`$ in its numerator. Although the formula started out as an order-of-magnitude estimate the fact that it works far better gives insight about the physics of recombination and cosmic background radiation.
To introduce our quintessence model as a time-dependent cosmological term, we start from the Einstein equation:
$$R_{\mu \nu }\frac{1}{2}Rg_{\mu \nu }=\mathrm{\Lambda }(t)g_{\mu \nu }+8\pi GT_{\mu \nu }=8\pi G𝒯_{\mu \nu }$$
(4)
where $`\mathrm{\Lambda }(t)`$ depends on time as will be specified later and $`T_\nu ^\mu =\mathrm{diag}(\rho ,p,p,p)`$. Using the Robertson-Walker metric, the ‘00’ component of Eq.(4) is
$$\left(\frac{\dot{R}}{R}\right)^2+\frac{k}{R^2}=\frac{8\pi G\rho }{3}+\frac{1}{3}\mathrm{\Lambda }$$
(5)
while the ‘ii’ component is
$$2\frac{\ddot{R}}{R}+\frac{\dot{R}^2}{R^2}+\frac{k}{R^2}=8\pi Gp+\mathrm{\Lambda }$$
(6)
Energy-momentum conservation follows from Eqs.(5,6) because of the Bianchi identity $`D^\mu (R_{\mu \nu }\frac{1}{2}g_{\mu \nu })=D^\mu (\mathrm{\Lambda }g_{\mu \nu }+8\pi GT_{\mu \nu })=D^\mu 𝒯_{\mu \nu }=0`$.
Note that the separation of $`𝒯_{\mu \nu }`$ into two terms, one involving $`\mathrm{\Lambda }(t)`$, as in Eq(4), is not meaningful except in a phenomenological sense because of energy conservation.
In the present cosmic era, denoted by the subscript ‘0’, Eqs.(5,6) become respectively:
$$\frac{8\pi G}{3}\rho _0=H_0^2+\frac{k}{R_0^2}\frac{1}{3}\mathrm{\Lambda }_0$$
(7)
$$8\pi Gp_0=2q_0H_0^2+H_0^2+\frac{k}{R_0^2}\mathrm{\Lambda }_0$$
(8)
where we have used $`q_0=\frac{\ddot{R}_0}{R_0H_0^2}`$ and $`H_0=\frac{\dot{R}_0}{R_0}`$.
For the present era, $`p_0\rho _0`$ for cold matter and then Eq.(8) becomes:
$$q_0=\frac{1}{2}\mathrm{\Omega }_M\mathrm{\Omega }_\mathrm{\Lambda }$$
(9)
where $`\mathrm{\Omega }_M=\frac{8\pi G\rho _0}{3H_0^2}`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }=\frac{\mathrm{\Lambda }_0}{3H_0^2}`$.
Now we can introduce the form of $`\mathrm{\Lambda }(t)`$ we shall assume by writing
$$\mathrm{\Lambda }(t)=bR(t)^P$$
(10)
where $`b`$ is a constant and the exponent $`P`$ we shall study for the range $`0P<3`$. This motivates the introduction of the new variables
$$\stackrel{~}{\mathrm{\Omega }}_M=\mathrm{\Omega }_M\frac{P}{3P}\mathrm{\Omega }_\mathrm{\Lambda },\stackrel{~}{\mathrm{\Omega }}_\mathrm{\Lambda }=\frac{3}{3P}\mathrm{\Omega }_\mathrm{\Lambda }$$
(11)
It is unnecessary to redefine $`\mathrm{\Omega }_C`$ because $`\stackrel{~}{\mathrm{\Omega }}_M+\stackrel{~}{\mathrm{\Omega }}_\mathrm{\Lambda }=\mathrm{\Omega }_M+\mathrm{\Omega }_\mathrm{\Lambda }`$. The case $`P=2`$ was proposed, at least for late cosmological epochs, in .
The equations for the first Doppler peak incorporating the possibility of non-zero $`P`$ are found to be the following modifications of Eqs.(1,2,3). For $`\mathrm{\Omega }_C=0`$
$$l_1=\pi \left(\frac{R_t}{R_0}\right)\left[\stackrel{~}{\mathrm{\Omega }}_M\left(\frac{R_0}{R_t}\right)^3+\stackrel{~}{\mathrm{\Omega }}_\mathrm{\Lambda }\left(\frac{R_0}{R_t}\right)^P\right]^{1/2}_1^{\frac{R_0}{R_t}}\frac{dw}{\sqrt{\stackrel{~}{\mathrm{\Omega }}_Mw^3+\stackrel{~}{\mathrm{\Omega }}_\mathrm{\Lambda }w^P}}$$
(12)
If $`\mathrm{\Omega }_C<0`$ the formula becomes
$`l_1`$ $`=`$ $`{\displaystyle \frac{\pi }{\sqrt{\mathrm{\Omega }_C}}}\left({\displaystyle \frac{R_t}{R_0}}\right)[\stackrel{~}{\mathrm{\Omega }}_M\left({\displaystyle \frac{R_0}{R_t}}\right)^3+\stackrel{~}{\mathrm{\Omega }}_\mathrm{\Lambda }\left({\displaystyle \frac{R_0}{R_t}}\right)^P+\mathrm{\Omega }_C\left({\displaystyle \frac{R_0}{R_t}}\right)^2]^{1/2}\times `$ (14)
$`\times \mathrm{sin}\left(\sqrt{\mathrm{\Omega }_C}{\displaystyle _1^{\frac{R_0}{R_t}}}{\displaystyle \frac{dw}{\sqrt{\stackrel{~}{\mathrm{\Omega }}_Mw^3+\stackrel{~}{\mathrm{\Omega }}_\mathrm{\Lambda }w^P+\mathrm{\Omega }_Cw^2}}}\right)`$
For the third possibility of a closed universe with $`\mathrm{\Omega }_C>0`$ the formula is:
$`l_1`$ $`=`$ $`{\displaystyle \frac{\pi }{\sqrt{\mathrm{\Omega }_C}}}\left({\displaystyle \frac{R_t}{R_0}}\right)[\stackrel{~}{\mathrm{\Omega }}_M\left({\displaystyle \frac{R_0}{R_t}}\right)^3+\stackrel{~}{\mathrm{\Omega }}_\mathrm{\Lambda }\left({\displaystyle \frac{R_0}{R_t}}\right)^P+\mathrm{\Omega }_C\left({\displaystyle \frac{R_0}{R_t}}\right)^2]^{1/2}\times `$ (16)
$`\times \mathrm{sinh}\left(\sqrt{\mathrm{\Omega }_C}{\displaystyle _1^{\frac{R_0}{R_t}}}{\displaystyle \frac{dw}{\sqrt{\stackrel{~}{\mathrm{\Omega }}_Mw^3+\stackrel{~}{\mathrm{\Omega }}_\mathrm{\Lambda }w^P+\mathrm{\Omega }_Cw^2}}}\right)`$
The dependence of $`l_1`$ on $`P`$ is illustrated for constant $`\mathrm{\Omega }_M=0.3`$ in Fig. 1(a), and for the flat case $`\mathrm{\Omega }_C=0`$ in Fig. 1(b). For illustration we have varied $`0P<3`$ but as will become clear later in the paper (see Fig 3 below) only the much more restricted range $`0P<0.2`$ is possible for a fully consistent cosmology when one considers evolution since the nucleosynthesis era.
We have introduced $`P`$ as a parameter which is real and with $`0P<3`$. For $`P0`$ we regain the standard cosmological model. But now we must investigate other restrictions already necessary for $`P`$ before precision cosmological measurements restrict its range even further.
Only for certain $`P`$ is it possible to extrapolate the cosmology consistently for all $`0<w=(R_0/R)<\mathrm{}`$. For example, in the flat case $`\mathrm{\Omega }_C=0`$ which our universe seems to approximate, the formula for the expansion rate is
$$\frac{1}{H_0^2}\left(\frac{\dot{R}}{R}\right)^2=\stackrel{~}{\mathrm{\Omega }}_Mw^3+\stackrel{~}{\mathrm{\Omega }}_\mathrm{\Lambda }w^P$$
(17)
This is consistent as a cosmology only if the right-hand side has no zero for a real positive $`w=\widehat{w}`$. The root $`\widehat{w}`$ is
$$\widehat{w}=\left(\frac{3(1\mathrm{\Omega }_M)}{P3\mathrm{\Omega }_M}\right)^{\frac{1}{3P}}$$
(18)
If $`0<\mathrm{\Omega }_M<1`$, consistency requires that $`P<3\mathrm{\Omega }_M`$.
In the more general case of $`\mathrm{\Omega }_C0`$ the allowed regions of the $`\mathrm{\Omega }_M\mathrm{\Omega }_\mathrm{\Lambda }`$ plot for $`P=0,1,2`$ are displayed in Fig. 2.
We see from Eq.(18) that if we do violate $`P<3\mathrm{\Omega }_M`$ for the flat case then there is a $`\widehat{w}>0`$ where the cosmology undergoes a bounce, with $`\dot{R}=0`$ and $`\dot{R}`$ changing sign. This necessarily arises because of the imposition of $`D^\mu 𝒯_{\mu \nu }=0`$ for energy conservation. For this example it occurs in the past for $`\widehat{w}>1`$. The consistency of big bang cosmology back to the time of nucleosynthesis implies that our universe has not bounced for any $`1<\widehat{w}<10^9`$. It is also possible to construct cosmologies where the bounce occurs in the future! Rewriting Eq.(18) in terms of $`\mathrm{\Omega }_\mathrm{\Lambda }`$:
$$\widehat{w}=\left(\frac{3\mathrm{\Omega }_\mathrm{\Lambda }}{3\mathrm{\Omega }_\mathrm{\Lambda }(3P)}\right)^{\frac{1}{3P}}$$
(19)
If $`P<3`$, then any $`\mathrm{\Omega }_\mathrm{\Lambda }<0`$ will lead to a solution with $`0<\widehat{w}<1`$ corresponding to a bounce in the future. If $`P>3`$ the condition for a future bounce is $`\mathrm{\Omega }_\mathrm{\Lambda }<\left(\frac{P3}{3}\right)`$. What this means is that for the flat case $`\mathrm{\Omega }_C=0`$ with quintessence $`P>0`$ it is possible for the future cosmology to be qualitatively similar to a non-quintessence closed universe where $`\dot{R}=0`$ at a finite future time with a subsequent big crunch.
Another constraint on the cosmological model is provided by nucleosynthesis which requires that the rate of expansion for very large $`w`$ does not differ too much from that of the standard model.
The expansion rate for $`P=0`$ coincides for large $`w`$ with that of the standard model so it is sufficient to study the ratio:
$`(\dot{R}/R)_P^2/(\dot{R}/R)_{P=0}^2`$ $`\stackrel{w\mathrm{}}{}`$ $`(3\mathrm{\Omega }_MP)/((3P)\mathrm{\Omega }_M)`$ (20)
$`\stackrel{w\mathrm{}}{}`$ $`(4\mathrm{\Omega }_RP)/((4P)\mathrm{\Omega }_R)`$ (21)
where the first limit is for matter-domination and the second is for radiation-domination (the subscript R refers to radiation).
The overall change in the expansion rate at the BBN era is therefore
$`(\dot{R}/R)_P^2/(\dot{R}/R)_{P=0}^2\stackrel{w\mathrm{}}{}(3\mathrm{\Omega }_MP)/((3P)\mathrm{\Omega }_M)\times (4\mathrm{\Omega }_R^{trans}P)/((4P)\mathrm{\Omega }_R^{trans})`$ (22)
where the superscript ”trans” refers to the transition from radiation domination to matter domination. Putting in the values $`\mathrm{\Omega }_M=0.3`$ and $`\mathrm{\Omega }_R^{trans}=0.5`$ leads to $`P<0.2`$ in order that the acceleration rate at BBN be within 15% of its value in the standard model, equivalent to the contribution to the expansion rate at BBN of one chiral neutrino flavor.
Thus the constraints of avoiding a bounce ($`\dot{R}=0`$) in the past, and then requiring consistency with BBN leads to $`0<P<0.2`$.
We may now ask how this restricted range of $`P`$ can effect the extraction of cosmic parameters from observations. This demands an accuracy which has fortunately begun to be attained with the Boomerang data. If we choose $`l_1=197`$ and vary $`P`$ as $`P=0,0.05,0.10,0.15,0.20`$ we find in the enlarged view of Fig 3 that the variation in the parameters $`\mathrm{\Omega }_M`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ can be as large as $`\pm 3\%`$. To guide the eye we have added the line for deceleration parameter $`q_0=0.5`$ as suggested by . In the next decade, inspired by the success of Boomerang (the first paper of true precision cosmology) surely the sum $`(\mathrm{\Omega }_M+\mathrm{\Omega }_\mathrm{\Lambda })`$ will be examined at much better than $`\pm 1\%`$ accuracy, and so variation of the exponent of $`P`$ will provide a useful parametrization of the quintessence alternative to the standard cosmological model with constant $`\mathrm{\Lambda }`$.
Clearly, from the point of view of inflationary cosmology, the precise vanishing of $`\mathrm{\Omega }_C=0`$ is a crucial test and its confirmation will be facilitated by comparison models such as the present one.
Acknowledgments
We thank L.H. Ford and S. Glashow for useful discussions and S. Weinberg for provocative questions. This work was supported in part by the US Department of Energy under Grant No. DE-FG02-97ER-41036.
Figure 1.
Dependence of $`l_1`$ on $`P`$ for (a) fixed $`\mathrm{\Omega }_M=0.3`$; (b) fixed $`\mathrm{\Omega }_C=0`$.
Figure 2.
Regions of the $`\mathrm{\Omega }_M\mathrm{\Omega }_\mathrm{\Lambda }`$ plot where there is a future bounce (small dot lattice), no bounce (unshaded) and a past bounce (large dot lattice) for (a) $`P=0`$; (b) $`P=1`$; and (c) $`P=2`$.
Figure 3.
Enlarged view of $`\mathrm{\Omega }_M\mathrm{\Omega }_\mathrm{\Lambda }`$ plot to exhibit sensitivity to $`0P0.2`$.
Contours are (right to left) $`P=0,0.05,0.10,0.15,0.20`$.
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# Numerical Simulations in Cosmology II: Spatial and Velocity Biases
## 1 Introduction
The distribution of galaxies is likely biased with respect to the dark matter. Therefore, the galaxies can be used to probe the matter distribution only if we understand the bias. Although the problem of bias has been studied extensively in the past (e.g., Kaiser 1984; Davis et al., 1985; Dekel & Silk 1986), new data on high redshift clustering and the anticipation of coming measurements have recently generated substantial theoretical progress in the field. The breakthrough in analytical treatment of the bias was the paper by Mo & White (1996), who showed how bias can be predicted in the framework of the extended Press-Schechter approximation. More elaborate analytical treatment has been developed by Catelan et al. (1998ab), Porciani et al.(1998), and Sheth & Lemson (1998). Effects of nonlinearity and stochasticity were considered in Dekel & Lahav (1998) (see also Toruya & Suto (2000)).
Valuable results are produced by “hybrid” numerical methods in which low-resolution N-body simulations (typical resolution $`20`$kpc) are combined with semi-analytical models of galaxy formation (e.g. Diaferio et al., ; Benson et al., 1999; Somerville et al., 1999). Typically, results of these studies are very close to those obtained with brute-force approach of high-resolution ($`\genfrac{}{}{0pt}{}{_<}{^{}}2`$kpc) N-body simulations (e.g., Colín et al., 1999a; Ghigna et al., ). This agreement is quite remarkable because the methods are very different. It may indicate that the biases of galaxy-size objects are controlled by the random nature of clustering and merging of galaxies and by dynamical effects, which cause the merging, because those are the only common effects in those two approaches.
Direct N-body simulations can be used for studies of the biases only if they have very high mass and force resolution. Because of numerous numerical effects, halos in low-resolution simulations do not survive in dense environments of clusters and groups (e.g., Moore, Katz & Lake 1996; Tormen, Diaferio & Syer, 1998; Klypin et al., 1999). Estimates of the needed resolution are given in Klypin et al. (1999). Indeed, recent simulations, which have sufficient resolution have found hundreds of galaxy-size halos moving inside clusters (Ghigna et al., 1998; Colín et al., 1999a; Moore et al., 1999; Okamoto & Habe, 1999).
It is very difficult to make accurate and trustful predictions of luminosities for galaxies, which should be hosted by dark matter halos. Instead of luminosities or virial masses we suggest to use circular velocities $`V_c`$ for both numerical and observational data. For a real galaxy its luminosity tightly correlate with the the circular velocity. So, one has a good idea what is the circular velocity of the galaxy. Nevertheless, direct measurements of circular velocities of a large complete sample of galaxies are extremely important because it will provide a direct way of comparing theory and observations. This lecture is mostly based on results presented in Colín et al. (1999ab) and Kravtsov & Klypin (1999).
## 2 Oh, Bias, Bias
There are numerous aspects and notions related with the bias. One should be really careful to understand what what type of bias is used. Results can be dramatically different. We start with introducing the overdensity field. If $`\overline{\rho }`$ is the mean density of some component (e.g., the dark matter or halos), then for each point $`𝐱`$ in space we have $`\delta (𝐱)[\rho (𝐱)\overline{\rho }]/\overline{\rho }`$. The overdensity can be decomposed into the Fourier spectrum, for which we can find the power spectrum $`P(k)=|\delta _𝐤|^2`$. We then can find the correlation function $`\xi (r)`$ and the rms fluctuation of $`\delta (R)`$ smoothed on a given scale $`R`$. We can construct the statistics for each component: dark matter, galaxies, or halos with given properties. Each statistics gives its own definition of bias $`b`$:
$$P_h(k)=b_P^2P_h(k),\xi _h(r)=b_\xi ^2\xi _{\mathrm{dm}}(r),\delta _h(R)=b_\delta \delta _{\mathrm{dm}}(R).$$
(1)
The three estimates of the bias $`b`$ are related. In special case, when the bias is linear, local, and scale independent all three forms of bias are all equal. In general case they are different and they are complicated nonlinear functions of scale, mass of the halos or galaxies, and redshift. The dependence on the scale is not local in the sense that the bias in a given position in space may depend on environment (e.g., density and velocity dispersion) on a larger scale. Bias has memory: it depends on the local history of fluctuations. There is another complication: bias very likely is not a deterministic function. One source of this stochasticity is that it is nonlocal. Dependence on the history of clustering may also introduce some random effect.
There are some processes, which we know create and affect the bias. At high redshifts there is statistical bias: in a Gaussian correlated field high density regions are more clustered than the field itself (Kaiser, 1984). Mo & White (1996) showed how the extended Press-Schechter formalism can be used for derivation of the bias of the dark matter halos. In the limit of small perturbations on large scales the bias is (Catelan et al., 1998a; Toruya & Suto, 2000)
$$b(M,z,z_f)=1+\frac{\nu ^21}{\delta _c(z,z_f)}.$$
(2)
Here $`\nu =\delta _c(z,z_f)/\sigma (M,z)`$ is the relative amplitude of a fluctuation on scale $`M`$ in units of the rms fluctuation $`\sigma (M,z)`$ of the density field at redshift $`z`$. Parameter $`z_f`$ is the redshift of halo formation. The critical threshold of the top-hat model is $`\delta _c(z,z_f)=\delta _{c,0}D(z)/D(z_f)`$, where $`D`$ is the growth factor of perturbations and $`\delta _{c,0}=1.69`$. At high redshifts, parameter $`\nu `$ for galaxy-size fluctuations is very large and $`\delta _c`$ is small. As the result, galaxy-size halos are expected to be more clustered (strongly biases) as compared to the dark matter. The bias is larger for more massive objects. As fluctuations grow, new forming galaxy-size halos do not represent as high peaks as at large redshifts and the bias tends to decrease. It also looses its sensitivity.
At later stages another process start to change the bias. In groups and cluster progenitors the merging and destruction of halos reduces the number of halos. This does not happen in the field where number of halos of given mass may only increase with time. As the result, the number of halos inside groups and cluster progenitors is reduced relatively to the field. This produces (anti)bias: there is relatively smaller number of halos as compared with the dark matter. This merging bias does not depend on mass of halos and it has a tendency to slow down once a group becomes a cluster with large relative velocity of halos (Kravtsov & Klypin, 1999).
Here is a list of different types of biases. We classify them into three groups: (1) measures of bias (2) terms related with the description of biases, (3) physical precesses, which produce or change the bias.
* Measures of bias
1. bias measured in a statistical sense (e.g., ratio of correlation functions $`\xi _h(r)=b^2\xi _{\mathrm{dm}}(r)`$)
2. bias measured point-by-point (e.g., $`\delta _h(𝐱)\delta _m(𝐱)`$ diagrams)
* Description of biases
1. local and nonlocal bias. For example, $`b(R)=\sigma _h(R)/\sigma _m(R)`$ is the local bias. If $`b=b(R;\stackrel{~}{R})`$, the bias is nonlocal, where $`\stackrel{~}{R}`$ is some other scale or scales.
2. linear and nonlinear bias. If in $`\xi _h(r)=b^2\xi _{\mathrm{dm}}(r)`$ the bias $`b`$ does not depend on $`\xi _{\mathrm{dm}}`$, it is the linear bias.
3. scale dependent and scale independent bias. If $`b`$ does not depend on scale at which the bias is estimated, the bias is scale independent. Note that in general, the bias can be nonlinear and scale independent, but this highely nonlikely.
4. stochastic and deterministic.
* Physical precesses, which produce or change the bias
1. statistical bias. Bias, which arises when a specific subset of points is selected from a Gaussian field.
2. merging bias. Bias produced due to merging and destruction of halos.
3. physical bias. Any bias due to physical processes inside forming galaxies.
## 3 Spatial bias
Colín et al. (1999a) have simulated different cosmological models and using the simulations studied halo biases. Most of the results presented here are for currently favored $`\mathrm{\Lambda }`$CDM model with the following parameters: $`\mathrm{\Omega }_0=1\mathrm{\Omega }_\mathrm{\Lambda }=0.3`$, $`h=0.7`$, $`\mathrm{\Omega }_b=0.032`$, $`\sigma _8=1`$. The model was simulated with $`256^3`$ particles in a 60$`h^1`$Mpc box. Formal mass and force resolutions are $`m_1=1.1\times 10^9h^1M_{}`$ and $`2h^1`$kpc. Bound Density Maximum halo finder was used to identify halos with at least 30 bound particles. For each halo we find maximum circular velocity $`V_c=\sqrt{GM(<r)/r}`$.
In figure 1 we compare the evolution of the correlation functions of the dark matter and halos. There are remarkable differences between halos and the dark matter. The correlation functions of the dark matter always increases with time (but the rate is different on different scales) and it never is a power-law. The correlation functions of the halos at redshifts goes down and then starts to increase again. It is accurately described by a power-law with slope $`\gamma =(1.51.7)`$. Figure 2 presents a comparison of the theoretical and observational data on correlation functions and power spectra. The dark matter clearly predicts much too high amplitude of clustering. The halos are much closer to the observational points and predict antibias. For the correlation function the antibias appears on scales $`r<5h^1`$Mpc; for the power spectrum the scales are $`k>0.2h\mathrm{Mpc}^1`$. One may get an impression that the antibias starts at longer waves in the power spectrum $`\lambda =2\pi /k30h^1\mathrm{Mpc}`$ as compared with $`r5h^1`$Mpc in the correlation function. There is no contradiction: sharp bias at small distances in the correlation function when Fourier transformed to the power spectrum produces antibias at very small wavenumbers. Thus, the bias should be taken into account at long waves when dealing with the power spectra. There is an inflection point in the power spectrum where the nonlinear power spectrum start to go upward (if one moves from low to high $`k`$) as compared with the prediction of the linear theory. Exact position of this point may have been affected by the finite size of the simulation box $`k_{\mathrm{min}}=0.105h^1`$Mpc, but effect is expected to be small.
At $`z=0`$ the bias almost does not depend on the mass limit of the halos. There is a tendency of more massive halos to be more clustered at very small distances $`r<200h^1`$kpc, but at this stage it is not clear that this is not due to residual numerical effects around centers of clusters. The situation is different at high redshift. At very high redshifts $`z>3`$ galaxy-size halos are very strongly (positively) biased. For example, at $`z=5`$ the correlation function of halos with $`v_c>150\mathrm{k}\mathrm{m}/\mathrm{s}`$ was 15 times larger than that of the dark matter at $`r=0.5h^1`$Mpc (see Fig.8 in Colín et al. (1999a)). The bias was also very strongly mass-dependent with more massive halos being more clustered. At smaller redshifts the bias was quickly declining. Around $`z=12`$ (exact value depends on halo circular velocity) the bias crossed unity and became less than unity (antibias) at later redshifts.
Evolution of bias is illustrated by Figure 4. The figure shows that at all epochs the overdensity of halos tightly correlates with the overdensity of the dark matter. The slope of the relation depends on the dark matter density and evolves with time. At $`z>1`$ halos are biased ($`\delta _h>\delta _m`$) in overdense regions with $`\delta _m>1`$ and antibiased in underdense regions with $`\delta _m<0.5`$ At low redshifts there is an antibias at large overdensities and almost no bias at low densities.
Figure 5 shows the density profiles for a cluster with mass $`2.5\times 10^{14}h^1M_{}`$. There is antibias on scales below $`300h^1\mathrm{kpc}`$. This is an example of the merging and destruction bias. Some of the halos have merged or were destroyed by the central cD halo of the cluster. As the result, there is smaller number of halos in the central part as compared with what we would expect if the number density of halos follower the density of the dark matter (the full curve in the bottom panel). Note that in the outer parts of the cluster halos closely follow the dark matter.
## 4 Velocity bias
There are two statistics, which measure velocity biases – differences in velocities of the galaxies (halos) and the dark matter. For a review of results and references see Colín et al. (1999). Two-particle or pairwise velocity bias (PVB) measures the relative velocity dispersion in pairs of objects with given separation $`r`$: $`b_{12}=\sigma _{\mathrm{halo}\mathrm{halo}}(r)/\sigma _{\mathrm{dm}\mathrm{dm}}(r)`$. Figure 6 (left panel) shows this bias. It is very sensitive to the number of pairs inside clusters of galaxies, where relative velocities are largest. Removal of few pairs can substantially change the value of the bias. This “removal” happens when halos merge or are destroyed by central cluster halos.
One-point velocity bias is estimated as a ratio of the rms velocity of halos to that of the dark matter: $`b_1=\sigma _{halos}/\sigma _{dm}`$. It is typically applied to clusters of galaxies where it is measured at different distances from the cluster center. For analysis of the velocity bias in clusters Colín et al. (1999a) have selected twelve most massive clusters in a simulation of the $`\mathrm{\Lambda }`$CDM model. The most massive cluster had virial mass $`6.5\times 10^{14}h^1M_{}`$ comparable to that of the Coma cluster. The cluster had 246 halos with circular velocities larger than 90 km/s. There were three Virgo-type clusters with virial masses in the range $`(1.62.4)\times 10^{14}h^1M_{}`$ and with approximately 100 halos in each cluster. Just as the spatial bias, PVB is positive at large redshifts (except for the very small scales) and decreases with the redshift. At lower redshifts it does not evolve much and stays below unity (antibias) at scales below $`5h^1`$Mpc on the level $`b_{12}(0.60.8)`$.
Figure 6 shows one-point velocity bias in clusters at $`z=0`$. Note that the sign of the bias is now different: halos move slightly faster than the dark matter. The bias is stronger in the central parts $`b_1=1.21.3`$ and goes to almost no bias $`b_11`$ at the virial radius and above. Both the antibias in the pairwise velocities and positive one-point bias are produced by the same physical process – merging and destruction of halos in central parts of groups and clusters. The difference is in the different weighting of halos in those two statistics. Smaller number of high-velocity pairs significantly changes PVB, but it does not affect much the one-point bias because it is normalized to the number of halos at a given distance from the cluster center. At the same time, merging preferentially happens for halos, which move with smaller velocity at a given distance from the cluster center. Slower halos have shorter dynamical times and have smaller apocenters. Thus, they have better chance to be destroyed and merged with the central cD halo. Because low-velocity halos are eaten up by the central cD, velocity dispersion of those, which survive, is larger. Another way of addressing the issue of the velocity bias is to use the Jeans equations. If we have a tracer population, which is in equilibrium in a potential produced by mass $`M(<r)`$, then
$$r\sigma _r^2(r)\left[\frac{d\mathrm{ln}\sigma _r^2(r)}{d\mathrm{ln}r}+\frac{d\mathrm{ln}\rho (r)}{d\mathrm{ln}r}+2\beta (r)\right]=GM(<r),$$
(3)
where $`\rho `$ is the number density of the tracer, $`\beta `$ is the velocity anisotropy, and $`\sigma _r`$ is the rms radial velocity. The r.h.s. of the equation is the same for the dark matter and the halos. If the term in the brackets would be the same, there would be no velocity bias. But there is systematic difference between the halos and the dark matter: the slope of the distribution halos in a cluster $`\frac{d\mathrm{ln}\rho (r)}{d\mathrm{ln}r}`$ is smaller than that of the dark matter (see Colín et al., 1998, Ghigna et al., 1999). The reason for the difference of the slopes is the same – merging with the central cD. Other terms in the equation also have small differences, but the main contribution comes from the slope of the density. Thus, as long as we have spatial antibias of the halos, there should be a small positive one-point velocity bias in clusters and a very strong antibias in pairwise velocity. Exact values of the biases are still under debate, but one thing seems to be certain: one bias does not go without the other.
The velocity bias in clusters is difficult to measure because it is small. The Figure 6 may be misleading because it shows the average trend, but it does not give the level of fluctuations for a single cluster. Note that the errors in the plots correspond to the error of the mean obtained by averaging over 12 clusters and two close moments of time. Fluctuations for a single cluster are much larger. Figure 6 shows results for three Virgo-type clusters in the simulation. The noise is very large because of both poor statistics (small number of halos) and the noise produced by residual non-equilibrium effects (substructure). Comparable (but slightly smaller) value of $`b_v`$ was recently found in simulations by Ghigna et al. (1999, astro-ph/9910166) for a cluster in the same mass range as in Figure 6. Unfortunately, it is difficult to make detailed comparison with their results because Ghigna et al. (1999) use only one hand-picked cluster for a different cosmological model. Very likely their results are dominated by the noise due to residual substructure. Results of another high-resolution simulation by Okamoto & Habe (1999) are consistent with our results.
## 5 Conclusions
There is a number of physical processes, which can contribute to the biases. In our papers we explore dynamical effects in the dark matter itself, which result in differences of the spatial and velocity distribution of the halos and the dark matter. Other effects related to the formation of luminous parts of galaxies also can produce or change biases. At this stage it is not clear how strong are those biases. Because there is a tight correlation between the luminosity and circular velocity of galaxies, any additional biases are limited by the fact that galaxies “know” how much dark matter they have.
Biases in the halos are reasonably well understood and can be approximated on a few Megaparsec scales by analytical models. We find that the biases in the distribution of the halos are sufficient to explain within the framework of standard cosmological models the clustering properties of galaxies on a vast ranges of scales from 100 kpc to dozens Megaparsecs. Thus, there is neither need nor much room for additional biases in the standard cosmological model.
In any case, biases in the halos should be treated as benchmarks for more complicated models, which include non-gravitational physics. If a model can not reproduce biases of halos or it does not have enough halos, it should be rejected, because it fails to have correct dynamics of the main component of the Universe – the dark matter.
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# Phys. Rev. D 63, 014017 (2001). Effective Lagrangian induced by the anomalous Wess-Zumino action and the exotic resonance state with 𝐼^𝐺(𝐽^{𝑃𝐶})=1⁻(1⁻⁺) in the 𝜌𝜋, 𝜂𝜋, 𝜂'𝜋, and 𝐾^∗𝐾̄+𝐾̄^∗𝐾 channels
## Abstract
A simple model for the exotic waves with $`I^G(J^{PC})=1^{}(1^+)`$ in the reactions $`VPVP`$, $`VPPP`$, and $`PPPP`$ is constructed beyond the scope of the quark-gluon approach. The model satisfies unitarity and analyticity and uses as a “priming” the “anomalous” nondiagonal $`VPPP`$ interaction which couples together the four channels $`\rho \pi `$, $`\eta \pi `$, $`\eta ^{}\pi `$, and $`K^{}\overline{K}+\overline{K}^{}K`$. The possibility of the resonancelike behavior of the $`I^G(J^{PC})=1^{}(1^+)`$ amplitudes belonging to the $`\{10\}\{\overline{10}\}`$ and $`\{8\}`$ representations of $`SU(3)`$ as well as their mixing is demonstrated explicitly in the $`1.31.6`$ GeV mass range which, according to the current experimental evidence, is really rich in exotics.
PACS number(s): 12.39.Fe, 13.25.Jx, 13.75.Lb
I. INTRODUCTION
Phantoms of manifestly exotic states with $`I^G(J^{PC})=1^{}(1^+)`$ have more and more agitated the experimental and theoretical communities \[1-20\]. They were discovered in the $`1.31.6`$ GeV mass range in the $`\eta \pi `$, $`\eta ^{}\pi `$, $`\rho \pi `$, $`b_1\pi `$, and $`f_1\pi `$ systems produced in $`\pi ^{}p`$ collisions at high energies and in $`N\overline{N}`$ annihilation at rest \[1-12\].
Theoretical considerations concerning the mass spectra and decay properties of exotic hadrons have been based, in the main, on the MIT-bag model, constituent gluon model, flux-tube model, QCD sum rules, lattice calculations, and various selection rules. The more recent discussions of these constituent quark-gluon models and selection rules in conformity with the observed $`J^{PC}=1^+`$ phenomena can be found in Refs. \[1-20\], together with extensive analyses of the current experimental data and comprehensive references. A resonance character of the observed exotic signals and also the more popular assumption about their hybrid $`(q\overline{q}g)`$ nature are the subject of much attention and require further careful investigations \[1,6,7,12,14-20\].
Let us note that evidence for the possible existence of an exotic $`J^{PC}=1^+`$ state coupled to the $`\eta \pi `$ and $`\rho \pi `$ channels and belonging to the icosuplet representation of $`SU(3)`$ was obtained for the first time by using the bootstrap technique of Schechter and Okubo more than 35 years ago (see also Ref. ).
Current algebra and effective chiral Lagragians are also important sources of theoretical information on exotic partial waves. It is sufficient to remember the prediction obtained within the framework of these approaches for the $`\pi \pi `$ $`s`$-wave scattering length with isospin $`I=2`$ . Constructing with the help of the effective chiral Lagragians the series expansions of the scattering amplitudes in powers of external momenta, one can reveal explicitly exotic contributions already among the lower order terms of these series. Can at least some of these contributions found at low energies turn out to be the manifestations (“the tails”) of high-lying exotic resonances? It is well known that, for example, for the $`\pi \pi `$ scattering channels involving the $`\sigma `$ and $`\rho `$ resonances, one can self-consistently (in the sense of agreement with experiment) sew together the resonancelike and low-energy behaviors of the scattering amplitudes by using the successfully selected unitarization scheme for the original chiral contributions, together with general analyticity requirements \[25-34\]. In other words, for these channels, there exist a good many of the model constructions which show that the low-energy contributions calculated within the effective chiral Lagragians framework may in principle transform with increasing energy into resonances with the experimentally established parameters. In the present work we continue in this way and construct a model satisfying unitarity and analyticity for an exotic wave with $`J^{PC}=1^+`$ in the reaction $`\rho \pi \eta \pi `$ and in the related reactions $`\rho \pi \eta ^{}\pi `$, $`\rho \pi \rho \pi `$, $`\eta \pi \eta \pi `$, $`\eta \pi (K^{}\overline{K}+\overline{K}^{}K)`$, and so on, using as a “priming” the tree exotic amplitudes generated by a simplest “anomalous” effective interaction of the vector ($`V`$) and pseudoscalar ($`P`$) mesons. The interaction is induced by the anomalous Wess-Zumino chiral Lagrangian and is proportional to $`ϵ_{\mu \nu \tau \kappa }`$.
At the tree level, the standard nonlinear chiral Lagrangian describing the low-energy dynamics of the pseudoscalar mesons belonging to the $`SU(3)`$ octet generates the $`PPPP`$ scattering amplitudes possessing only the usual quantum numbers $`J^{PC}=0^{++}`$ and $`1^{}`$ in the $`s`$ channel . However, already in the next order of chiral perturbation theory, the $`J^{PC}=1^+`$ exotic contributions arise in these amplitudes at the expense of the finite parts of the one-loop diagrams. In so doing they turn out to be different from zero only owing to $`SU(3)`$ symmetry breaking for pseudoscalar masses. The resonances, with which such contributions might be associated, have to possess rather odd properties. All their coupling constants to the octet of pseudoscalar mesons must vanish in the $`SU(3)`$ symmetry limit. Therefore, it seems more reasonable to assume that if the exotic resonances with $`J^{PC}=1^+`$ exist, then they are of another origin. In such a case, the resources for their possible generation, which still remain within the effective chiral Lagrangian framework, seem to involve the “anomalous” interactions of the vector and pseudoscalar mesons \[36-43\]. Some indirect evidence in favor of this assumption has been given by the analysis of the $`PPPP`$ scattering amplitudes carried out in the framework of the linear $`SU(3)\times SU(3)`$ $`\sigma `$ model involving only scalars and pseudoscalars . There operate the repulsive forces in the $`J^{PC}=1^+`$ channels in this model, and any resonance states do not arise.
In Sec. II, the general properties of the $`VPPP`$ reaction amplitudes are briefly discussed within the framework of the unitary symmetry assumption. In Sec. III, a simple model for the $`I=1`$ $`p`$-wave (exotic) reaction amplitudes $`VPPP`$, $`VPVP`$, and $`PPPP`$ is constructed. The model takes into account as a “priming” the nondiagonal $`VPPP`$ interaction which couples together the four channels $`\rho \pi `$, $`\eta \pi `$, $`\eta ^{}\pi `$, and $`K^{}\overline{K}+\overline{K}^{}K`$. It is essentially the summing up of all the $`s`$-channel loop diagrams with the $`VP`$ and $`PP`$ intermediate states. In Sec. IV, the possibility of the resonance-like behavior of the $`I^G(J^{PC})=1^{}(1^+)`$ amplitudes belonging to the $`\{10\}\{\overline{10}\}`$ and $`\{8\}`$ representations of $`SU(3)`$ as well as their mixing is demonstrated explicitly in the $`1.31.6`$ GeV mass range. In quark-gluon language, the $`\{10\}\{\overline{10}\}`$ representation of $`SU(3)`$ first occurs in the $`qq\overline{q}\overline{q}`$ sector, whereas the states with $`J^{PC}=1^+`$ belonging to the octet representation of $`SU(3)`$ may in principle correspond to both $`qq\overline{q}\overline{q}`$ and $`q\overline{q}g`$ configurations.
II. GENERAL PROPERTIES OF THE $`VP\mathbf{}PP`$ AMPLITUDE
The general Lorentz and $`SU(3)`$ structure for the amplitude of the reaction $`V_a(k)+P_b(q_1)P_c(q_2)+P_d(q_3)`$, where $`V_a`$ and $`P_a`$ are the members of the vector and pseudoscalar octets taken in the Cartesian basis ($`a=1,\mathrm{},8`$), <sup>1</sup><sup>1</sup>1$`V_a=(\rho _1,\rho _2,\rho _3,K_4^{},K_5^{},K_6^{},K_7^{},\omega _8)`$ and $`P_a=(\pi _1,\pi _2,\pi _3,K_4,K_5,K_6,K_7,\eta _8)`$. and $`k,q_1,q_2`$, and $`q_3`$ are the four-momenta of the particles in the reaction, has the form
$`M_{ab;cd}^{(\lambda )}=iϵ_{\mu \nu \tau \kappa }e_{(\lambda )}^\mu q_1^\nu q_2^\tau q_3^\kappa [f_{abm}d_{mcd}A(s,t,u)+d_{abm}f_{mcd}B(s,t,u)+(u_{ab})_{cd}C(s,t,u)].`$ (1)
Here $`f_{abc}`$ and $`d_{abc}`$ are the standard structure constants of $`SU(3)`$ , $`(u_{ab})_{cd}=f_{cam}d_{mbd}d_{cam}f_{mbd}`$ , $`e_{(\lambda )}^\mu `$ is a $`\mu `$ component of the $`V`$ meson polarization vector with helicity $`\lambda `$$`s=(k+q_1)^2`$$`t=(kq_2)^2`$, and $`u=(kq_3)^2`$. From Bose symmetry it follows that the invariant amplitude $`A(s,t,u)`$ is antisymmetric under the interchange of the $`t`$ and $`u`$ variables, whereas the invariant amplitudes $`B(s,t,u)`$ and $`C(s,t,u)`$ are symmetric. Note that Eq. (1) can be obtained in the usual way \[45-48\] by applying $`SU(3)`$ symmetry, together with $`P`$ and $`C`$ invariance.
The first and second terms in Eq. (1) correspond to the octet transition amplitudes $`\{8_a\}\{8_s\}`$ and $`\{8_s\}\{8_a\}`$ which we shall designate for short by $`A_{as}`$ and $`A_{sa}`$, respectively; as usual, $`\{8_s\}`$ and $`\{8_a\}`$ mean the symmetric and antisymmetric octet representations of $`SU(3)`$ which occur in the direct production of $`\{8\}\times \{8\}`$. The third term in Eq. (1) describes transitions via the mutually conjugate representations $`\{10\}`$ and $`\{\overline{10}\}`$ with the amplitudes $`A_{10}`$ and $`A_{\overline{10}}`$ appearing in the combination $`A_{10}A_{\overline{10}}`$. In other words, it describes the transitions from the initial $`VP`$ icosuplet $`\{10\}\{\overline{10}\}`$ to the final $`PP`$ icosuplet $`\{10\}+\{\overline{10}\}`$. The transition amplitudes between the self-conjugate representations $`A_{as}`$ and $`A_{sa}`$ can be expanded into the partial waves with $`J^{PC}=2^{++},4^{++},\mathrm{}`$ and $`J^{PC}=1^{},3^{},\mathrm{}`$, respectively. Hence they do not contain any explicitly exotic contributions. As for the $`\eta _8\pi `$ final state, it dos not occur in the $`\{8_a\}`$ but can belong to the representations $`\{8_s\}`$, $`\{10\}`$, and $`\{\overline{10}\}`$ . The $`SU(3)`$ exotic meson amplitudes $`A_{10}`$ and $`A_{\overline{10}}`$ can be expanded into partial waves with $`J^P=1^{},3^{},\mathrm{}`$ . The isotriplet amplitudes of $`A_{10}A_{\overline{10}}`$ correspond to two sets of the reactions with opposite $`G`$ parity in the $`s`$ channel: (a) $`\rho \eta _8\pi \pi `$, $`\rho \eta _8K\overline{K}`$, $`\omega _8\pi \pi \pi `$, $`\omega _8\pi K\overline{K}`$ and (b) $`\rho \pi \eta _8\pi `$, $`K^{}\overline{K}\eta _8\pi `$, $`\overline{K}^{}K\eta _8\pi `$. The partial amplitudes of the reactions belonging to set (a) possess the nonexotic quantum numbers $`I^G(J^{PC})=1^+(1^{},3^{},\mathrm{})`$ and, therefore, only hidden $`SU(3)`$ exotics. The reactions belonging to set (b) are purely exotic because they contain the partial waves with $`I^G(J^{PC})=1^{}(1^+,3^+,\mathrm{})`$. In particular, it is these reactions that will be the subject of our attention in the following.
Let us now write down the amplitude for the reaction $`V_a(k)+P_b(q_1)P_c(q_2)+P_d(q_3)`$ involving the $`V_0`$ and $`P_0`$ $`SU(3)`$ singlets:
$`N_{ab;cd}^{(\lambda )}=iϵ_{\mu \nu \tau \kappa }e_{(\lambda )}^\mu q_1^\nu q_2^\tau q_3^\kappa [\delta _{a0}f_{bcd}D(s,t,u)+\delta _{b0}f_{acd}E(s,t,u)+\delta _{c0}f_{abd}F(s,t,u)],`$ (2)
where $`a,b,c,d`$ are the flavor indices running now over $`0,1,\mathrm{},8`$, $`f_{ab0}=0`$, $`V_0=\omega _0`$, and $`P_0=\eta _0`$. By the isoscalar particles with the definite masses we shall mean the pseudoscalar mesons $`\eta =\eta _8\mathrm{cos}\theta _P\eta _0\mathrm{sin}\theta _P`$ and $`\eta ^{}=\eta _8\mathrm{sin}\theta _P+\eta _0\mathrm{cos}\theta _P`$ with the mixing angle $`\theta _P20^{}`$ and the vector mesons $`\omega =\sqrt{1/3}\omega _8+\sqrt{2/3}\omega _0`$ and $`\varphi =\sqrt{2/3}\omega _8\sqrt{1/3}\omega _0`$ with “ideal mixing”. Equation (2) describes the transitions via the $`SU(3)`$ octet intermediate states. The first two terms in Eq. (2) do not contribute to $`\eta \pi `$ and $`\eta ^{}\pi `$ production because they correspond to the transitions into the final states belonging to the $`\{8_a\}`$ representation which does not contain the $`\eta _8\pi `$ system. The third term in Eq. (2) describes $`\eta \pi `$ and $`\eta ^{}\pi `$ production via the $`SU(3)`$ singlet components of the $`\eta `$ and $`\eta ^{}`$. Under $`tu`$ interchange, the invariant amplitudes $`D(s,t,u)`$ and $`E(s,t,u)`$ are symmetric, while the invariant amplitude $`F(s,t,u)`$ does not possess a definite symmetry. Thus, the first two terms in Eq. (2) can be expanded into partial waves with $`J^{PC}=1^{},3^{},\mathrm{}`$ and the third term into partial waves with $`J^{PC}=1^+,2^{++},3^+,4^{++},\mathrm{}`$, of which the odd waves $`1^+,3^+,\mathrm{}`$ are exotic. In principle, Eqs. (1) and (2) permit the $`VPPP`$ reaction channels involving the $`\omega `$, $`\varphi `$, $`\eta `$, and $`\eta ^{}`$ mesons to be considered in the most general form.
III. MODEL FOR THE $`I^G(J^{PC})=1^{}(1^+)`$ WAVES IN THE REACTIONS $`VP\mathbf{}PP\mathbf{,}PP\mathbf{}PP`$, AND $`VP\mathbf{}VP`$
Consider the $`SU(3)`$ symmetric effective Lagrangian for the pointlike $`VPPP`$ interaction which also possesses additional nonet symmetry with respect to the vector mesons,
$`L(VPPP)=ihϵ_{\mu \nu \tau \kappa }\text{Tr}(\widehat{V}^\mu ^\nu \widehat{P}^\tau \widehat{P}^\kappa \widehat{P})+i\sqrt{1/3}h^{}ϵ_{\mu \nu \tau \kappa }\text{Tr}(\widehat{V}^\mu ^\nu \widehat{P}^\tau \widehat{P})^\kappa \eta _0,`$ (3)
where $`\widehat{P}=_{a=1}^8\lambda _aP_a/\sqrt{2}`$, $`\widehat{V}^\mu =_{a=0}^8\lambda _aV_a^\mu /\sqrt{2}`$, and $`\lambda _a`$ are the Gell-Mann matrices . The tree amplitudes for the reactions $`VPPP`$ generated by the Lagrangian (3) are given by Eqs. (1) and (2) with the following sets of the invariant amplitudes (here we omit their arguments for short): $`(A,B,C,D)=h(0,2,1,\sqrt{6})`$ and $`(E,F)=h^{}(\sqrt{2/3},\sqrt{2/3})`$. The presence of the amplitudes $`C`$ and $`F`$ such as above implies (see the discussion in Sec. II) that the Lagrangian (3) generates tree exotic amplitudes with $`I^G(J^{PC})=1^{}(1^+)`$ belonging to the $`\{10\}\{\overline{10}\}`$ representation of $`SU(3)`$ for the inelastic reactions $`\rho \pi \eta _8\pi `$, $`K^{}\overline{K}\eta _8\pi `$, and $`\overline{K}^{}K\eta _8\pi `$ as well as the amplitudes belonging to the octet representation of $`SU(3)`$ for the reaction $`\rho \pi \eta _0\pi `$, $`K^{}\overline{K}\eta _0\pi `$, and $`\overline{K}^{}K\eta _0\pi `$. In the next orders, these tree amplitudes induce as well the $`I^G(J^{PC})=1^{}(1^+)`$ exotic ones for the elastic processes $`\rho \pi \rho \pi `$, $`\eta \pi \eta \pi `$, and so on. In this connection it is of interest to consider the following $`4\times 4`$ system of scattering amplitudes for the coupled exotic channels of the reactions $`VPVP`$, $`VPPP`$ and $`PPPP`$:
$`T_{ij}=\left[\begin{array}{cccc}T(\rho \pi \rho \pi )\hfill & T(\rho \pi \eta \pi )\hfill & T(\rho \pi \eta ^{}\pi )\hfill & T(\rho \pi K^{}K)\hfill \\ T(\eta \pi \rho \pi )\hfill & T(\eta \pi \eta \pi )\hfill & T(\eta \pi \eta ^{}\pi )\hfill & T(\eta \pi K^{}K)\hfill \\ T(\eta ^{}\pi \rho \pi )\hfill & T(\eta ^{}\pi \eta \pi )\hfill & T(\eta ^{}\pi \eta ^{}\pi )\hfill & T(\eta ^{}\pi K^{}K)\hfill \\ T(K^{}K\rho \pi )\hfill & T(K^{}K\eta \pi )\hfill & T(K^{}K\eta ^{}\pi )\hfill & T(K^{}KK^{}K)\hfill \end{array}\right].`$ (8)
Here the subscripts $`i,j=1,2,3,4`$ are the labels of the $`\rho \pi `$, $`\eta \pi `$, $`\eta ^{}\pi `$, and $`K^{}K`$ channels, respectively, and the abbreviation $`K^{}K`$ implies just the $`\overline{K}^{}K`$ and $`K^{}\overline{K}`$ channels. The corresponding matrix of the $`VPPP`$ coupling constants generated by the Lagrangian (3) has the form
$`h_{ij}=h\left[\begin{array}{cccc}0& \alpha & \beta & 0\\ \alpha & 0& 0& \gamma \\ \beta & 0& 0& \delta \\ 0& \gamma & \delta & 0\end{array}\right],`$ (13)
where
$`\alpha =\sqrt{{\displaystyle \frac{1}{3}}}\mathrm{cos}\theta _P{\displaystyle \frac{h^{}}{h}}\sqrt{{\displaystyle \frac{2}{3}}}\mathrm{sin}\theta _P,\beta =\sqrt{{\displaystyle \frac{1}{3}}}\mathrm{sin}\theta _P+{\displaystyle \frac{h^{}}{h}}\sqrt{{\displaystyle \frac{2}{3}}}\mathrm{cos}\theta _P,`$ (14)
$`\gamma =\sqrt{{\displaystyle \frac{2}{3}}}\mathrm{cos}\theta _P+{\displaystyle \frac{h^{}}{h}}\sqrt{{\displaystyle \frac{1}{3}}}\mathrm{sin}\theta _P,\delta =\sqrt{{\displaystyle \frac{2}{3}}}\mathrm{sin}\theta _P{\displaystyle \frac{h^{}}{h}}\sqrt{{\displaystyle \frac{1}{3}}}\mathrm{cos}\theta _P.`$ (15)
In the following we shall consider three natural limiting cases: $`(i)`$ $`h^{}=0`$, i.e., when all exotic amplitudes belong to the $`\{10\}\{\overline{10}\}`$ representation of $`SU(3)`$, $`(ii)`$ $`h=0`$, i.e., when all exotic amplitudes belong to the octet representation of $`SU(3)`$, and $`(iii)`$ $`h^{}=h`$, when the original pointlike $`VPPP`$ interaction possesses nonet symmetry with respect to the pseudoscalar mesons.
To satisfy the unitarity condition for the coupled channel amplitudes, we sum up all the possible chains of the $`s`$-channel loop diagrams the typical examples of which are shown in Fig. 1. Such an old-fashioned field theory way of the unitarization is well known in the literature (see, for example, Refs. \[52-55,32\]). Notice that in case $`(i)`$ and in case $`(ii)`$ the whole complex of the unitarized amplitudes, in fact, can be constructed by using only the amplitudes for the loop diagrams shown explicitly in Fig. 1. The point is that in these cases the denominator of the corresponding geometrical series for any channel turns out to be proportional to the sum of diagrams $`(b)`$, $`(c)`$, $`(d)`$, and $`(e)`$ in Fig. 1, and the loop diagrams of the type $`(f)`$ and $`(g)`$, or $`(h)`$ and $`(i)`$, play a role of a “priming” in the corresponding elastic channels like diagram $`(a)`$ in the inelastic $`\rho \pi \eta \pi `$ channel. However, in case $`(iii)`$ the situation is considerably more complicated.
Before summing the diagrams, let us make two remarks about the model itself.
First, generally speaking, the pointlike exotic contributions due to the $`VPPP`$ interaction might be modified by the tree diagrams involving $`V`$ meson exchanges if one takes into account the “anomalous” Lagrangian for the $`VVP`$ interaction and the ordinary one for the $`VPP`$ interaction. However, such a considerable complication of the original exotic amplitudes actually does not lead to any new possibilities (or degrees of freedom) to obtain the resonancelike behavior of the complete unitarized amplitudes. This only burdens the model by additional technical difficulties and makes it much less transparent in comparison with the one based only on the Lagrangian (3). For example, after such a modification of the original exotic amplitudes, the above-mentioned obvious unitarization scheme need be changed, say, by some version of the Padé approximation because of the impossibility of the direct calculation and summation of higher loops.
Second, the effective coupling constant $`h`$ occurring in the Lagrangian (3) is not the unambiguously definite value in the theory with the “anomalous” chiral Lagrangians (comprehensive discussions of this point may be found in Refs. \[36-43\]). Actually, one may only claim that it is not too large in the scale defined by the combination $`2g_{\rho \pi \pi }g_{\omega \rho \pi }/m_\rho ^2284`$ GeV<sup>-3</sup> . Therefore, we consider the coupling constants $`h`$ and $`h^{}`$ as free parameters of the model in the region of their relatively small values.
The summing up of the loop diagram chains can be easily carried out by using the matrix equation for the auxiliary amplitudes $`\stackrel{~}{T}_{ij}`$,
$$\stackrel{~}{T}_{ij}=h_{ij}+h_{im}\mathrm{\Pi }_{mn}\stackrel{~}{T}_{nj},$$
(16)
which is shown schematically in Fig. 2. The auxiliary amplitudes $`\stackrel{~}{T}_{ij}`$ pertain to the hypothetical case when all the particles in the reactions are spinless, but otherwise they are the exact analogs of the physical amplitudes $`T_{ij}`$ designated in Eq. (4). So in Eq. (8) the matrix $`h_{ij}`$ is given by Eqs. (5), (6), and (7), and $`\mathrm{\Pi }_{ij}`$ is the diagonal matrix of the loops, i.e., $`\mathrm{\Pi }_{ij}=\delta _{ij}\mathrm{\Pi }_j`$, where $`\mathrm{\Pi }_1`$, $`\mathrm{\Pi }_2`$, $`\mathrm{\Pi }_3`$, and $`\mathrm{\Pi }_4`$ correspond to the four independent $`s`$-channel loops involving the $`\rho \pi `$, $`\eta \pi `$, $`\eta ^{}\pi `$, and $`K^{}K`$ intermediate states, respectively (for a moment, all for the spinless case). Notice that if all the particles are spinless, then the $`h_{ij}`$ and $`\mathrm{\Pi }_{ij}`$ in Eq. (8) are dimensionless, as well as the $`\stackrel{~}{T}_{ij}`$ themselves. For the matrix elements $`h_{im}\mathrm{\Pi }_{mj}=h_{ij}\mathrm{\Pi }_j`$ it is convenient to introduce the following compact notation \[look at Eq. (5)\]:
$`h_{im}\mathrm{\Pi }_{mj}=h_{ij}\mathrm{\Pi }_j=h\left[\begin{array}{cccc}0& \alpha _2& \beta _3& 0\\ \alpha _1& 0& 0& \gamma _4\\ \beta _1& 0& 0& \delta _4\\ 0& \gamma _2& \delta _3& 0\end{array}\right],`$ (21)
where $`\alpha _1=\alpha \mathrm{\Pi }_1`$, $`\beta _3=\beta \mathrm{\Pi }_3`$, and so on. The solution of Eq. (8) has the form
$$\stackrel{~}{T}_{ij}=[(\widehat{1}\widehat{h}\widehat{\mathrm{\Pi }})^1]_{im}h_{mj},$$
(22)
where $`\widehat{1}`$ is the $`4\times 4`$ identity matrix and the matrix $`\widehat{h}\widehat{\mathrm{\Pi }}`$ is defined by the relations of Eq. (9). Next, we define
$$\overline{D}=\text{det}(\widehat{1}\widehat{h}\widehat{\mathrm{\Pi }})=1h^2(\alpha _1\alpha _2+\beta _1\beta _3+\gamma _2\gamma _4+\delta _3\delta _4)+h^4(\alpha _1\delta _4\beta _1\gamma _4)(\alpha _2\delta _3\beta _3\gamma _2).$$
(23)
Let us now write down, as an example, the explicit expressions for the amplitudes of the following five reactions:
$$\stackrel{~}{T}(\rho \pi \rho \pi )=h^2[\alpha \alpha _2+\beta \beta _3h^2(\alpha \delta _4\beta \gamma _4)(\alpha _2\delta _3\beta _3\gamma _2)]/\overline{D},$$
(24)
$$\stackrel{~}{T}(\rho \pi \eta \pi )=h[\alpha h^2(\alpha \delta _3\beta _3\gamma )\delta _4]/\overline{D},$$
(25)
$$\stackrel{~}{T}(\rho \pi \eta ^{}\pi )=h[\beta +h^2(\alpha _2\delta \beta \gamma _2)\gamma _4]/\overline{D},$$
(26)
$$\stackrel{~}{T}(\rho \pi K^{}K)=h^2[\alpha _2\gamma +\beta _3\delta ]/\overline{D}.$$
(27)
$$\stackrel{~}{T}(\eta \pi \eta \pi )=h^2[\alpha \alpha _1+\gamma \gamma _4h^2(\alpha \delta _3\beta _3\gamma )(\alpha _1\delta _4\beta _1\gamma _4)]/\overline{D},$$
(28)
In cases $`(i)`$ and $`(ii)`$, the combination $`h^2(\alpha \delta \beta \gamma )=0`$ \[see Eqs. (5), (6), and (7)\], so that the contributions proportional to that vanish in Eqs. $`(10)(16)`$ and, as one can see, all the formulas are essentially simplified \[for example, the numerator in Eq. (13) becomes simply equal to $`h\alpha `$, since, according to Eq. (9), $`h^2(\alpha \delta _3\beta _3\gamma )=h^2(\alpha \delta \beta \gamma )\mathrm{\Pi }_3`$\].
Let us now take into account the spin of the particles. Consider the three different processes
$$\rho ^0(k)+\pi ^{}(q_1)\rho ^0(k^{})+\pi ^{}(q_1^{}),$$
(29)
$$\rho ^0(k)+\pi ^{}(q_1)\eta (q_2)+\pi ^{}(q_3),$$
(30)
$$\eta (p)+\pi ^{}(q)\eta (q_2)+\pi ^{}(q_3).$$
(31)
Let $`Q=k+q_1=k^{}+q_1^{}=q_2+q_3=p+q`$ and $`s=Q^2`$. Straightforward calculations with the help of the Lagrangian (3) of arbitrary terms of the relevant diagram series results in the following Lorentz structures and angular dependences for the corresponding physical amplitudes:
$$T^{(\lambda ^{},\lambda )}(\rho ^0\pi ^{}\rho ^0\pi ^{})=ϵ_{\mu ^{}\nu ^{}\tau ^{}\sigma }e_{(\lambda ^{})}^\mu ^{}q_1^\nu ^{}k^\tau ^{}ϵ_{\mu \nu \tau }^\sigma e_{(\lambda )}^\mu q_1^\nu k^\tau \stackrel{~}{T}^{}(\rho \pi \rho \pi )=$$
(32)
$$(\delta _{\lambda ,+1}+\delta _{\lambda ,1})(\delta _{\lambda ^{},+1}+\delta _{\lambda ^{},1})(s|\stackrel{}{q}_1|^2/2)(\lambda \lambda ^{}+\mathrm{cos}\theta )\stackrel{~}{T}^{}(\rho \pi \rho \pi ),$$
$$T^{(\lambda )}(\rho ^0\pi ^{}\eta \pi ^{})=ϵ_{\mu \nu \tau \sigma }e_{(\lambda )}^\mu q_1^\nu q_2^\tau q_3^\sigma \stackrel{~}{T}^{}(\rho \pi \eta \pi )=$$
(33)
$$(\delta _{\lambda ,+1}+\delta _{\lambda ,1})i\sqrt{s/2}|\stackrel{}{q}_1||\stackrel{}{q}_3|\mathrm{sin}\theta \stackrel{~}{T}^{}(\rho \pi \eta \pi ),$$
$$T(\eta \pi ^{}\eta \pi ^{})=|\stackrel{}{q}|^2\mathrm{cos}\theta \stackrel{~}{T}^{}(\eta \pi \eta \pi ),$$
(34)
where $`\lambda `$ ($`\lambda ^{}`$) is the initial (final) $`\rho `$ meson helicity and $`\theta `$ is the angle between the momenta of the initial and final pions in the reaction c.m. system. Certainly the dimensions of all physical amplitudes $`T`$ in Eqs. $`(20)(22)`$ are the same: the amplitudes are dimensionless. At the same time, as is seen from Eqs. $`(20)(22)`$, the invariant amplitudes $`\stackrel{~}{T}^{}`$ have different dimensions in the $`VPVP`$, $`VPPP`$, and $`PPPP`$ channels. These invariant amplitudes are obtained directly from the corresponding auxiliary amplitudes $`\stackrel{~}{T}`$ \[see Eqs. $`(10)(16)`$\] by substituting the physical dimensional coupling constants $`h`$ and $`h^{}`$ from the Lagrangian (3) and the following expressions for the $`p`$-wave loop integrals:
$`\mathrm{\Pi }_i={\displaystyle \frac{1}{16\pi }}{\displaystyle \frac{2}{3}}F_i\times \{\begin{array}{c}\hfill 4s,i=1,4(VP\text{loops}),\\ \hfill 1,i=2,3(PP\text{loops}),\end{array}`$ (37)
where
$$F_i=C_{1i}+sC_{2i}+\frac{s^2}{\pi }\underset{m_{i+}^2}{\overset{\mathrm{}}{}}\frac{[P_i(s^{})]^3ds^{}}{\sqrt{s^{}}s^{\mathrm{\hspace{0.17em}2}}(s^{}si\epsilon )}=C_{1i}+sC_{2i}+$$
(38)
$$\frac{(sm_{i+}^2)^{3/2}(sm_i^2)^{3/2}}{8\pi s^2}\left[\mathrm{ln}\left(\frac{\sqrt{sm_i^2}\sqrt{sm_{i+}^2}}{\sqrt{sm_i^2}+\sqrt{sm_{i+}^2}}\right)+i\pi \right]+$$
$$\frac{1}{4\pi }\{\frac{1}{2m_{i+}m_i}\mathrm{ln}\left(\frac{m_{i+}m_i}{m_{i+}+m_i}\right)[\frac{m_{i+}^4m_i^4}{s^2}\frac{3m_{i+}^2m_i^2}{2s}(m_{i+}^2+m_i^2)+$$
$$\frac{3}{8}(m_{i+}^4+m_i^4+6m_{i+}^2m_i^2)+\frac{s(m_{i+}^2+m_i^2)}{16m_{i+}^2m_i^2}(m_{i+}^410m_{i+}^2m_i^2+m_i^4)]+$$
$$\frac{m_{i+}^2m_i^2}{2s}\frac{5}{8}(m_{i+}^2+m_i^2)+\frac{s(3m_{i+}^4+3m_i^4+38m_{i+}^2m_i^2)}{48m_{i+}^2m_i^2}\}.$$
Here $`P_i(s)=[(sm_{i+}^2)(sm_i^2)/(4s)]^{1/2}`$, $`m_{i+}`$ ($`m_i`$) is the sum (the difference) of the particle masses in channel $`i`$, and $`C_{1i}`$ and $`C_{2i}`$ are the subtraction constants. Note that the expression (24) is valid for $`sm_{i+}^2`$. In the regions $`m_i^2<s<m_{i+}^2`$ and $`sm_i^2`$, it changes according to analytic continuation .
IV. ANALYSIS OF THE POSSIBLE RESONANCE PHENOMENA
First of all let us note that a number of free parameters in the present model can be reduced essentially, leaving its potentialities almost unchanged. So we shall assume that $`C_{11}=C_{14}`$, $`C_{21}=C_{24}`$ for the $`VP`$ loops and $`C_{12}=C_{13}`$, $`C_{22}=C_{23}`$ for $`PP`$ loops. Moreover, near a feasible resonance, the smooth $`s`$ dependence of the combinations $`C_{1i}+sC_{2i}`$ is not of crucial importance. Thus, as the essential free parameters we can leave only the $`C_{11}`$ and $`C_{12}`$ ones, setting $`C_{21}=C_{22}=0`$. Just this will be done in most variants considered below. A simplest way to discover “by hand” a possible resonance situation is that to find zero of Re$`(\overline{D})`$ at fixed values of $`h`$, $`h^{}`$, and $`\sqrt{s}`$, for example, at $`\sqrt{s}=1.43`$ GeV \[see Eqs. $`(11)(16)`$\]. In so doing the left free subtraction constants $`C_{11}`$ and $`C_{12}`$ are not uniquely determined. For example, in cases $`(i)`$ and $`(ii)`$, the condition Re$`(\overline{D})=0`$ gives only a relation of the type $`C_{12}=(\xi _1+\xi _2C_{11})/(\xi _3+\xi _4C_{11})`$, where $`\xi _i`$ are the known numbers. However, this is not the weak point of the model; on the contrary, this allows the shapes of the resonance curves and the relations between the absolute cross section values in the different channels to be easily changed by changing $`C_{11}`$.
According the detailed analysis performed in Refs. \[36-40,43\], the acceptable tentative values of the parameter $`\stackrel{~}{h}F_\pi ^3h`$ (where $`F_\pi 130`$ MeV) lie within the range $`|\stackrel{~}{h}|0.4`$. To illustrate the existence of the resonance phenomena in our toy model we are guided by the values of $`\stackrel{~}{h}`$ (and $`\stackrel{~}{h}^{}F_\pi ^3h^{}`$) near 0.1. We would like particularly to emphasize that, in fact, the resonance phenomena are possible in the present model for any $`|\stackrel{~}{h}|0.4`$. However, as $`|\stackrel{~}{h}|`$ (and/or $`|\stackrel{~}{h}^{}|`$) increases from 0.1 to 0.4, the distinct resonancelike enhancements in the reaction cross sections move into the region $`\sqrt{s}11.3`$ GeV. Note that the unitarized amplitudes essentially depend on the second and fourth powers of coupling constants and therefore are very sensitive to changes of $`|\stackrel{~}{h}|`$ and $`|\stackrel{~}{h}^{}|`$.
In Figs. 3 and 4, we show the typical energy dependences, which occur in our model for cases $`(i)`$, $`(ii)`$, and $`(iii)`$, for the four reaction cross sections $`\sigma (\rho ^0\pi ^{}\rho ^0\pi ^{})`$, $`\sigma (\rho ^0\pi ^{}\eta \pi ^{})`$, $`\sigma (\rho ^0\pi ^{}\eta ^{}\pi ^{})`$, and $`\sigma (\rho ^0\pi ^{}K^0K^{})`$ and for the phases of the $`\rho \pi \rho \pi `$ and $`\rho \pi \eta \pi `$ amplitudes (note that the inelastic amplitude $`\rho \pi \eta \pi `$ is defined only up to the sign). Figure 3, together with Table I, and Fig. 4, together with Table II, illustrate the resonance effects when they concentrate mainly in the regions $`\sqrt{s}1.31.4`$ GeV and $`\sqrt{s}1.51.6`$ GeV, respectively. As a rule, the channel $`\rho \pi \eta \pi `$ is dominant in case $`(i)`$, when all considered amplitudes belong to the $`\{10\}\{\overline{10}\}`$ representation of $`SU(3)`$. In case $`(ii)`$, when the amplitudes belong to the octet representation of $`SU(3)`$, the main channels are the $`\rho \pi \rho \pi `$ and $`\rho \pi \eta ^{}\pi `$ ones. In case $`(iii)`$, when $`h^{}=h`$ and the Lagrangian (3) possesses additional nonet symmetry, the cross sections for all channels except the $`K^{}K`$ one turn out to be comparable, and the general situation is rather complicated. The branching ratios of the presented resonancelike enhancements to the $`\rho ^0\pi ^{}(\rho ^0\pi ^{}`$, $`\eta \pi ^{}`$, $`\eta ^{}\pi ^{}`$, $`K^0K^{})`$ channels are listed in Tables I and II. Such characteristics for the complicated broad resonance structure can be defined as follows. For example, $`B(\rho ^0\pi ^{})=\overline{\sigma }(\rho ^0\pi ^{}\rho ^0\pi ^{})/\mathrm{\Sigma }`$, $`B(\eta \pi ^{})=\overline{\sigma }(\rho ^0\pi ^{}\eta \pi ^{})/\mathrm{\Sigma }`$, etc., where $`\mathrm{\Sigma }=2\overline{\sigma }(\rho ^0\pi ^{}\rho ^0\pi ^{})`$ \+ $`\overline{\sigma }(\rho ^0\pi ^{}\eta \pi ^{})`$ \+ $`\overline{\sigma }(\rho ^0\pi ^{}\eta ^{}\pi ^{})`$ \+ $`2\overline{\sigma }(\rho ^0\pi ^{}K^0K^{})`$ and every $`\overline{\sigma }`$ is the integral of the corresponding cross section over the $`\sqrt{s}`$ range from 1.2 to 1.8 GeV, where an enhancement concentrates.
Let us now compare the cross section values shown in Figs. 3 and 4 with those of $`a_2(1320)`$ resonance production. Using the tabular branching ratios we get $`\sigma (\rho ^0\pi ^{}a_2\rho ^0\pi ^{})5.7`$ mb and $`\sigma (\rho ^0\pi ^{}a_2\eta \pi ^{})2.36`$ mb at $`\sqrt{s}=m_{a_2}=1.32`$ GeV. Taking also into account the ratio of the factors $`(2J+1)/|\stackrel{}{k}|^2`$ for the $`a_2(1320)`$ resonance and for the $`J=1`$ enhancement found at $`\sqrt{s}1.31.4`$ GeV, or at $`\sqrt{s}1.51.6`$ GeV, we can conclude that we are certainly dealing with the resonancelike behavior of the $`I^G(J^{PC})=1^{}(1^+)`$ exotic waves in the region $`1.3\sqrt{s}1.6`$ GeV, at least, in the $`\rho \pi `$, $`\eta \pi `$, and $`\eta ^{}\pi `$ channels.
The most appreciable manifestation of an $`I^G(J^{PC})=1^{}(1^+)`$ exotic state in the mass region $`1.31.4`$ GeV has been observed in the $`\eta \pi ^0`$ channel in the charge exchange reaction $`\pi ^{}p\eta \pi ^0n`$ at 32, 38, and 100 GeV/$`c`$ (currently the exotic states of such a type are denoted most commonly as $`\pi _1`$). It was found that the intensity of the $`\pi _1`$ signal at its maximum in this reaction is only 3.5 times smaller than the corresponding intensity of the $`a_2(1320)`$ signal. It is very essential that $`a_2(1320)`$ production and $`\pi _1`$ production both proceed in this case via a single mechanism, namely, via the Reggeized $`\rho `$ exchange. If the $`\pi _1`$ really represents a complicated structure of the $`qq\overline{q}\overline{q}`$ or $`q\overline{q}g`$ type, then the fact that the $`\pi _1`$ production cross section in the charge exchange reaction has been found fully comparable with that of the conventional $`q\overline{q}`$ resonance $`a_2(1320)`$ is certainly strong evidence for the resonance nature of the observed exotic signal.
On the other hand, the intensity of $`\pi _1`$ production in the $`\eta \pi ^{}`$ channel in the reactions $`\pi ^{}p\eta \pi ^{}p`$ at 37 GeV/$`c`$ and 18 GeV/$`c`$ has been found to be about 15 and, respectively, 30 times smaller than the $`a_2(1320)`$ production intensity. This is evidently due to the more complicated mechanism of the reaction $`\pi ^{}p\eta \pi ^{}p`$ than that of the charge exchange reaction. In fact, there are three competing Regge exchanges with natural parity in this reaction: the $`\rho `$ exchange, the $`f_2`$ exchange, and the Pomeron one. Also, as is known, the last two are dominant in the case of $`a_2(1320)`$ production . Note that $`\pi _1`$ production can proceed via the Pomeron mechanism only owing to the octet component of the $`\pi _1`$. However, if this component is small, that is, if the $`\pi _1`$ belongs mainly to the $`\{10\}\{\overline{10}\}`$ representation of $`SU(3)`$, then $`\pi _1`$ production via Pomeron exchange has to be suppressed.
Another opportunity to observed $`\pi _1`$ and $`a_2(1320)`$ formation with comparable cross sections appears by using photoproduction (and electroproduction) processes, for example, $`\gamma p\rho ^0\pi ^{}\mathrm{\Delta }^{++}\pi ^+\pi ^{}\pi ^{}\mathrm{\Delta }^{++}`$, $`\gamma p\rho ^0\pi ^+n\pi ^+\pi ^{}\pi ^+n`$, $`\gamma p\eta \pi ^+n`$, and so on, which go at low momentum transfer mainly via the Reggeized one-pion exchange mechanism. Indeed, the existing data on the reactions $`\gamma p\rho ^0\pi ^{}\mathrm{\Delta }^{++}\pi ^+\pi ^{}\pi ^{}\mathrm{\Delta }^{++}`$ and $`\gamma p\rho ^0\pi ^+n\pi ^+\pi ^{}\pi ^+n`$ show a clear signature of the $`a_2(1320)`$ resonance and the appreciable enhancements in the $`3\pi `$ mass spectra in the range $`1.52`$ GeV . However, they do not yet allow certain conclusions to be made concerning the presence of the exotic wave in the $`\rho \pi `$ system and further investigations of the above reactions are needed.
At the present time an extensive program of the search for the exotic $`\pi _1`$ states in photoproduction experiments with high statistics and precision is planned for the Jefferson Laboratory \[6,14,15,18-20\]. A careful study of the $`\pi _1\gamma \pi `$ radiative decays in hadroproduction from nuclei via the Primakoff one-photon exchange mechanism is also planned with the CERN COMPASS spectrometer . A collection of the data on $`\pi _1`$ photoproduction, electroproduction, and hadroproduction and on the decays of the $`\pi _1`$ into $`\rho \pi `$ will also allow for the first time to verify the vector meson dominance model for states with exotic quantum numbers.
Summarizing we conclude that our calculation gives a further new reason in favor of the plausibility of the existence of an explicitly exotic resonance with $`I^G(J^{PC})=1^{}(1^+)`$ in the mass range $`1.31.6`$ GeV. Currently two exotic states at 1.4 and 1.6 GeV are extensively discussed in the literature \[7-20\]. In our scheme one does not succeed in simultaneously generating both the 1.4 and the 1.6 resonances, although the variants with a “fine structure” exist \[see, for example, the cross sections for case $`(iii)`$ in Figs. 3 and 4\]. Such a “fine structure” will be smoothed by the experimental resolution and we cannot certainly say about two resonances. The question may be raised as to whether this is a crucial result. It is not improbable that the inclusion of the $`b_1\pi `$ and $`f_1\pi `$ channels, where the exotic signals have also been found, can change the situation. However, the issue of the additional $`b_1\pi `$ and $`f_1\pi `$ channels remains open in the effective chiral Lagrangian approach. Notice also that at present the situation with the two exotic resonances at 1.4 and 1.6 GeV is not yet finally arranged.
ACKNOWLEDGMENTS
The present work was supported in part by grant INTAS-RFBR IR-97-232.
TABLE I. The parameter values of the model for the curves in Fig. 3. $`C_{11}`$ and $`C_{12}`$ are in GeV<sup>2</sup>; the other parameters are dimensionless; $`C_{21}=C_{22}=0`$ in all cases. Also, in the four right columns, the branching ratios of the resonancelike enhancement obtained to the partial channels are presented.
| Cases | $`F_\pi ^3h`$ | $`F_\pi ^3h^{}`$ | $`C_{11}`$ | $`C_{12}`$ | $`B(\rho ^0\pi ^{})`$ | $`B(\eta \pi ^{})`$ | $`B(\eta ^{}\pi ^{})`$ | $`B(K^0K^{})`$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- |
| $`(i)`$ | 0.10746 | 0 | 0.17 | 1.25 | 0.2377 | 0.3614 | 0.0125 | 0.0754 |
| $`(ii)`$ | 0 | 0.10746 | 0.34 | 0.67 | 0.3259 | 0.0890 | 0.1924 | 0.0289 |
| $`(iii)`$ | 0.10746 | 0.10746 | 0.49 | 0.50 | 0.2534 | 0.3619 | 0.1276 | 0.0019 |
TABLE II. The parameter values of the model for the curves in Fig. 4. $`C_{11}`$ and $`C_{12}`$ are in GeV<sup>2</sup>; the other parameters are dimensionless; $`C_{21}=C_{22}=0`$ in cases $`(i)`$ and $`(ii)`$ and $`C_{21}=C_{22}=0.11`$ in case $`(iii)`$. Also, in the four right columns, the branching ratios of the resonancelike enhancement obtained to the partial channels are presented.
| Cases | $`F_\pi ^3h`$ | $`F_\pi ^3h^{}`$ | $`C_{11}`$ | $`C_{12}`$ | $`B(\rho ^0\pi ^{})`$ | $`B(\eta \pi ^{})`$ | $`B(\eta ^{}\pi ^{})`$ | $`B(K^0K^{})`$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- |
| $`(i)`$ | 0.10746 | 0 | 0.18 | 0.76 | 0.1616 | 0.4634 | 0.0217 | 0.0959 |
| $`(ii)`$ | 0 | 0.08417 | 0.33 | 0.78 | 0.3032 | 0.0804 | 0.2184 | 0.0474 |
| $`(iii)`$ | 0.10746 | 0.10746 | 0.11 | 0.11 | 0.2429 | 0.3686 | 0.1356 | 0.0050 |
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# 1 Introduction
## 1 Introduction
Recently there has been renewed interest in the study of supergravity theories and in particular gauged supergravities due to the AdS/CFT correspondence as well as the domain-wall/QFT correspondence . Extended supergravity theories with $`N`$ extended local supersymmeties have a global $`SO(N)`$ invariance. Gauged supergravity theories arise when a subgroup of the $`R`$-symmetry group or the automorphism group of the supersymmetry algebra is gauged by the vector fields in the graviton supermultiplet. There are also other ways to obtain gauged models, like gauging the isometries of the vector multiplet moduli space as well as the hypermultiplet moduli space. The procedure of gauging does not change the particle content of the theory, but it introduces new terms proportional to the square of the gauge coupling constant in the action. So gauging necessarily induces either a cosmological constant (for $`N3`$) in which case, the supersymmetric ground state is an AdS space or in the presence of scalar fields (for $`N4`$), a scalar potential which is unbounded from below and it may or may not have critical points. Though the potential is inverted and unbounded from below, in the presence of gravity, this becomes a perfectly consistent theory having stable ground state configurations. If a background has killing spinors, then one can show that it is a stable background by applying Witten’s positivity of energy argument . Therefore, it is important to understand the nature of the ground states of these theories as well as the relationship between gauged supergravities and consistent compactifications/truncations of higher dimensional supergravity theories.
In four dimensions, there are two versions of $`N=4`$ supergravity theories, one with a global $`SO(4)`$ symmetry and the other one with a global $`SU(4)`$ symmetry . The equations of motion as derived from the two versions are equivalent by using field redefinition and duality transformations. However, when one considers the gauged models corresponding to the respective local internal symmetries, one finds that the two versions are inequivalent. The $`N=4`$, $`SO(4)`$ gauged model has one coupling constant and the scalar potential which is generated by gauging, is unbounded from below. On the otherhand, one can consider the Freedman-Schwarz model (FS) , where one considers gauging a $`SU(2)\times SU(2)`$ subgroup of $`SU(4)`$ internal symmetry with two independent gauge coupling constants. Here the scalar potential is again unbounded from below and has no critical points. But in both the cases, there exist stable vacuum configurations preserving some amount of supersymmetry. The electro-vac solutions in gauged $`SU(2)\times SU(2)`$ supergravity is one such well known example. Other backgrounds in the FS model, like domain walls, strings, pure axionic gravity etc have also been recently obtained and they preserve either half or one fourth of the supersymmetry . Nonabelian solitons and black holes as stable vacuum configurations were also shown to exist. In related work on strings in curved backgrounds , exact supersymmetric solutions of $`D=4`$ gauged supergravities have been constructed by using the techniques of conformal field theory and the connection between gauged supergravities and non-critical strings have been discussed.
In this work, we shall concentrate on the Euclidean Freedman Schwarz (EFS) model in $`D=4`$ which has recently been obtained by Volkov . The two theories (FS and EFS) are different as they are obtained from compactification of the ten dimensional theory on different group manifolds and it is to be noted that they are not just related by analytic continuations. $`D=4`$, $`N=4`$ gauged $`SU(2)\times SU(2)`$ FS model can be embedded into $`N=1`$ supergravity in ten dimensions as an $`S^3\times S^3`$ compactification with the group manifold being $`SU(2)\times SU(2)`$ . Previously also, a Kaluza-Klein (KK) interpretation for the $`SU(2)\times SU(2)`$ gauged supergravity was given in , where the model was identified as part of the effective $`D=4`$ field theory for the heterotic string theory on $`S^3\times S^3`$. These two KK interpretations are essentially the same upto consistent truncations. One can also consider another reduction of the $`N=1`$ ten dimensional theory on the group manifold $`SU(2)\times SU(1,1)`$ so that the geometry of the internal space-time is $`S^3\times AdS_3`$ with the signature $`(+,+,+,+,+,)`$ and the corresponding four dimensional theory becomes an Euclidean theory. As the scalar curvature of $`S^3`$ is positive and that of $`AdS_3`$ is negative, the dilaton or equivalently the scalar potential in the corresponding four dimensional theory becomes proportional to $`g_1^2g_2^2`$, where, $`g_1`$ and $`g_2`$ are the gauge coupling constants corresponding to $`SU(2)`$ and $`SU(1,1)`$ respectively. Since the potential is proportional to the square of the difference of the gauge couplings, one can consider a variety of cases, where the potential can be positive, negative or zero. The dimensional reduction on the above group manifold is consistent in the sense that for a given four dimensional configuration which satisfies the four dimensional equations of motion and supersymmety variations, the corresponding uplifted version also satisfies the ten dimensional equations of motion as well as the supersymmetry variations.
The paper is organized as follows: In section 2, we discuss the four dimensional gauged EFS model as obtained from the dimensional reduction of the corresponding ten dimensional theory. In section 3 and 4, we explicitly obtain the new background solutions and illustrate that they preserve either half or one fourth of the original supersymmety. The Euclidean solutions we have obtained, include the interesting cases of domain wall, $`E^2\times S^2`$, $`E^2\times AdS_2`$, $`E^1\times S^3`$ where $`E^1`$ and $`E^2`$ denote one and two dimensional Euclidean spaces. We also show that the four dimensional gravitational instanton solutions like Eguchi-Hanson is a solution of the EFS model with nontrivial (anti)self-dual abelian gauge fields belonging to the $`U(1)`$ of $`SU(2)`$ and the noncompact $`SU(1,1)`$ groups. In section 5, we summarize our results.
## 2 The Euclidean Freedman-Schwarz model
In this section we set our notations and briefly review some necessary aspects of the EFS model which will be necessary for our analysis. The field content of the EFS model is same as that of the FS model. In the EFS model, the four dimensional gravity multiplet contains the graviton $`E_\mu ^m`$, four majorana spin $`\frac{3}{2}`$ gravitinos $`\mathrm{\Psi }_\mu ^I(I=1,\mathrm{}4)`$, three nonabelian vector fields $`A_\mu ^a(a=1,2,3)`$ belonging to $`SU(2)`$ with gauge coupling $`g_1`$, three nonabelian pseudovector gauge fields $`\dot{A}_\mu ^a`$ belonging to $`SU(1,1)`$ group with gauge coupling constant $`g_2`$, four majorana spin $`\frac{1}{2}`$ fields $`\chi ^I`$, the axion $`𝐚`$ and the dilaton $`\mathrm{\Phi }`$. Here the Greek indices $`\mu ,\nu ,\mathrm{}`$ refer to the base space indices and latin indices $`m,n,\mathrm{}`$ refer to the tangent space indices. The bosonic part of the ten dimensional theory contains the metric $`\widehat{g}_{MN}(M,N,\mathrm{}=1,\mathrm{}10)`$, the three form antisymmetric tensor $`\widehat{H}_{MNP}`$ and the dilaton $`\widehat{\mathrm{\Phi }}`$. The fermionic field contents are the ten dimensional gravitino $`\widehat{\mathrm{\Psi }}_M`$ and the gaugino $`\widehat{\chi }`$. We consider vanishing spinor background fields, however their supersymmetric variations do not vanish and they are important for our considerations.
The bosonic part of the ten dimensional action corresponding to $`N=1`$ supergravity is given by,
$`S_{10}`$ $`=`$ $`{\displaystyle \sqrt{\widehat{g}}d^{10}\widehat{x}\left(\frac{1}{4}\widehat{R}\frac{1}{2}_M\widehat{\mathrm{\Phi }}^M\widehat{\mathrm{\Phi }}\frac{1}{12}e^{2\widehat{\mathrm{\Phi }}}\widehat{H}_{MNP}\widehat{H}^{MNP}\right)}`$ (1)
where $`\widehat{R}`$ is the curvature scalar in $`D=10`$. The equations of motion following from this action are given by,
$`\widehat{}_M\widehat{}^M\widehat{\mathrm{\Phi }}+{\displaystyle \frac{1}{6}}e^{2\widehat{\mathrm{\Phi }}}\widehat{H}_{MNP}\widehat{H}^{MNP}`$ $`=`$ $`0`$ (2)
$`\widehat{}_M(e^{2\widehat{\mathrm{\Phi }}}\widehat{H}^{MNP})`$ $`=`$ $`0`$ (3)
$`\widehat{R}_{MN}2_M\widehat{\mathrm{\Phi }}_N\widehat{\mathrm{\Phi }}e^{2\widehat{\mathrm{\Phi }}}\widehat{H}_{MPQ}\widehat{H}_N^{PQ}+{\displaystyle \frac{1}{12}}e^{2\widehat{\mathrm{\Phi }}}\widehat{g}_{MN}\widehat{H}_{PQS}\widehat{H}^{PQS}`$ $`=`$ $`0`$ (4)
The dimensional reduction of the above ten dimensional theory in terms of suitable parametization has been discussed in a recent paper by Volkov . The corresponding four dimensional equations of motion for metric, dilaton, axion and gauge fields are respectively given by,
$`R_{\mu \nu }2_\mu \mathrm{\Phi }_\nu \mathrm{\Phi }+2e^{4\mathrm{\Phi }}_\mu 𝐚_\nu 𝐚`$
$`2e^{2\mathrm{\Phi }}\left[\eta _{ab}^{(1)}\left(F_{\mu \lambda }^aF_\nu ^{b\lambda }{\displaystyle \frac{1}{4}}g_{\mu \nu }F_{\lambda \rho }^aF^{b\lambda \rho }\right)+\eta _{ab}^{(2)}\left(\dot{F}_{}^{a}{}_{\mu \lambda }{}^{}\dot{F}_\nu ^{b\lambda }{\displaystyle \frac{1}{4}}g_{\mu \nu }\dot{F}_{}^{a}{}_{\lambda \rho }{}^{}\dot{F}^{b\lambda \rho }\right)\right]`$
$`2g_{\mu \nu }U(\mathrm{\Phi })=0`$ (5)
$`_\mu ^\mu \mathrm{\Phi }2e^{4\mathrm{\Phi }}_\mu 𝐚^\mu 𝐚{\displaystyle \frac{1}{2}}e^{2\mathrm{\Phi }}\left[\eta _{ab}^{(1)}F_{\mu \nu }^aF^{b\mu \nu }+\eta _{ab}^{(2)}\dot{F}_{\mu \nu }^a\dot{F}^{b\mu \nu }\right]+2U(\mathrm{\Phi })=0`$ (6)
$`_\mu (e^{4\mathrm{\Phi }}^\mu 𝐚)+{\displaystyle \frac{1}{2}}\left[\eta _{ab}^{(1)}F_{\mu \nu }^aF^{b\mu \nu }+\eta _{ab}^{(2)}\dot{F}_{\mu \nu }^a\dot{F}^{b\mu \nu }\right]=0`$ (7)
$`_\rho (e^{2\mathrm{\Phi }}F^{a\rho \mu })+g_1e^{2\mathrm{\Phi }}f_{}^{a}{}_{bc}{}^{}A_{}^{b}{}_{\rho }{}^{}F^{c\rho \mu }2F^{a\mu \rho }_\rho 𝐚=0`$ (8)
$`_\rho (e^{2\mathrm{\Phi }}\dot{F}^{a\rho \mu })+g_2e^{2\mathrm{\Phi }}\dot{f}_{}^{a}{}_{bc}{}^{}\dot{A}_{}^{b}{}_{\rho }{}^{}\dot{F}^{c\rho \mu }2\dot{F}^{a\mu \rho }_\rho 𝐚=0`$ (9)
where $`U(\mathrm{\Phi })`$ is the dilaton potential given by,
$`U(\mathrm{\Phi })`$ $`=`$ $`{\displaystyle \frac{1}{8}}(g_1^2g_2^2)e^{2\mathrm{\Phi }}`$ (10)
The structure constants are given by,
$$f_{}^{c}{}_{ab}{}^{}=\eta ^{(1)cd}ϵ_{dab};\dot{f}_{}^{c}{}_{ab}{}^{}=\eta ^{(2)cd}ϵ_{dab}$$
(11)
$`ϵ_{abc}`$ is the antisymmetric tensor, $`\eta _{ab}^{(1)}`$ and $`\eta _{ab}^{(2)}`$ are the cartan metrics corresponding to $`SU(2)`$ and $`SU(1,1)`$ respectively, where $`\eta _{ab}^{(1)}=diag(1,1,1)`$ and $`\eta _{ab}^{(2)}=diag(1,1,1)`$. The dual field strengths in the four dimensional Euclidean theory are defined as,
$`F_{\mu \nu }^a`$ $`=`$ $`{\displaystyle \frac{1}{2}}\sqrt{g}ϵ_{\mu \nu \lambda \rho }F^{a\lambda \rho }`$ (12)
and,
$`\dot{F}_{\mu \nu }^a`$ $`=`$ $`{\displaystyle \frac{1}{2}}\sqrt{g}ϵ_{\mu \nu \lambda \rho }\dot{F}^{a\lambda \rho }`$ (13)
The above equations of motion can be obtained from the four dimensional Euclidean Freedman-Schwarz action,
$`S_4={\displaystyle }\sqrt{g}d^4x[{\displaystyle \frac{R}{4}}{\displaystyle \frac{1}{2}}_\mu \mathrm{\Phi }^\mu \mathrm{\Phi }+{\displaystyle \frac{1}{2}}e^{4\mathrm{\Phi }}_\mu 𝐚^\mu 𝐚{\displaystyle \frac{1}{4}}e^{2\mathrm{\Phi }}(\eta _{ab}^{(1)}F_{\mu \nu }^aF^{b\mu \nu }+\eta _{ab}^{(2)}\dot{F}_{\mu \nu }^a\dot{F}^{b\mu \nu })`$
$`{\displaystyle \frac{1}{2}}𝐚(\eta _{ab}^{(1)}F_{\mu \nu }^aF^{b\mu \nu }+\eta _{ab}^{(2)}\dot{F}_{\mu \nu }^a\dot{F}^{b\mu \nu })+{\displaystyle \frac{1}{8}}(g_1^2g_2^2)e^{2\mathrm{\Phi }}]`$ (14)
Similarly the ten dimensional spinors can also be consistently reduced and the four dimensional supersymmetry variations are given by,
$`\delta \chi =\left({\displaystyle \frac{1}{\sqrt{2}}}\gamma ^\mu _\mu \mathrm{\Phi }{\displaystyle \frac{1}{\sqrt{2}}}e^{2\mathrm{\Phi }}\gamma _5\gamma ^\mu _\mu 𝐚\right)ϵ+{\displaystyle \frac{1}{2}}e^\mathrm{\Phi }\left({\displaystyle \frac{1}{2}}\eta _{ab}^{(1)}\gamma ^\alpha \gamma ^\beta F_{\alpha \beta }^a\alpha ^b{\displaystyle \frac{1}{2}}\gamma _5\eta _{ab}^{(2)}\gamma ^\alpha \gamma ^\beta \dot{F}_{\alpha \beta }^a\dot{\alpha }^b\right)ϵ`$
$`+{\displaystyle \frac{1}{4}}e^\mathrm{\Phi }(g_1g_2\gamma _5)ϵ,`$
$`\delta \mathrm{\Psi }_\mu =\left(_\mu +{\displaystyle \frac{1}{4}}\omega _\mu ^{\alpha \beta }\gamma _\alpha \gamma _\beta {\displaystyle \frac{g_1}{2}}\eta _{ab}^{(1)}\alpha ^aA_\mu ^b+{\displaystyle \frac{g_2}{2}}\eta _{ab}^{(2)}\dot{\alpha }^a\dot{A}_\mu ^b+{\displaystyle \frac{1}{2}}e^{2\mathrm{\Phi }}\gamma _5_\mu 𝐚\right)ϵ+`$
$`{\displaystyle \frac{1}{2\sqrt{2}}}e^\mathrm{\Phi }\left(\eta _{ab}^{(1)}F_{\lambda \nu }^a\alpha ^b+\gamma _5\eta _{ab}^{(2)}\dot{F}_{\lambda \nu }^a\dot{\alpha }^b\right)\gamma ^\lambda \gamma ^\nu \gamma _\mu ϵ+{\displaystyle \frac{1}{4\sqrt{2}}}e^\mathrm{\Phi }(g_1+g_2\gamma _5)\gamma _\mu ϵ`$ (15)
where $`ϵ`$ is the Majorana spinor corresponding to the supersymmetry transformation parameter. Here, $`\gamma _5=\gamma ^0\gamma ^1\gamma ^2\gamma ^3`$ and $`\{\gamma _5,\gamma ^\alpha \}=0`$. $`\gamma ^\alpha `$ are the four dimensional tangent space gamma matrices satisfying the usual anticommutation relation $`\{\gamma ^\alpha ,\gamma ^\beta \}=2\eta ^{\alpha \beta }`$ with $`\eta ^{\alpha \beta }=diag(+1,+1,+1,+1)`$. $`\omega _\mu ^{\alpha \beta }`$ are the spin connections and $`\alpha ^a`$ and $`\dot{\alpha }^a`$ are the $`4\times 4`$ matrices which generate the Lie algebra of the group $`SU(2)`$ and $`SU(1,1)`$ respectively with the properties,
$`\alpha ^a\alpha ^b`$ $`=`$ $`ϵ^{abc}\eta _{cd}^{(1)}\alpha ^d\eta ^{(1)ab}`$ (16)
$`\dot{\alpha }^a\dot{\alpha }^b`$ $`=`$ $`ϵ^{abc}\eta _{cd}^{(2)}\dot{\alpha }^d+\eta ^{(2)ab}`$ (17)
$`[\alpha ^a,\dot{\alpha }^a]`$ $`=`$ $`0`$ (18)
The corresponding matrix representation is given by,
$$\alpha ^a=i\tau ^aI_2(a=1,2,3);\dot{\alpha }^b=I_2\tau ^b(b=1,2);\dot{\alpha }^3=iI_2\tau ^3$$
(19)
where $`\tau `$’s are Pauli matrices.
In next section, we shall show that one can construct many interesting stable new vacua for the EFS model which are supersymmetric and are consistent with the four dimensional background equations of motion.
## 3 Supersymmetric configurations in EFS model
In the examples below, we choose specific $`U(1)`$ directions thereby spontaneously breaking $`SU(2)`$ to $`U(1)`$ and in the noncompact case, $`SU(1,1)`$ to $`U(1)`$, similar to Freedman-Gibbons electro-vac solutions with constant dilaton. This corresponds to setting the other two gauge fields of $`SU(2)`$ triplets or the $`SU(1,1)`$ triplet to zero vacuum expectation value. So whenever we have nonzero gauge fields, they are basically abelian.
### 3.1 Euclidean domain walls
First we consider the four dimensional Euclidean domain wall obtained by analytically continuing the Lorenztian domain walls with the field configurations,
$`ds^2=U(y)(dt^2+dx_1^2+dx_2^2)+U^1(y)dy^2,`$
$`\mathrm{\Phi }={\displaystyle \frac{1}{2}}\mathrm{ln}U(y),U(y)=m|yy_0|,`$
$`A_\mu ^a=0,\dot{A}_\mu ^a=0,𝐚=0,`$ (20)
This background is singular at $`y=y_0`$. Since there is no other matter field present here other than dilaton, the above configuration represents pure dilaton gravity in Euclidean space. We now study the supersymmetric properties of this background. For $`g_10`$ (in fact the gauge coupling constant and the mass parameter are related by $`g_1^2=2m^2`$) and $`g_2=0`$, if we substitute the above background in the supersymmetry equations, we find that the fermionic variations vanish provided the supersymmetry parameters satisfy,
$$ϵ=\gamma _3ϵ,ϵ=U(y)^{\frac{1}{4}}ϵ_0$$
(21)
where $`ϵ_0`$ is a constant spinor. These conditions break half of the supersymmetry. Thus we see that there exist nontrivial killing spinors preserving $`N=2`$ supersymmetry for pure dilatonic Euclidean domain wall background.
### 3.2 $`E^2\times S^2`$
Here we consider the case where dilaton is constant. We also take one of the $`U(1)`$ gauge fields of the $`SU(2)`$ part to be nonvanishing and the geometry as that of $`E^2\times S^2`$. Corresponding field configurations are given by,
$`ds^2=d\psi ^2+d\chi ^2+{\displaystyle \frac{1}{B}}\left(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2\right)`$
$`F^a=\delta ^{a3}Q\mathrm{sin}\theta d\theta d\varphi ;\mathrm{\Phi }=\mathrm{\Phi }_0=constant`$
$`\dot{F}=0,𝐚=constant,g_1=0`$ (22)
where $`\frac{1}{\sqrt{B}}`$ is the constant radius of the two sphere. This configuration satisfies the equations of motion with $`Q=\frac{1}{g_2}`$. The above solution is analogous to the magnetic solution in the charged Nariai black hole background , but here we are in the Euclidean space with $`F^2=2B^2Q^2`$ which is strictly positive.
Next, we discuss the supersymmetric property of this background. The only nonvanishing spin connection is $`\omega _\varphi ^{23}=\mathrm{cos}\theta `$. The components of the killing spinor equations are given by,
$`_\psi ϵ=0`$ (23)
$`_\chi ϵ=0`$ (24)
$`_\theta ϵ+{\displaystyle \frac{1}{2}}\gamma _5\gamma _2ϵ=0`$ (25)
$`_\varphi ϵ{\displaystyle \frac{1}{2}}\mathrm{cos}\theta \gamma _2\gamma _3ϵ+{\displaystyle \frac{1}{2}}\mathrm{sin}\theta \gamma _5\gamma _3ϵ=0`$ (26)
with the condition
$$\left(\gamma _5\gamma _2\gamma _3\alpha ^31\right)ϵ=0.$$
The complete set of killing spinors which are the solution to the above equations are
$$ϵ=e^{\frac{1}{2}\theta \gamma _5\gamma _2}e^{\frac{1}{2}\varphi \gamma _3\gamma _2}\left(\frac{\gamma _5\gamma _2\gamma _3\alpha ^3+1}{2}\right)ϵ_0$$
(27)
where $`ϵ_0`$ is some constant spinor and $`\alpha ^3`$ is in $`SU(2)`$. The operator $`\left[\gamma _5\gamma _2\gamma _3\alpha ^3+1\right]`$ acts as a projection operator, hence breaks $`\frac{1}{2}`$ supersymmetry.
### 3.3 $`E^2\times AdS_2`$
Here we consider the geometry $`E^2\times AdS_2`$ with constant dilaton and the nonvanishing $`U(1)`$ gauge field corresponding to the noncompact part of $`SU(1,1)`$. This choice of gauge field is necessary so that the background equations of motion are satisfied.
The field configurations are given by,
$`ds^2=d\psi ^2+d\chi ^2+{\displaystyle \frac{1}{B}}\left(r^2dt^2+{\displaystyle \frac{dr^2}{r^2}}\right)`$
$`\dot{F}^a=\delta ^{\dot{a}3}Qdtdr,\mathrm{\Phi }=\mathrm{\Phi }_0=constant`$
$`F=0,𝐚=constant,g_2=0`$ (28)
where $`\frac{1}{\sqrt{B}}`$ corresponds to the radius of the AdS space. The above background fields satisfy the equations of motion with $`Q=\frac{1}{g_1}`$. The only nonzero spin connection is given by $`\omega _t^{23}=r`$. Considering that the fermionic supersymmetry variations vanish, one gets the following equations for the $`\psi `$, $`\chi `$, $`t`$ and $`r`$ components for the killing spinor:
$`_\psi ϵ=0`$ (29)
$`_\chi ϵ=0`$ (30)
$`_tϵ+{\displaystyle \frac{r}{2}}\gamma _2\gamma _3ϵ+{\displaystyle \frac{1}{2}}r\gamma _2ϵ=0`$ (31)
$`_rϵ+{\displaystyle \frac{1}{2r}}\gamma _3ϵ=0`$ (32)
The projector is given by $`\left[\gamma _5\gamma _2\gamma _3\dot{\alpha }^3+1\right]`$ where $`\dot{\alpha }^\mathrm{𝟑}`$ is along the noncompact direction in $`SU(1,1)`$. The killing spinors are
$`ϵ={\displaystyle \frac{1}{2}}r^{\frac{1}{2}}\left(\gamma _5\gamma _2\gamma _3\dot{\alpha }^3+1\right)ϵ_{}+{\displaystyle \frac{1}{2}}\left[r^{\frac{1}{2}}r^{\frac{1}{2}}\gamma _2t\right]\left(\gamma _5\gamma _2\gamma _3\dot{\alpha }^3+1\right)ϵ_+`$ (33)
where $`ϵ_{}`$ are constant spinors and they satisfy the conditions $`\gamma _3ϵ_{}=ϵ_{}`$. So as before, this background preserves one half of the supersymmetry.
These last two EFS solutions are analogous to the electro-vac solution in FS model. The later one, $`E^2\times AdS_2`$, can be mapped to the Lorentzian section by applying the transformations as dicussed in .
### 3.4 $`E^1\times S^3`$
Here, we would like to obtain a background analogous to the pure axionic gravity solution in FS model . So, we take nontrivial axion field while the dilaton as well as gauge fields are vanishing. The field configuration is given by,
$`ds^2=d\psi ^2+{\displaystyle \frac{1}{Q}}\left(d\chi ^2+\mathrm{sin}^2\chi \left(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2\right)\right)`$
$`𝐚=\pm \sqrt{Q}\psi ;\mathrm{\Phi }=0`$
$`A=0;\dot{A}=0;g_1=0`$ (34)
The equations of motion of these background fields are consistent with $`Q=\frac{g_2^2}{8}`$. To study the supersymmetry properties, we need to calculate the spin connections on three sphere. The nonzero components are given by,
$`\omega _\theta ^{21}=\mathrm{cos}\chi ;\omega _\varphi ^{31}=\mathrm{cos}\chi \mathrm{sin}\theta ;\omega _\varphi ^{32}=\mathrm{cos}\theta `$ (35)
The projector is given by,
$$ϵ=\gamma _0ϵ,$$
(36)
and the components of the killing equations are given by,
$`_\psi ϵ=0`$ (37)
$`_\chi ϵ+{\displaystyle \frac{1}{2}}\gamma _5\gamma _1ϵ=0`$ (38)
$`_\theta ϵ+{\displaystyle \frac{1}{2}}\mathrm{cos}\chi \gamma _2\gamma _1ϵ+{\displaystyle \frac{1}{2}}\mathrm{sin}\chi \gamma _5\gamma _2ϵ=0`$ (39)
$`_\varphi ϵ+{\displaystyle \frac{1}{2}}\mathrm{cos}\chi \mathrm{sin}\theta \gamma _3\gamma _1ϵ+{\displaystyle \frac{1}{2}}\mathrm{sin}\chi \mathrm{sin}\theta \gamma _5\gamma _3ϵ+{\displaystyle \frac{1}{2}}\mathrm{cos}\theta \gamma _3\gamma _2ϵ=0`$ (40)
The complete solution of the killing spinor equation is given by,
$`ϵ=e^{\frac{1}{2}\chi \gamma _5\gamma _1}e^{\frac{1}{2}\theta \gamma _2\gamma _1}e^{\frac{1}{2}\varphi \gamma _3\gamma _2}\left[{\displaystyle \frac{\gamma _0+1}{2}}\right]ϵ_0`$ (41)
Hence this choice of field configurations breaks $`\frac{1}{2}`$ of the supersymmetry.
We find that $`E^1\times AdS_3`$ background can also be a solution but then one has to consider imaginary axion field to solve of the background equations of motion.
## 4 Gravitational Instantons
In this section we consider gravitational instantons which are solutions in Euclidean gravity. It has been noted in that with vanishing dilaton, axion, gauge fields and for $`g_1=g_2`$, the flat gravitational instantons (cosmological constant being zero) are vacua of EFS model. Here we show that even in the presence of (anti)self-dual gauge fields, the Eguchi-Hanson instanton satisfying the flat space Einstein equations is a consistent background of this EFS model and it preserves certain fraction of the supersymmetry.
This is one of the examples, where we keep both the gauge coupling constants and we take them to be equal, $`g_1=g_2`$. We take nonzero $`U(1)`$ gauge field belonging to the noncompact part of $`SU(1,1)`$ and the nonvanishing $`U(1)`$ gauge field of the $`SU(2)`$ part could be any of the triplet (let us choose $`A^3`$ to be nonzero). The field configuration is given by,
$`ds^2={\displaystyle \frac{dr^2}{1\frac{a^4}{r^4}}}+{\displaystyle \frac{r^2}{4}}\left(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2\right)+{\displaystyle \frac{r^2}{4}}\left(1{\displaystyle \frac{a^4}{r^4}}\right)\left(d\psi +\mathrm{cos}\theta d\varphi \right)^2`$
$`F={\displaystyle \frac{2}{r^4}}\left(e^3e^0+e^1e^2\right)=\dot{F}`$
$`\mathrm{\Phi }=0;𝐚=0;g_1=g_2`$ (42)
where the vierbeins are
$`e^0={\displaystyle \frac{dr}{\sqrt{1\frac{a^4}{r^4}}}};e^1={\displaystyle \frac{r}{2}}\left(\mathrm{sin}\psi d\theta \mathrm{sin}\theta \mathrm{cos}\psi d\varphi \right)`$
$`e^2={\displaystyle \frac{r}{2}}\left(\mathrm{cos}\psi d\theta \mathrm{sin}\theta \mathrm{sin}\psi d\varphi \right);e^3={\displaystyle \frac{r}{2}}\sqrt{1{\displaystyle \frac{a^4}{r^4}}}\left(d\psi +\mathrm{cos}\theta d\varphi \right).`$ (43)
One can immediately note that the gauge field strengths in (42) are anti-self-dual. The spin connections which are also anti-self-dual can be calculated from (43) and these are
$`\omega _\theta ^{10}=\omega _\theta ^{23}={\displaystyle \frac{1}{2}}\sqrt{1{\displaystyle \frac{a^4}{r^4}}}\mathrm{sin}\psi ;\omega _\varphi ^{10}=\omega _\varphi ^{23}={\displaystyle \frac{1}{2}}\sqrt{1{\displaystyle \frac{a^4}{r^4}}}\mathrm{sin}\theta \mathrm{cos}\psi `$
$`\omega _\theta ^{20}=\omega _\theta ^{31}={\displaystyle \frac{1}{2}}\sqrt{1{\displaystyle \frac{a^4}{r^4}}}\mathrm{cos}\psi ;\omega _\varphi ^{20}=\omega _\varphi ^{31}={\displaystyle \frac{1}{2}}\sqrt{1{\displaystyle \frac{a^4}{r^4}}}\mathrm{sin}\theta \mathrm{sin}\psi `$
$`\omega _\psi ^{30}=\omega _\psi ^{12}={\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{a^4}{r^4}}\right);\omega _\varphi ^{30}=\omega _\varphi ^{12}=\mathrm{cos}\theta \omega _\psi ^{30}.`$ (44)
With the above choice of background fields, the supersymmetry variations give the projector conditions,
$`\left(1\gamma _5\right)ϵ=0`$
$`\left(\alpha ^3+\dot{\alpha }^3\right)ϵ=0`$ (45)
With these projectors, the killing spinor equations are really simplified:
$`_\mu ϵ=0.`$ (46)
So the killing spinors are independent of $`r,\theta ,\varphi ,\psi `$. Because of the twin supersymmetric conditions, the solution preserves $`\frac{1}{4}`$ of the supersymmetry. However, once the gauge field backgrounds are switched off the second condition in eq.(45) will drop out and the pure gravitational instanton background will become half supersymmetric.
## 5 Summary
In this work, we have obtained stable vacuum configurations in $`D=4`$, $`N=4`$ $`SU(2)\times SU(1,1)`$ gauged supergravity theory (EFS model) which has been obtained from dimensional reduction of $`D=10`$, $`N=1`$ supergravity on $`S^3\times AdS_3`$. We have obtained stable domain wall solutions preserving half of the supersymmetry. We have then considered vacua like $`E^2\times S^2`$, $`E^2\times AdS_2`$ with constant dilaton and nontrivial $`U(1)`$ gauge fields which are analogous to the electrovac solutions in FS model preserving half of the supersymmetry. We also obtained geometry like $`E^1\times S^3`$ with nontrivial axion and vanishing dilaton field preserving half of the supersymmetry. Finally we obtained Euclidean gravitational instanton, namely Eguchi-Hanson instanton with nontrivial abelian gauge fields and vanishing dilaton and axion fields as a consistent vaccum configuration breaking one fourth of the supersymmetry. The dilaton potential in this case vanishes as the two gauge coupling constants are equal. So our findings of a rich variety of vacua for the EFS model makes the theory more interesting and worth exploring further.
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# 1 Introduction
## 1 Introduction
The Standard Model (SM) of electroweak interactions, while eminently successful in describing most available data, is rightly not regarded as the final theory. Apart from aesthetic drawbacks such as the lack of explanations for either the fermion masses or the relative strength of the gauge couplings, it also suffers from a few technical problems. To overcome these lacunae, many authors have, over the years, proposed models going beyond the SM. Two of the most attractive classes of such models are those incorporating grand unification and/or supersymmetry . Both, as well as other models, predict additional particle states. In this article we shall concentrate on the scalar sector of such theories.
Within the SM, baryon ($`B`$) and lepton ($`L`$) number conservation come about due to accidental symmetries. In other words, such conservation is not guaranteed by any theoretical reasons, but are rather the consequences of the choice of the particle content<sup>1</sup><sup>1</sup>1Indeed, nonperturbative effects within the SM itself do break $`B+L`$ symmetry.. In extensions of the SM, such an accidental occurrence is obviously not guaranteed. For example, even in the simplest grand unified theories (GUTs), both the gauge and the scalar sector interactions violate each of $`B`$ and $`L`$. The corresponding particles, namely the diquarks and leptoquarks have been studied in the literature to a considerable extent. Simultaneous breaking of both $`B`$ and $`L`$ symmetry can be disastrous though as that would lead to rapid proton decay. Within GUTs, gauge boson-mediated proton decay is naturally suppressed on account of their large mass; in the scalar sector, the particle content can be chosen such that there is no diquark-leptoquark mixing, at least as far as the light sector is concerned.
In the case of the Minimal Supersymmetric Standard Model (MSSM), on the other hand, we do not have the option of demanding the ‘offending’ fields (the supersymmetric partners of the SM fermions) to be superheavy. Ruling out the undesirable terms necessitates the introduction of a discrete symmetry, $`R(1)^{3(BL)+2S}`$ (with $`S`$ denoting the spin of the field) . Apart from ruling out both $`B`$ and $`L`$ violating terms in the superpotential, this symmetry has the additional consequence of rendering the lightest supersymmetric partner absolutely stable. However, such a symmetry is ad hoc. Hence, it is of interest to consider possible violations of this symmetry especially since it has rather important experimental consequences, not the least of which concerns the detection of the supersymmetric partners.
As we shall see later, there are certain similarities between $`R\text{/}`$ interactions on the one hand and leptoquarks and/or diquarks on the other. An example of this is afforded by the explanations of the anomalous large-$`Q^2`$ data reported by the Hera collaborations. However, the $`R\text{/}`$-MSSM, being a richer (low-energy) theory, offers a larger set of possibilities, both in the context of the Hera events as well as other anomalies like the ones observed at Karmen or at Kamiokande. Hence, in our discussions, we shall concentrate primarily on the supersymmetric case, and point out, as special cases, the corresponding results for theories with leptoquarks or diquarks.
It is only natural that such particles would be looked for in existing colliders. Indeed, the best lower bounds on the masses of such particles have been obtained from analyses Tevatron data. Pair production of leptoquarks (or equivalently squarks decaying through an $`L`$ violating interaction) leads to a final state comprising a dilepton pair alongwith jets . More interestingly, in the supersymmetric case, gluino production cross-section is larger and, in addition, can lead to like-sign dileptons, thereby making the signal stand out even more . However, all such analyses, of necessity, make certain assumptions that are not necessarily true. Similar is the case of searches at Hera which depend on the size of the $`L`$-violating coupling . To minimize the dependence on such assumptions, it is thus necessary to consider further experiments.
In this article we investigate the production of such scalars (supersymmetric or otherwise) at a photon collider. The lepton (or baryon) number violating decays may result in a significant excess in dijet plus dilepton or 4-jet final states. An analysis of such data would allow us not only to detect such particles but also to measure their masses and, to an extent, the branching fractions. The plan of the rest of the article is as follows. In the next section we discuss briefly the photon photon collider and the production cross-section of these scalar fermions at a such a collider. In section 3, we focus on the $`R\text{/}`$ violating couplings and the decays of the squarks and the sleptons. Section 4 will be devoted to comparing the signals and the possible SM backgrounds. We will explore the discovery/ exclusion limits on the SUSY parameter space using our signal and background analysis in section 5. Finally, we summarise in section 6.
## 2 Pair-production of Scalars at a Photon Collider
To the leading order in perturbation theory, the cross-section for charged particle pair-production at a gamma gamma collider is completely model independent This is quite unlike the case for the $`e^+e^{}`$ machine where the coupling of the scalar to the $`Z`$ as well as its Yukawa couplings have a bearing on the answer. Thus, model dependence appears only in the decay channels, and hence is easier to analyse. Let us consider, first, the case of monochromatic photon beams. If the center-of-mass energy be $`\sqrt{s}`$, and the product of the photon helicities (circular polarization)<sup>2</sup><sup>2</sup>2We do not consider here the possibility of linear polarization. The most general expression can be found in Ref. . be $`P_{\gamma \gamma }`$, the pair-production cross-section for a scalar of charge $`Q`$ and mass $`m`$ is given by
$$\begin{array}{ccc}\hfill \frac{\mathrm{d}\sigma }{\mathrm{d}\mathrm{\Omega }}& =& \frac{Q^4N_c\alpha ^2}{s}\beta \left[\left(1+P_{\gamma \gamma }\right)\left\{\frac{1}{2}\frac{\beta ^2\mathrm{sin}^2\theta }{1\beta ^2\mathrm{cos}^2\theta }\right\}+\frac{\beta ^4\mathrm{sin}^4\theta }{(1\beta ^2\mathrm{cos}^2\theta )^2}\right].\hfill \end{array}$$
(1)
In eq.(1), $`\beta (14m^2/\widehat{s})`$ is the velocity of the scalars in the center of mass frame and $`\theta `$ the corresponding scattering angle. The colour factor $`N_c=1(3)`$ for sleptons (squarks). The unpolarised cross-section can easily be obtained from the above expression by setting $`P_{\gamma \gamma }=`$ 0. One can check easily that, for small scalar masses, the dominant mode is the $`P_{\gamma \gamma }=1`$ mode. On the other hand, the $`P_{\gamma \gamma }=1`$ mode dominates for scalar masses near the kinematic limit .
In reality though, high energy monochromatic photon beams are extremely unlikely. In fact, the only known way to obtain very high energy photon beams is to induce laser back-scattering off an energetic $`e^\pm `$ beam . The reflected photon beam carries off only a fraction ($`y`$) of the $`e^\pm `$ energy with
$$\begin{array}{ccc}\hfill y_{\mathrm{max}}& =& \frac{z}{1+z}\hfill \\ \hfill z& & \frac{4E_eE_L}{m_e^2}\mathrm{cos}^2\frac{\theta _{eL}}{2}\hfill \end{array}$$
(2)
where $`E_{e(L)}`$ are the energies of the incident $`e^\pm `$ beam and the laser respectively and $`\theta _{eL}`$ is the incidence angle. One can, in principle, increase the photon energy by increasing the energy of the laser beam. However, increasing $`E_L`$ also enhances the probability of electron positron pair creation through laser and scattered-photon interactions, and consequently results in beam degradation. An optimal choice of $`z`$ taking care of this is $`z=2(1+\sqrt{2})`$, and this is the value that we adopt in our analysis.
The cross-sections for a realistic photon-photon collider can then be obtained by convoluting the fixed-energy cross-sections of eq.(1) with the appropriate photon spectrum. For circularly polarized lasers scattering off polarised electron beams, the number-density $`n(y)`$ and average helicity $`\xi (y)`$ for the scattered photons are given by
$$\begin{array}{ccc}\hfill \frac{dn}{dy}& =& \frac{2\pi \alpha ^2}{m_e^2z\sigma _C}C(y)\hfill \\ \hfill \xi (y)& =& \frac{1}{C(y)}\left[P_e\left\{\frac{y}{1y}+y(2r1)^2\right\}P_l(2r1)\left(1y+\frac{1}{1y}\right)\right]\hfill \\ \hfill C(y)& & \frac{y}{1y}+(1y)4r(1r)2P_eP_lrz(2r1)(2y)\hfill \end{array}$$
(3)
where $`ry/z/(1y)`$ and $`\sigma _C`$ is the total Compton cross-section. In the study of polarized beams, one fact needs to be borne in mind. While full (100%) polarisation is possible for a laser, it is unlikely to be realized for electrons. In the rest of this article we shall consider the electron polarization, wherever applicable, to be 90%.
In Fig.1, we present the total cross-section for slepton pair production at a photon collider wherein the parent $`e^+e^{}`$ (or $`e^{}e^{}`$) collider operates at a center of mass energy of $`1\mathrm{TeV}`$. We present the results for three combinations of incident laser and electron polarizations ($`L_1e_1L_2e_2`$). For the entire mass range, at least one of the two polarised ($`++++`$ and $`++`$) cross-sections wins over the unpolarised case. When scalar masses are less than $`230\mathrm{GeV}`$, cross-section for all three cases are comparable with $`\sigma _{++++}>\sigma _{\mathrm{unp}}>\sigma _{++}`$. But for scalar masses above $`230\mathrm{GeV}`$ this hierarchy is just the opposite. In this region, $`\sigma _{++}`$ falls off more slowly than the other two. Depending on $`m_{\stackrel{~}{f}}`$, $`\sigma _{++}`$ can be 5–8 times larger than $`\sigma _{\mathrm{unp}}`$ in this mass range.
## 3 R-parity Violating Decays of the Sfermions.
In this section we present a very brief overview of the phenomenology of $`R`$-parity violation within the MSSM with particular emphasis on the additional decay channels available to the sfermions. As has been noted in the literature, unless a discrete symmetry is introduced explicitly, the superpotential, in addition to the normal Yukawa terms, can also contain the following terms:
$$𝒲_{R/}=\mu _iL_iH_2+\lambda _{ijk}L_L^iL_L^j\overline{E}_R^k+\lambda _{ijk}^{}L_L^iQ_L^j\overline{D}_R^k.+\lambda _{ijk}^{\prime \prime }\overline{U}_R^i\overline{D}_R^j\overline{D}_R^k.$$
(4)
In the above, while $`L_L^i`$ and $`Q_L^i`$ denote the left-handed lepton and quark doublet superfields respectively, $`E_R^i,U_R^i`$ and $`D_R^i`$ denote the corresponding right-handed superfields. The couplings $`\lambda _{ijk}`$ are antisymmetric under the exchange of the first two indices, while $`\lambda _{ijk}^{\prime \prime }`$ are antisymmetric under the exchange of the last two. Since the 36 couplings $`\lambda _{ijk}`$ and $`\lambda _{ijk}^{}`$ violate lepton-number while the other 9 ($`\lambda _{ijk}^{\prime \prime }`$) violate baryon number, simultaneous existence of both the sets of operators would lead to a rapid proton decay and is hence strongly disfavoured. We will thus consider either $`L`$ violating or $`B`$ violating couplings. Furthermore, even within one such subset, non-zero values of more than one coupling could lead to large flavour-changing neutral current amplitudes . We will thus restrict ourselves to cases where only one such coupling is dominant<sup>3</sup><sup>3</sup>3This assumption, of course, prevents us from considering certain spectacular collider signatures..
Examining the individual couplings vis a vis direct $`R\text{/}`$ decays of sfermions, we easily notice:
* the couplings $`\lambda _{ijk}`$ connect either a charged slepton to a lepton-neutrino pair or a sneutrino to two charged leptons. While sneutrino pair production is irrelevant in the context of photon colliders, charged sleptons will lead to a final state with a lepton pair and missing energy-momentum, a state with a large SM background emanating from $`W`$-pair production;
* the couplings $`\lambda _{ijk}^{}`$ lead to both squark-quark-lepton and slepton-quark-quark vertices. While the first resembles a leptoquark vertex, the second (apart from colour factors) mimics a diquark vertex (although there is no violation of baryon number);
* the couplings $`\lambda _{ijk}^{\prime \prime }`$ lead to squark-quark-quark vertices, and again mimic diquarks as far as direct $`R\text{/}`$ decays are concerned.
Thus direct decays through $`\lambda _{ijk}`$ are of no concern to us. Since we shall not address the question of cascading decays in this article, we do not consider such couplings any further. Similarly, as far as direct decays are concerned, the phenomenology of $`\lambda _{ijk}^{\prime \prime }`$ is very similar to that of slepton pair production and subsequent decay through some $`\lambda _{ijk}^{}`$. Hence it suffices to consider the case of a single non-zero $`\lambda _{ijk}^{}`$. Analogous results for $`\lambda _{ijk}^{\prime \prime }`$ can easily be deduced from those that we present.
Expressed in terms of the component fields, the relevant part of the Lagrangian reads
$$\begin{array}{ccc}\hfill _\lambda ^{}& =& \lambda _{ijk}^{}[\stackrel{~}{\nu }_L^i\overline{d}_R^kd_L^j+\stackrel{~}{d}_L^j\overline{d}_R^k\nu _L^i+(\stackrel{~}{d}_R^k)^{}(\overline{\nu }_L^i)^cd_L^j\hfill \\ & & \stackrel{~}{e}_L^i\overline{d}_R^ku_L^j\stackrel{~}{u}_L^j\overline{d}_R^ke_L^i(\stackrel{~}{d}_R^k)^{}(\overline{e}_L^i)^cu_L^j]+h.c\hfill \end{array}$$
(5)
Bounds on these couplings can be obtained from various low-energy observables . These include, for example, meson decay widths , neutrino masses , rates for neutrinoless double beta decay etc. The bounds generally scale with the sfermion mass, and for $`m_{\stackrel{~}{f}}=100\mathrm{GeV}`$ range from $`0.02`$ to $`0.8`$.
In presence of such a $`R\text{/}`$ term, a sfermion can decay into two SM fermions. This mode then competes with the $`R`$-conserving ones namely $`\stackrel{~}{f}f+\chi _i^0,f^{}+\chi _i^\pm `$. The partial decay widths for the latter are of course determined by the masses of the neutralinos (charginos) and their couplings to the sfermion. These quantities, in turn, are determined by the gaugino mass parameters $`M_i`$, the higgsino mass $`\mu `$ and the ratio of the two higgs vacuum expectation values $`\mathrm{tan}\beta `$. However, if decays into the top quark are not allowed, the dependence on the last two parameters is negligible. Furthermore, if we assume gaugino mass unification, as happens in GUTs or supergravity inspired scenarios, only one of the parameters $`M_i`$ is independent. In Fig.2, we plot the contours of constant $`R\text{/}`$-branching ratios for sleptons and up-type squarks in $`M_2`$-$`m_{\stackrel{~}{f}}`$ plane. The particular value of $`\lambda _{ijk}^{}`$ chosen here satisfies the strongest (barring the case of $`\lambda _{111}^{}`$) of the low energy constraints for $`m_{\stackrel{~}{f}}\stackrel{>}{}\mathrm{\hspace{0.25em}100}\mathrm{GeV}`$. To reduce the number of parameters, we assume the GUT relation between the $`U(1)`$ and $`SU(2)`$ gaugino masses but retain the gluino mass as a free parameter. The lower limit on the squark mass has been chosen so as to be consistent<sup>4</sup><sup>4</sup>4In actuality, their bound is somewhat larger than $`200\mathrm{GeV}`$, but has been obtained for a special case. with the bounds quoted by the CDF collaboration .
Let us concentrate first on Fig.2$`a`$, where we consider the case of a slepton and $`\mu =500\mathrm{GeV}`$. The higgsino mass being large, two of the neutralinos (as well as one chargino state) are irrelevant over the entire range of parameter space displayed. The slepton may only decay into the two lighter neutralinos (mainly gauginos) and the lighter chargino. The straight lines thus reflect contours of equal kinematic suppression in the decay widths. The curves change somewhat for a smaller value of $`\mu `$ (see Fig.2$`b`$) since now the neutralino (chargino) sector mixings ensure that more decay channels are available to the slepton (especially if it is on the heavier side).
For smaller masses, the contours for the squarks (Figs.2$`c,d`$) are qualitatively similar. The main difference arises on account of the gluino, which we assume, for the purpose of the graph, to be significantly lighter than that stipulated by gaugino mass unification. Since the decay into gluino proceeds through strong interactions, it dominates almost immediately on crossing the kinematic threshold. It must be pointed out though that, for the stop, this decay (or, for that matter, the decays into neutralinos) almost never reaches the kinematic threshold in the parameter range of interest.
## 4 Signals and Backgrounds
We are now in a position to discuss the signals that we are interested in and the backgrounds thereof. As already discussed, we focus on the direct $`R\text{/}`$–decays of the sfermions. Thus, the sleptons decay into two quarks each, while for the squarks we concentrate on $`l(e,\mu )+jet`$. Thus, the pair produced charged scalars will give rise to $`l^+l^{}`$ plus 2 jet or 4 jet in the final state.
The obvious background consists of the SM process $`\gamma \gamma 4\mathrm{fermions}`$. This has contributions from both ‘resonant’ (such as $`W^+W^{}`$ production) and nonresonant diagrams. Another potential source is heavy quark ($`b,c`$) production followed by decays of these quarks. However, in general such decay products are soft and such backgrounds can be eliminated by imposing simple kinematical cuts. Also to be considered are the contributions from both single– and double–resolved processes. These turn out to be negligible though. The large number of diagrams contributing to the background are calculated using the helicity amplitude package madgraph . To estimate the number of events and their distribution(s), we use a parton-level Monte-Carlo event generator.
### 4.1 Squark Production
We begin with the squarks as the analysis is simpler. To be specific, we consider only the up-type ($`Q=2/3`$) squarks as the production cross-section is 16 times than that for the down-type squarks. The only direct $`R\text{/}`$ decay channel is therefore into a charged lepton and a down-type quark with the resultant final state consisting of two hard jets alongwith two hard leptons.
For given quark and lepton flavours, there are forty SM diagrams contributing to the the process $`\gamma \gamma \mathrm{}^+\mathrm{}^{}q\overline{q}`$. These can be divided into three topological classes:
* 12 of the type $`\gamma \gamma \mathrm{}^+\mathrm{}^{}\gamma ^{}(Z^{})`$ with the off-shell boson emanating from any of the three lepton lines and subsequently going into the quark pair;
* 12 of the type $`\gamma \gamma q\overline{q}\gamma ^{}(Z^{})`$;
* 16 diagrams with a ‘$`t`$-channel’ $`\gamma (Z)`$ exchange (nonresonant topology).
Clearly, for small angle scatterings, each of these diagrams can lead to a large contribution. To eliminate these, we require that both jets and leptons should be relatively central:
$$|\eta _{j,\mathrm{}}|<2.5.$$
(6)
This also ensures that these would be well within the detector. Clearly, the loss in signal would be marginal as the final state there arises from decay of two scalar particles, with even the production cross-section not being strongly peaked. We also must ensure that the jets and the leptons are separated well enough to be detectable as individual entities. To this end, we adapt the well-known cone algorithm for jet separation to a parton-level analysis. Defining $`\mathrm{\Delta }R_{ab}\sqrt{(\mathrm{\Delta }\eta _{ab})^2+(\mathrm{\Delta }\varphi _{ab})^2}`$ where $`\mathrm{\Delta }\eta `$ is the difference of the rapidities of the two particles and $`\mathrm{\Delta }\varphi `$ is their azimuthal separation, we demand that
$$\mathrm{\Delta }R_{jj}0.7,\mathrm{\Delta }R_{jl}>0.5,\mathrm{\Delta }R_{ll}>0.2.$$
(7)
Detectability also requires that these particles (jets) must have a minimum momentum. Over and above this, it should be noted that the signal events would typically be characterized by all the four particles having relatively large transverse momenta ($`p_T`$). On the other hand, the SM background has a large component wherein at least one fermion pair has relatively small $`p_T`$. A cut on the particle momenta is thus called for. We find that rather than imposing the same requirement on all the particles, it is better to order them (leptons and jets individually) according to their $`p_T`$. However, in doing this, one must take into account the detector resolution effects. We incorporate this into our analysis by means of a rather conservative smearing of energies<sup>5</sup><sup>5</sup>5The expected angular resolutions are too fine to be of any concern to us. :
$$\begin{array}{ccccc}\hfill \frac{\delta E_j}{E_j}& =& \frac{0.4}{\sqrt{E_j/1\mathrm{GeV}}}+0.02\hfill & & \mathrm{for}\mathrm{jets},\hfill \\ \hfill \frac{\delta E_{\mathrm{}}}{E_{\mathrm{}}}& =& \frac{0.15}{\sqrt{E_{\mathrm{}}/1\mathrm{GeV}}}+0.01\hfill & & \mathrm{for}\mathrm{leptons}.\hfill \end{array}$$
(8)
We then demand that
$$p_T^{j1},p_T^\mathrm{}1>25\mathrm{GeV},\mathrm{p}_\mathrm{T}^{\mathrm{j2}},\mathrm{p}_\mathrm{T}^\mathrm{}2>20\mathrm{GeV},$$
(9)
where $`j1`$ denotes the jet with larger transverse momentum etc.. Alongwith the separability requirement (eq.7), this also serves to eliminate the bulk of contributions from the $`\gamma \gamma \mathrm{}^+\mathrm{}^{}\gamma ^{}`$ and $`\gamma \gamma q\overline{q}\gamma ^{}`$ diagrams.
As we have discussed before, the background also contains diagrams of the form $`\gamma \gamma f\overline{f}Z^{}`$ with the $`Z`$ going into the other pair of fermions. Eliminating these necessitates that we discard events where either the lepton-lepton or the jet-jet invariant mass is close to $`m_Z`$:
$$m_{\mathrm{}\mathrm{}},m_{jj}[80\mathrm{GeV},100\mathrm{GeV}].$$
(10)
Imposing the above set of cuts, we can then calculate both the signal and background. For the latter, we need to sum over all possible light quarks<sup>6</sup><sup>6</sup>6We do not consider here the special case that the squark decays into a lepton and $`b`$-quark thus making it possible for us to tag the corresponding jet.. Doing this, we obtain, for unpolarized beams and $`\sqrt{s_{ee}}=1\mathrm{TeV}`$, an integrated cross-section of $`\sigma _{\mathrm{Bkgd}}=2.32\mathrm{fb}`$. To further improve the signal-to-noise ratio, we can use a particular feature of the signal event topology. In the limit of infinite momentum resolution, we can, for each of these events, find one particular lepton–jet pairing such that the corresponding invariant masses are exactly the same. Provided such a pairing is unambiguous, this invariant mass is then the mass of the squark. Clearly, the background events would not show the same characteristics, and this can be used to our advantage. We then retain only those events for which
$$\left|M_\mathrm{}j^{(1)}M_\mathrm{}j^{(2)}\right|10\mathrm{GeV}$$
(11)
for at least one such pairing. The corresponding average mass is then treated as our determination of $`m_{\stackrel{~}{q}}`$. With this additional restriction, the background cross-section drops to $`\sigma _{\mathrm{Bkgd}}=0.14\mathrm{fb}`$. It might be argued, and rightly too, that a $`m_{\stackrel{~}{q}}`$ independent criterion as in eq.(11) is not the most efficient one. Indeed, with the fractional energy resolution growing with the jet/lepton energies, we would do better by optimizing this cut for each squark value that we might be interested in. However, we omit to do so in order to keep the analysis a simple one.
With the above set of kinematical cuts, we have been able to reduce the background to insignificant levels. It remains to be seen how much of the signal is retained and whether mass reconstruction is possible. We tackle the questions in the reverse order. The impediments to mass reconstruction come from two sources. The first is of course the effect of the resolution smearing. A second source of ambiguity is the possibility that both set of reconstructions could satisfy eq.(11). In such a case, we retain the pairing with the larger average ($`M_{\mathrm{av}}`$) of the reconstructed masses. As is easily ascertained from Fig.3$`a`$, the mass reconstruction works quite well for relatively smaller values of $`m_{\stackrel{~}{q}}`$. For heavier squarks though, the peaks are not as sharp. This feature is easy to understand. As $`m_{\stackrel{~}{q}}`$ increases, the squarks are produced with smaller and smaller average momentum. Close to the kinematic threshold, the mass difference (eq.11), in the limit of infinite resolution, would vanish identically for both the pairings. Of course, if the resolution was really infinite, our algorithm—namely the larger of the two reconstructed masses—would still make the correct identification. However, because of the resolution smearing, migration of events do occur. Thus, sensitivity could be increased were we to compare, bin by bin, the data and the SM expectations, thereby using a statistical discriminator (such as a $`\chi ^2`$ test). In this article, though, we will not attempt such an analysis. Rather, we effect the somewhat cruder strategy of integrating over five contiguous bins<sup>7</sup><sup>7</sup>7As the lepton and quark momenta are smeared using a Gaussian distribution, the signal invariant mass distribution has not a sharp peak in the relevant mass bin. Rather it also shows a Gaussian structure. To take into account this smearing effect, when calculating the number of signal events for a particular squark or slepton mass, we not only consider the number of events in the relevant mass bin but also add to it the contribution from the four adjacent bins. centering around $`M_{\mathrm{av}}`$ and ascribing any excess therein to a squark of mass $`M_{\mathrm{av}}`$. The partial loss of information inherent in such a procedure makes our results to be somewhat conservative. In Fig.3$`b`$, we present the number of events expected (as a function of $`m_{\stackrel{~}{q}}`$) after the imposition of all of the above cuts (we assume here 100% branching into a lepton and a quark). Comparing it to Fig.1, we see that the loss in signal, while significant, is not debilitating. On the other hand, hardly any background events are expected within a given bin (the surviving background cross-section of $`0.14\mathrm{fb}`$ corresponds to just 7 events over the entire mass range for the assumed luminosity of $`50\mathrm{fb}^1`$). Consequently, observation of even an handful of such events, concentrated within a small mass range, could be construed as the evidence for squark production and subsequent decay through an $`R`$-parity violating interaction, albeit with smaller branching fractions.
As we have commented in Section 3, in the supersymmetric case, the squark may decay through $`R`$-conserving interactions into a quark and a chargino/neutralino. Although the latter will ultimately decay into SM particles through $`R\text{/}`$ interactions, the event shape would be considerably different from that we have considered here. Confining ourselves solely to the analysis of 4-fermion final states, we can use the information of Fig.3$`b`$ to obtain exclusion/discovery limits in the $`m_{\stackrel{~}{q}}`$–branching fraction plane. Since the number of background events is almost zero, the required branching fraction $`Br`$ is approximately given by
$$\mathrm{Br}=\sqrt{\frac{n_{\mathrm{req}}}{N_s}}$$
where $`N_s`$ is the number of signal events corresponding to $`\mathrm{Br}=1`$ and $`n_{\mathrm{req}}=5(3)`$ for discovery (95% exclusion). In Fig. 4 we present the corresponding contours for two different choices of the initial state polarisation. As the polarised $`(++`$) cross-section dominates over the unpolarised for large $`m_{\stackrel{~}{q}}`$, the required branching ratios are consequently smaller. For example, close to the kinematic limit (say $`m_{\stackrel{~}{q}}400`$), we would make a discovery even with a branching fraction 50% provided we work with the correct beam polarisation. On the other hand, with an unpolarised initial state, and $`m_{\stackrel{~}{q}}\stackrel{>}{}\mathrm{\hspace{0.25em}390}\mathrm{GeV}`$, signal events number less than five even for a 100% branching.
### 4.2 Slepton Production
The production cross-section for charged sleptons is higher (by a factor of $`27/16`$) than that for $`up`$-type squarks. However, since they decay into two quarks each, the resultant final state of 4 jets is more complicated than that in the previous subsection. In Table. 1, we list the various subprocesses contributing to this background.
Despite the profusion of diagrams and the attendant complications, certain kinematical cuts suggest themselves. Drawing from the experience of the previous subsection, we demand that
$$|\eta _j|<2.5,$$
(12)
and
$$\mathrm{\Delta }R_{jj}0.7.$$
(13)
An important feature is that the background receives $`𝒪(\alpha \alpha _s)`$ contributions and that many more subprocesses contribute to it as compared to that in the previous subsection. Consequently, the background is much larger and we need to impose somewhat stricter $`p_T`$ requirements. Once again ordering the jets by their transverse momentum, we demand that
$$p_T^{j1},p_T^{j2}>40\mathrm{GeV},\mathrm{p}_\mathrm{T}^{\mathrm{j3}},\mathrm{p}_\mathrm{T}^{\mathrm{j4}}>15\mathrm{GeV}.$$
(14)
This helps us to eliminate the bulk of the $`𝒪(\alpha \alpha _s)`$ contributions. The $`𝒪(\alpha ^2)`$ contributions, on the other hand, are dominated by the resonant contributions. In the present case, these are of two types : ($`i`$) $`\gamma \gamma f\overline{f}Z^{}`$ as before, and ($`ii`$) $`\gamma \gamma W^{}W^{}`$. To eliminate both sets, we demand that
$$m_{ik}[75\mathrm{GeV},95\mathrm{GeV}]$$
(15)
for any of the six pairings. Comparing it to the analogous cut of the last subsection, it might seem that a somewhat harder cut is called for. However, such a course would entail the loss of a significant fraction of the signal as well. Unless we design $`m_\stackrel{~}{\mathrm{}}`$-specific cuts, eq.(15) was found to be an optimal choice. With these set of selection criteria, the processes of Table. 1 lead to a total of $`\sigma _{\mathrm{Bkgd}}170\mathrm{fb}`$ (for $`\sqrt{s_{ee}}=1\mathrm{TeV}`$).
As in the case of the squark, we again take recourse to mass reconstruction, demanding that
$$\left|M_{ij}M_{kl}\right|10\mathrm{GeV}$$
(16)
for at least one of the three pairings. Clearly, the combinatorial factor for the SM background is higher in this case than that for squark decays. Consequently, the reduction ($`\sigma _{\mathrm{Bkgd}}20\mathrm{fb}`$) is not as pronounced. However, with most of this being concentrated at relatively small $`M_{\mathrm{av}}`$ (see Figs.5a,b), where the signal cross-section is on the higher side too, even this background is not a major source of concern. Also shown in Figs.5 are the signal profiles for particular values of $`m_{\stackrel{~}{l}}`$. The mass reconstruction is slightly worse here than that in the previous subsection. This is not unexpected as ($`i`$) jet energy resolution is significantly worse than lepton energy resolution; and ($`ii`$) here three different jet pairings are possible as compared to only two pairings in the other case. Still, there is a significant excess of signal events over the background. Since the latter is rather insensitive to beam polarization, a careful choice of the same can be used to further enhance this excess. An example of this is provided by Fig.5$`b`$.
To maximize sensitivity, we could either opt for mass-dependent kinematical cuts or compare events bin by bin and employ a statistical discriminator such as a $`\chi ^2`$ test. However, once again, we adopt the simpler course of summing over five contiguous bins centered around the slepton mass of interest and compare it with the background. Since the latter is no longer vanishingly small, we cannot simplify our analysis as in the previous subsection. Instead, the required branching fraction is now given by
$$Br.=\sqrt{\frac{n\sqrt{N_b}}{N_s}}$$
where $`N_b`$ and $`N_s`$ are the number of background and signal events (summed over the 5 bins). $`n=5(2)`$ for discovery (95 % C.L. exclusion). Of course, this algorithm is valid strictly in the large $`N_b`$ limit; for small $`N_b`$, we use the appropriate Poisson limit. In Fig.6, we present the contours for two different polarization choices. For small sfermion masses, this channel is less sensitive compared to the one considered in the previous subsection, inspite of the larger production cross-section. This of course can be ascribed primarily to the much larger background and to a smaller extent to the combinatorial problem. For larger sfermion masses though, the background count is $`\stackrel{<}{}𝒪(1)`$ for both cases. Hence the exclusion curves are quite similar in this region. Of course, with a larger integrated luminosity, this argument would cease to hold and the squark production channel would outperform the present one even for $`m_{\stackrel{~}{f}}`$ close to the kinematic limit.
#### 4.2.1 Sleptons decaying into $`b`$ quarks
The above results can be significantly improved if the slepton were to decay into a $`b`$-quark and a light ($`u`$-type) quark. Since $`b`$-jets can be distinguished from those coming from light quarks or gluons, the number of processes contributing to the SM background is reduced considerably. Looking at Table 1, we see that, in the limit of ideal identification <sup>8</sup><sup>8</sup>8This approximation is quite valid given that misidentification probability $`\stackrel{<}{}\mathrm{\hspace{0.25em}0.005}`$ even in a hadronic environment ., only one subprocess each of types (2 & 4) and two subprocesses each of types (6 & 8) may contribute. Consequently, the background, with the an identical set of cuts, is now reduced to $`\stackrel{<}{}\mathrm{\hspace{0.25em}0.1}\mathrm{fb}`$.
This enormous reduction in the background more than offsets the reduction in the signal on account of the less-than-ideal $`b`$-tagging efficiency. For the latter, we use the conservative value $`ϵ_b=0.6`$ per $`b`$-jet . An additional improvement occurs in the mass reconstruction on account of the smaller number of combinatorial possibilities. The situation is thus quite analogous to that of squark production, albeit with a smaller effective production cross-section (the ratio being $`27ϵ_b^2/16`$). That the curves in Fig.7$`a`$ are not exactly parallel to those of Fig.3$`b`$ is due to the facts that the kinematic cuts are not exactly the same and that the energy resolution is slightly worse in the present case. The same also explains the shapes of the exclusion/discovery curves in Fig.7$`b`$.
Interestingly, despite the smaller backgrounds, $`b`$-tagging does not seem to help much for large slepton masses (compare the curves of Fig.6 and Fig.7$`b`$). This is easily understood on noticing the fact that the SM background is essentially zero for such invariant masses. Consequently, the sensitivity of the channel is determined by the signal size alone. With a less-than-ideal efficiency, $`b`$-tagging obviously reduces the signal size without gaining in terms of the background. Foregoing tagging would improve the efficiency (for large $`m_{\stackrel{~}{l}}`$) to the levels of Fig.6, but at the cost of determining the nature of the coupling.
## 5 Bounds on the SUSY Parameter Space
In the previous sections we have presented the reach of a photon-photon collider in a model independent way. In other words, we have not made any assumptions about the decay modes allowed to these scalars. Once the rest of the spectrum in the theory (and their couplings with the scalar in question) is known, the branching fractions to a pair of SM fermions can then be computed in terms of the Yukawa coupling (as, for example, was done in section 3). The branching ratios can then be combined with the exclusion/discovery plots in a straightforward manner to yield the bounds on the relevant parameter space of the theory. For the sake of completeness, we present the outcome of one such exercise here.
As an illustrative example, we choose the case of a slepton decaying into a $`b`$ and a light quark. The low energy bounds only imply
$$\lambda _{113}^{}<0.021\frac{m_{\stackrel{~}{f}}}{100\mathrm{GeV}},$$
will all other relevant $`\lambda _{ij3}^{}`$s allowed to be larger. Thus we may, safely, choose $`\lambda ^{}=0.02`$ over the entire range of interest. In Fig.8, we present the discovery contours that are obtained by combining the results of our simulation (Fig. 7$`b`$) with the branching fractions of Fig. 2$`a`$.
The parameter space bounds as obtained from the other modes discussed in this paper are very similar. The case of the slepton decaying into two non-$`b`$ quarks would, typically, lead to constraints slightly weaker than those of Fig.8. This is particularly true for smaller values of $`m_{\stackrel{~}{l}}`$. However, for slepton masses close to the kinematic limit, $`b`$-tagging actually reduces the sensitivity. Bounds derived from squark production, on the other hand, would depend crucially on the gluino mass. For small squark masses, decays into gluinos are unlikely and hence the bounds would be stronger than the ones described by Fig.8. For large squark masses, on the other hand, the gluino channel might open up thus reducing the $`R\text{/}`$ branching fraction drastically (see Figs.2) and thereby decreasing the efficacy of the search strategy discussed here.
Similar analyses obtain for leptoquarks and diquarks as well. In fact, the case of a generic diquark is almost exactly the same as that for the slepton, but for the modification in the cross section on account of a possibly differing electric charge. A generic leptoquark, on the other hand, provides for an additional two-body decay channel namely into a quark and a neutrino. While such a mode precludes mass reconstruction, it certainly can be used to improve the signal to noise ratio.
## 6 Conclusions
To summarise, we discuss scalar particle pair production at a photon photon collider and their subsequent 2-body decays through $`L`$ violating interactions.
We find that the use of a few well-chosen kinematic cuts can eliminate the SM backgrounds to a considerable extent. Mass reconstruction is possible and is quite accurate almost over the entire kinematic range. For pair-production close to the threshold, the use of polarized electron and laser beams (used to obtain the high energy photons) can increase the production cross-sections manifold without significantly altering the SM background.
As the production cross-section is completely model-independent, the signal strength is a very good measure of the branching fractions into these $`L`$-violating decay modes. Consequently, bounds on the parameter space can be obtained, and we have done so for the case of $`R`$-parity violating supersymmetric models. Although slepton production cross sections are the largest, the corresponding backgrounds are larger too. This renders the discovery reach for this channel to be somewhat worse than that for up-type squarks decaying through a similar $`LQD`$ coupling. However, if the slepton were to decay into a $`b`$-quark, then $`b`$-jet identification could be used to increase the sensitivity significantly and make it competitive with the up-squark channel.
It must be noted though that even if the $`L`$ violating 2-body decay modes be small, the daughter particles from the $`L`$-conserving channels will finally decay through a $`L`$-violating mode. Such a cascading process will leave its own tell tale signature. An analysis, albeit a more complicated one, of the same would only serve to complement the present one. The discovery/exclusion plots presented here are thus only conservative ones and can be improved.
Finally, a study of $`B`$ violating decays is but a straightforward extension of that presented in Section 4.2 of this article.
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# References
On Jordan angles and triangle inequality in Grassmannian
Yurii A.Neretin<sup>1</sup><sup>1</sup>1Supported by grants RFBR–98-01-00303, NWO 047-008-009
## Abstract
Let $`L,M,N`$ be $`p`$-dimensional subspaces in $`^n`$. Let $`\phi _j`$ be the angles between $`L`$ and $`M`$, let $`\psi _j`$ be the angles between $`M`$ and $`N`$, and $`\theta _j`$ be the angles between $`L`$ and $`M`$. Consider the orbit of the vector $`\psi ^p`$ with respect to permutations of coordinates and inversions of axises. Let $`Z`$ be the convex hull of this orbit. Then $`\theta \phi +Z`$. We discuss similar theorems for other symmetric spaces.
We also obtain formula for geodesic distance on any invariant convex Finsler metrics on classical symmetric space
We obtain a version of V.B.Lidskii theorem on spectrum of a sum of matrices for arbitrary classical Riemannian semisimple symmetric space <sup>2</sup><sup>2</sup>2The result of the paper was announced in .
1. Grassmannians. Fix positive integers $`pq`$. Consider the space $`^{p+q}`$ equipped with the standard scalar product. Denote by $`\mathrm{Gr}_{p,q}`$ the set of all $`p`$-dimensional linear subspaces in $`^{p+q}`$. The orthogonal group $`\mathrm{O}(p+q)`$ acts in $`^{p+q}`$ and hence it acts on $`\mathrm{Gr}_{p,q}`$. Obviously,
$$\mathrm{Gr}_{p,q}=\mathrm{O}(p+q)/\mathrm{O}(p)\times \mathrm{O}(q)$$
2. Jordan angles. Let $`L,M\mathrm{Gr}_{p,q}`$. Consider an orthonormal basis $`e_1,\mathrm{},e_pL`$ and an orthonormal basis $`f_1,\mathrm{},f_pM`$. Consider the matrix
$$\mathrm{\Lambda }=\mathrm{\Lambda }[L,M]$$
with the matrix elements $`<e_i,f_j>`$. Denote by
$$\lambda _1\mathrm{}\lambda _p$$
the singular values<sup>3</sup><sup>3</sup>3The singular values of a matrix $`A`$ are the eigenvalues of the matrix $`\sqrt{A^{}A}`$. of the matrix $`\mathrm{\Lambda }`$. Obviously, the numbers
$$\lambda _j=\lambda _j[L,M]$$
don’t depend on choice of the bases $`e_1,\mathrm{},e_pL`$ and $`f_1,\mathrm{},f_pM`$.
Proposition 1. Let $`L,M,L^{},M^{}\mathrm{Gr}_{p,q}`$. The following conditions are equivalent
i) $`\lambda _j[L,M]=\lambda _j[L^{},M^{}]`$ for all $`j`$
ii) There exists an element $`g\mathrm{O}(p+q)`$ such that $`gL=L^{},gM=M^{}`$.
Proof. The statement is obvious.
Proposition 2. Consider $`L,M\mathrm{Gr}_{p,q}`$. There exist orthonormal bases $`e_1,\mathrm{},e_pL`$ and $`f_1,\mathrm{},f_pM`$ such that
$`<e_i,f_j>=0\text{if}ij`$
$`<e_j,f_j>=\lambda _j`$
Proof. The statement is obvious.
Proposition 3.
$$a)\lambda _k[L,M]=\mathrm{max}_{\begin{array}{c}PL\end{array}}\mathrm{min}_{\begin{array}{c}vP\\ v=1\end{array}}\mathrm{max}_{\begin{array}{c}wM\\ w=1\end{array}}<v,w>$$
where the first maximum is given over all $`k`$-dimensional subspaces $`P`$ in $`L`$.
$$b)\lambda _k[L,M]=\mathrm{min}_{\begin{array}{c}QL\end{array}}\mathrm{max}_{\begin{array}{c}vQ\\ v=1\end{array}}\mathrm{max}_{\begin{array}{c}wM\\ w=1\end{array}}<v,w>$$
where the minimum is given over all subspaces $`Q`$ having codimension $`k`$ in $`L`$.
Proof. The statement is a corollary of the standard minimax characterizations of eigenvalues and singular values, see , .
Proposition 4. Denote by $`\mathrm{\Pi }_M`$ the orthogonal projector to the subspace $`M`$. Then the numbers $`\lambda _j[L,M]`$ are the singular values of the operator
$$\mathrm{\Pi }_M:LM$$
.
Proof. The statement is obvious.
Lemma 5. Let $`u_1,u_2,\mathrm{},u_p`$ be a ( nonorthogonal ) basis in L. Let $`v_1,v_2,\mathrm{},v_p`$ be a ( nonorthogonal ) basis in M. Denote by $`U`$ the matrix with the matrix elements $`<u_i,u_j>`$, denote by $`V`$ the matrix with the matrix elements $`<v_i,v_j>`$, denote by $`W`$ the matrix with the matrix elements $`<u_i,v_j>`$. Then the numbers $`\lambda _j^2[L,M]`$ coincides with eigenvalues of the matrix
$$U^1WV^1W^t$$
(1)
Proof. The statement is obvious.
The angles or stationary angles<sup>4</sup><sup>4</sup>4Other terms are complex distance or compound distance. Models of symmetric spaces given in and shows that these invariants for all classical Riemannian symmetric spaces are really angles. (K.Jordan, 1875) $`\psi _1\psi _2\mathrm{}\psi _p`$ between the subspaces $`L,M\mathrm{Gr}_{p,q}`$ are defined by
$$\psi _1=\mathrm{\Psi }_1[L,M]:=\mathrm{arccos}\lambda _1,\mathrm{},\psi _p=\mathrm{\Psi }_p[L,M]:=\mathrm{arccos}\lambda _p$$
Obviously, $`0\psi _j\pi /2`$. We also will use the notation
$$\mathrm{\Psi }[L,M]=(\mathrm{\Psi }_1[L,M],\mathrm{},\mathrm{\Psi }_p[L,M])$$
Remark. For all $`j`$ we have $`\mathrm{\Psi }_j[L,M]=\mathrm{\Psi }_j[M,L]`$.
3. The result of the paper. Denote by $`W_p`$ (Weyl group) the group of all transformations of $`^p`$ generated by permutations of the coordinates and by the transformations
$$(t_1,\mathrm{},t_p)(\sigma _1t_1,\mathrm{},\sigma _pt_p)\text{where}\sigma _j=\pm 1$$
Theorem A. Let $`\mathrm{}(x)`$ be a $`W_p`$-invariant norm in $`^p`$. Then the function
$$d(L,M):=\mathrm{}(\mathrm{\Psi }_1[L,M],\mathrm{},\mathrm{\Psi }_p[L,M])$$
is an $`\mathrm{O}(p+q)`$-invariant metric on $`\mathrm{Gr}_{p,q}`$.
Remark. The geodesic distance in $`\mathrm{Gr}_{p,q}`$ associated with the $`\mathrm{O}(p)\times \mathrm{O}(q)`$-invariant Riemannian metrics is given by the formula
$$\mathrm{dist}(L,M)=\sqrt{\mathrm{\Psi }_1[L,M]^2+\mathrm{}+\mathrm{\Psi }_p[L,M]^2}$$
Theorem B. Let $`L,M,N\mathrm{Gr}_{p,q}`$. Let $`\phi _j=\mathrm{\Psi }_j[L,M]`$, $`\psi _j=\mathrm{\Psi }_j[M,N]`$, $`\theta _j=\mathrm{\Psi }_j[L,N]`$ be the angles. Denote by $`𝒵`$ the convex hull of the $`W_p`$-orbit of the vector $`(\psi _1,\mathrm{},\psi _p)^p`$. Denote by $`𝒰`$ the shift of $`𝒵`$ by the vector $`(\phi _1,\mathrm{},\phi _p)`$.
Then there exists a vector $`(\theta _1^{},\mathrm{},\theta _p^{})𝒰`$ such that the collection of numbers $`(\mathrm{cos}\theta _1^{},\mathrm{},\mathrm{cos}\theta _p^{})`$ coincides up to permutation with the collection of numbers $`(\mathrm{cos}\theta _1,\mathrm{},\mathrm{cos}\theta _p)`$.
4. Infinitesimal angular structure. For $`L\mathrm{Gr}_{p,q}`$ we denote by $`T_L(\mathrm{Gr}_{p,q})`$ the tangent space to $`\mathrm{Gr}_{p,q}`$ at the point $`L`$. It is natural to identify elements $`\xi T_L(\mathrm{Gr}_{p,q})`$ with operators $`H`$ from $`L`$ to the orthogonal complement $`L^{}`$. We denote by
$$\rho _1[L;H]\mathrm{}\rho _p[L;H]$$
the singular values of the operator $`H:LL^{}`$.
Lemma 6. Let $`L,L^{}\mathrm{Gr}_{p,q}`$, $`HT_L(\mathrm{Gr}_{p,q})`$ , $`H^{}T_L^{}(\mathrm{Gr}_{p,q})`$. The following conditions are equivalent
i) $`\rho _j[L;H]=\rho _j[L^{};H^{}]`$ for all $`j`$
ii) There exists an operator $`g\mathrm{O}(p+q)`$ such that $`gL=L^{}`$, $`gH=H^{}`$.
Proof. The statement is obvious.
Remark. The $`\mathrm{O}(p+q)`$-invariant Riemannian metric in $`\mathrm{Gr}_{p,q}`$ is
$$\mathrm{tr}H^tH=\rho _j^2[L;H]$$
5. Relations between angles and infinitesimal angular structure. The following statement is obvious.
Proposition 7. Let $`M(t)`$ be a smooth $`C^{\mathrm{}}`$-curve in $`\mathrm{Gr}_{p,q}`$. Then
$$\underset{\epsilon +0}{lim}\frac{\mathrm{\Psi }_j[M(a+\epsilon ),M(a)]}{\epsilon }=\rho _j[M(a);M^{}(a)]$$
where $`\underset{\epsilon +0}{lim}`$ denotes the right limit at $`0`$.
6. Infinitesimal variation of angles. Let $`L,M\mathrm{Gr}_{p,q}`$. Let $`e_jL`$, $`f_jM`$ be orthonormal bases satisfying the conditions
$`<e_i,f_j>=0\text{if}ij`$
$`<e_j,f_j>=\mathrm{cos}\psi _j`$
Assume $`\psi _j`$ be pairwise different.
Consider an orthonormal basis $`r_1,\mathrm{},r_qM^{}`$ such that for all $`jp`$ the vectors $`e_j`$, $`f_j`$, $`r_j`$ span 2-dimensional plane and $`f_j`$ is situated in the angle between $`e_j`$ and $`r_j`$.
We have
$$e_j=f_j\mathrm{cos}\psi _jr_j\mathrm{sin}\psi _j$$
Let $`H:MM^{}`$ be a tangent vector to $`\mathrm{Gr}_{p,q}`$ at the point $`M`$. Let $`h_{ij}`$ be the matrix elements of $`H`$ in the bases $`f_1,\mathrm{},f_p`$ and $`r_1,\mathrm{},r_q`$.
Consider a $`C^{\mathrm{}}`$-smooth curve $`M(\epsilon )\mathrm{Gr}_{p,q}`$ such that
$$M(0)=H;M^{}(0)=H$$
Proposition 8. Assume $`\psi _00`$, $`\psi _p\pi /2`$, and $`\psi _{j+1}\psi _j`$ for all $`j`$. Then
$$\frac{d}{d\epsilon }\mathrm{\Psi }_j[L,M(\epsilon )]|_{\epsilon =0}=h_{jj}$$
(2)
Proof. Denote by $`f_j(\epsilon )`$ the unique vector in $`M(\epsilon )`$ having the form
$$f_j(\epsilon )=f_j+a_{jk}(\epsilon )r_k$$
Obviously,
$$a_{jk}(\epsilon )=\epsilon h_{jk}+O(\epsilon ^2)$$
For all $`i,jp`$ we have
$$<f_i(\epsilon ),f_j(\epsilon )>=\{\begin{array}{c}O(\epsilon ^2),\text{if}ij\\ 1+O(\epsilon ^2),\text{if}i=j\end{array}$$
and
$$<e_i,f_j(\epsilon )>=\{\begin{array}{c}\epsilon h_{ij}\mathrm{sin}\psi _i+O(\epsilon ^2),\text{if}ij\\ \mathrm{cos}(\psi _j+\epsilon h_{jj})+O(\epsilon ^2),\text{if}i=j\end{array}$$
Now we are ready to wright matrix (1) for the subspaces $`L`$, $`M(\epsilon )`$
$$\left(\begin{array}{cccc}\mathrm{cos}^2(\psi _1+\epsilon h_{11})+O(\epsilon ^2)& O(\epsilon )& \mathrm{}& O(\epsilon )\\ O(\epsilon )& \mathrm{cos}^2(\psi _2+\epsilon h_{22})+O(\epsilon ^2)& \mathrm{}& O(\epsilon )\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ O(\epsilon )& O(\epsilon )& \mathrm{}& \mathrm{cos}^2(\psi _p+\epsilon h_{pp})+O(\epsilon ^2)\end{array}\right)$$
This implies required statement.
Remark. Let us fix $`L\mathrm{Gr}_{p,q}`$.
a) The set of all $`M\mathrm{Gr}_{p,q}`$ such that $`\mathrm{\Psi }_0[L,M]=0`$ has codimension $`qp+1`$.
b) The set of all $`M\mathrm{Gr}_{p,q}`$ such that $`\mathrm{\Psi }_p[L,M]=\pi /2`$ has codimension 1.
c) The set of all $`M\mathrm{Gr}_{p,q}`$ such that $`\mathrm{\Psi }_{j+1}[L,M]=\mathrm{\Psi }_j[L,M]`$ has codimension 2.
7. Preliminaries. A $`p\times p`$ matrix $`A`$ is called bistochastic if for all $`k`$, $`l`$
$$\underset{i}{}a_{ik}=1;\underset{j}{}a_{lj}=1$$
We say that a real $`p\times p`$ matrix $`A`$ is quasistochastic if for all $`k`$, $`l`$
$$\underset{i}{}|a_{ik}|1;\underset{j}{}|a_{lj}|1$$
Proposition 9. (Birkhoff) The set of all bistochastic matrices is the convex hull of matrices of permutations<sup>5</sup><sup>5</sup>5i.e. matrices consisting of 0 and 1 and having strictly one 1 in each column and each row.
See , .
Lemma 10. The set of all quasistochastic matrices is the convex hull of the group $`W_p`$.
Proof. It is sufficient to describe extremal points of the set of all quasistochastic matrices.
a) Obviously, for any extremal point $`A`$
$$\underset{i}{}|a_{ik}|=1;\underset{j}{}|a_{lj}|=1$$
(3)
b) Let a matrix $`A`$ satisfies condition (3). Assume the matrix $`|a_{ij}|`$ be not an extremal point of the set of stochastic matrices. Then $`A`$ is not an extremal point of the set of quasistochastic matrices.
Hence any extremal point of the set of quasistochastic matrices is an element of $`W_p`$.
Example. Let $`U=(u_{ij})`$ and $`V=(v_{ij})`$ be matrices with norm $`1`$. Then the matrix $`W=(u_{ij}v_{ij})`$ is quasistochastic.
Lemma 11.<sup>6</sup><sup>6</sup>6This is a minor variation of Fan Ky theorem, see , . Let $`A`$ be a real $`p\times q`$ matrix. Let $`\lambda =(\lambda _1,\mathrm{},\lambda _p)`$ be its singular values. Then the convex hull of the $`W_p`$-orbit of $`\lambda `$ contains the vector $`(a_{11},\mathrm{},a_{pp})`$.
Proof. Indeed, the matrix $`A`$ can be represented in the form
$$A=U\left(\begin{array}{ccccccc}\lambda _1& 0& \mathrm{}& 0& 0& \mathrm{}& 0\\ 0& \lambda _2& \mathrm{}& 0& 0& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& \lambda _p& 0& \mathrm{}& 0\end{array}\right)V^t;U\mathrm{O}(p),V\mathrm{O}(q)$$
Hence
$$\left(\begin{array}{c}a_{11}\\ a_{22}\\ \mathrm{}\\ a_{pp}\end{array}\right)=\left(\begin{array}{cccc}u_{11}v_{11}& u_{12}v_{12}& \mathrm{}& u_{1p}v_{1p}\\ u_{21}v_{21}& u_{22}v_{22}& \mathrm{}& u_{2p}v_{2p}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ u_{p1}v_{p1}& u_{p2}v_{p2}& \mathrm{}& u_{pp}v_{pp}\end{array}\right)\left(\begin{array}{c}\lambda _1\\ \lambda _2\\ \mathrm{}\\ \lambda _p\end{array}\right)$$
Then we apply Lemma 10.
8. Proof of Theorems A–B. Fix nonnegative numbers
$$a_1\mathrm{}a_p$$
Fix arbitrary orthonormal basis
$$e_1,e_2\mathrm{},e_p,f_1,f_2\mathrm{},f_p,r_1,\mathrm{},r_{qp}^{p+q}$$
(4)
Consider the subspace $`L_a(s)\mathrm{Gr}_{p,q}`$ spanned by the vectors $`v_1(s),\mathrm{},v_p(s)`$ given by the formula
$$v_j(s)=\mathrm{cos}(a_js)e_j+\mathrm{sin}(a_js)f_j$$
We obtain a curve $`L_a(s)`$ in $`\mathrm{Gr}_{p,q}`$.
We say that a curve $`\gamma (t)`$ in $`\mathrm{Gr}_{p,q}`$ is a $`H`$-curve if in some orthonormal basis it has the form $`L_a(s)`$. We say that the numbers $`a_j`$ are the invariants of the $`H`$-curve $`\gamma `$.
We say that points $`L_a(s)`$, $`L_a(t)`$ on a $`H`$-curve are sufficiently near if
$$a_p|st|\pi /2$$
The following statements are obvious
Lemma 12. Consider sufficiently near points $`L(s_1)`$, $`L(s_2)`$, $`L(s_3)`$ on $`H`$-curve. Assume $`s_1<s_2<s_3`$. Then for all $`j`$
$$\mathrm{\Psi }_j[L(s_1),L(s_2)]+\mathrm{\Psi }_j[L(s_2),L(s_3)]=\mathrm{\Psi }_j[L(s_1),L(s_3)]$$
Lemma 13. Let $`L,M\mathrm{Gr}_{p,q}`$ and $`\mathrm{\Psi }_p[L,M]<\pi /2`$. Then there exists the unique $`H`$-curve $`\gamma (s)`$ joining $`L`$, $`M`$ such that $`L`$, $`M`$ are sufficiently near points on $`\gamma (s)`$.
Consider points $`L,M,N\mathrm{Gr}_{p,q}`$ having a general position. Denote by $`\theta _j`$ the angles between $`M`$ and $`N`$. Consider the $`H`$-curve $`\gamma (t)`$ such that $`\gamma (0)=M`$, $`\gamma (1)=N`$ and $`M`$, $`N`$ are sufficiently near points of the curve $`\gamma (t)`$. Then the invariants of the $`H`$-curve $`\gamma (s)`$ are $`\theta _1,\mathrm{},\theta _p`$. We assume, that for each $`s[0,1]`$
$$\mathrm{\Psi }_0[L,M(s)]0,\mathrm{\Psi }_p[L,M(s)]\pi /2$$
(5)
$$\mathrm{\Psi }_j[L,M(s)]\mathrm{\Psi }_{j+1}[L,M(s)]\text{for all}j$$
(6)
Denote by $`𝒵`$ the convex hull of $`W_p`$-orbit of the vector $`\theta `$.
By Proposition 8 and Lemma 11, we have
$$\mathrm{\Psi }[L,\gamma (1/n)]\mathrm{\Psi }[L,M]+\left(\frac{1}{n}+O(\frac{1}{n^2})\right)𝒵,n\mathrm{}$$
In the same way,
$`\mathrm{\Psi }[L,\gamma (2/n)]\mathrm{\Psi }[L,\gamma (1/n)]+\left({\displaystyle \frac{1}{n}}+O({\displaystyle \frac{1}{n^2}})\right)𝒵`$
$`\mathrm{\Psi }[L,\gamma (3/n)]\mathrm{\Psi }[L,\gamma (2/n)]+\left({\displaystyle \frac{1}{n}}+O({\displaystyle \frac{1}{n^2}})\right)𝒵`$
$`........`$
and $`O(1/n^2)`$ are uniform in $`k/n`$. Hence
$$\mathrm{\Psi }[L,\gamma (t)]\mathrm{\Psi }[L,M]+t\left(1+O(\frac{1}{n})\right)𝒵;n\mathrm{}$$
and hence
$$\mathrm{\Psi }[L,\gamma (t)]\mathrm{\Psi }[L,M]+t𝒵;$$
(7)
Assume that there exists the unique value $`\stackrel{~}{s}`$ that doesn’t satisfies the conditions (5)–(6). Consider a small $`\delta `$. Then for $`t>\stackrel{~}{s}`$ we have
$`\mathrm{\Psi }[L,\gamma (\stackrel{~}{s}\delta )]\mathrm{\Psi }[L,M]+(\stackrel{~}{s}\delta )𝒵`$
$`\mathrm{\Psi }[L,\gamma (\stackrel{~}{s}\delta )]\text{is close to}\mathrm{\Psi }[L,\gamma (\stackrel{~}{s}+\delta )]`$
$`\mathrm{\Psi }[L,\gamma (t)]\mathrm{\Psi }[L,\gamma (\stackrel{~}{s}+\delta )]+(t\stackrel{~}{s}+\delta )𝒵`$
and we again obtain (7).
A $`H`$-curve of general position contains a finite number of points $`\stackrel{~}{s}`$ that don’t satisfy conditions (5)–(6). Hence we can repeat our arguments.
This proves Theorem B.
Consider a $`W_p`$-invariant norm $`\mathrm{}()`$ on $`^p`$. By Lemma 11 for any $`x𝒵`$
$$\mathrm{}(x)\mathrm{}(\theta )$$
and this finishes the proof of Theorem A.
Corollary 14. $`H`$-curves are geodesics in any metrics having the form
$$d(L,M)=\mathrm{}(\mathrm{\Psi }_1[L,M],\mathrm{},\mathrm{\Psi }_p[L,M])$$
(8)
Remark. If the sphere $`\mathrm{}(x)=1`$ in $`^n`$ doesn’t contain a segment, then this geodesics is unique.
Corollary 15. Consider the Finsler metric on $`\mathrm{Gr}_{p,q}`$ given by the formula
$$F(L,H)=\mathrm{}(\rho _1[L,H],\mathrm{},\rho _1[L,H]);L\mathrm{Gr}_{p,q},HT_L$$
Then the associated geodesic distance is given by the expression (8).
9. Other symmetric spaces. In it was explained that arbitrary classical compact symmetric space is a Grassmannian in real, complex or quaternionic linear space. This allows to translate literally our results to all classical compact Riemannian symmetric spaces
$`\mathrm{U}(n)\times \mathrm{U}(n)/\mathrm{U}(n);\mathrm{U}(n)/\mathrm{O}(n);\mathrm{U}(2n)/\mathrm{Sp}(n)`$ (9)
$`\mathrm{U}(p+q)/\mathrm{U}(p)\times \mathrm{U}(q);\mathrm{O}(2n)/\mathrm{U}(n);\mathrm{O}(p+q)/\mathrm{O}(p)\times \mathrm{O}(q);\mathrm{O}(n)\times \mathrm{O}(n)/\mathrm{O}(n);`$
$`\mathrm{Sp}(p+q)/\mathrm{Sp}(p)\times \mathrm{Sp}(q);\mathrm{Sp}(n)/\mathrm{U}(n);\mathrm{Sp}(n)\times \mathrm{Sp}(n)/\mathrm{Sp}(n)`$
Three series of the type $`A`$ (i.e (9)) slightly differs from others: we have to replace the group $`W_p`$ by the symmetric group.
In the same way, all classical Riemannian noncompact symmetric spaces are open domains in Grassmannians (see ). This allows to extend our results to all classical Riemannian noncompact symmetric spaces
$`\mathrm{GL}(n,)/\mathrm{U}(n);\mathrm{GL}(n,)/\mathrm{U}(n);\mathrm{GL}(n,)/\mathrm{U}(n)`$
$`\mathrm{U}(p,q)/\mathrm{U}(p)\times \mathrm{U}(q);\mathrm{SO}^{}(2n)/\mathrm{U}(n);\mathrm{O}(p,q)/\mathrm{O}(p)\times \mathrm{O}(q);\mathrm{O}(n,)/\mathrm{O}(n);`$
$`\mathrm{Sp}(p,q)/\mathrm{Sp}(p)\times \mathrm{Sp}(q);\mathrm{Sp}(2n,)/\mathrm{U}(n);\mathrm{Sp}(n,)/\mathrm{Sp}(n)`$
10. Some examples. a) The space $`\mathrm{GL}(n,)/\mathrm{U}(n)`$. We realize points of the space as positive definite $`n\times n`$ complex matrices. The group $`\mathrm{GL}(n,)`$ acts on this space by the transformations
$$LgLg^{},g\mathrm{GL}(n,)$$
The angles<sup>7</sup><sup>7</sup>7hyperbolic angles $`\mathrm{\Psi }_j[L,M]`$ between points $`L`$ and $`M`$ are the solutions of the equation
$$det(Le^\psi M)=0$$
Denote by $`\mathrm{\Psi }[L,M]`$ the vector $`(\mathrm{\Psi }_1[L,M],\mathrm{},\mathrm{\Psi }_n[L,M])`$.
Let $`L`$, $`M`$, $`N`$ be points of our space. Consider all vectors in $`^n`$ that can be obtained from $`\mathrm{\Psi }[L,M]`$ by permutations of coordinates. Denote by $`𝒵`$ their convex hull. Then
$$\mathrm{\Psi }[L,N]\mathrm{\Psi }[L,M]+𝒵$$
(10)
b) Original Lidskii theorem. Consider the space $`𝒮`$ of hermitian $`n\times n`$ matrices. This space also is a (nonsemisimple) symmetric space. The group of isometries is the group of transformations
$$XUXU^{}+A\text{where}U\mathrm{U}(n),A𝒮$$
Let $`X,Y𝒮`$. The analogy of angles are the eigenvalues of $`XY`$. The analogy of Theorem B is the original Lidskii theorem .
The space $`𝒮`$ can be identified with the tangent space to $`\mathrm{GL}(n,)`$ at the point $`1`$. For $`X,Y𝒮`$ we define matrices
$$L=1+\epsilon A,M=1+\epsilon B\mathrm{GL}(n,)/\mathrm{U}(n)$$
where $`\epsilon `$ is small. Then the angles between $`L`$ and $`M`$ have the form $`\epsilon \lambda _j`$, where $`\lambda _j`$ are the eigenvalues of $`XY`$. Hence the inclusion (10) implies Lidskii theorem.
Lidskii theorem on singular values of sum of two matrices corresponds to the triangle inequality in a tangent space to $`\mathrm{U}(p+q)/\mathrm{U}(p)\times \mathrm{U}(q)`$ or $`\mathrm{U}(p,q)/\mathrm{U}(p)\times \mathrm{U}(q)`$
c) The space $`\mathrm{Sp}(2n,)/\mathrm{U}(n)`$. This spaces can be realized as the space of symmetric $`n\times n`$ complex matrices with norm $`<1`$ (see for instance ,5.1,6.3). For two points $`T,S`$ we define the expression
$$\mathrm{\Lambda }[T,S]=(1TT^{})^{1/2}(1TS)(1SS^{})^{1/2}$$
(11)
Let $`\lambda _j`$ be its singular values. Then the hyperbolic angles between $`T`$ and $`S`$ are given by the formula
$$\psi _j=\mathrm{arcosh}\lambda _j$$
(12)
The analogy of Theorem B is given by the formula (10).
d) Arazy norms. Denote by $`V_{fin}`$ the space of finite real sequences $`x=(x_1,\mathrm{},x_N,0,0,\mathrm{})`$. Consider a norm $`\mathrm{}`$ on a space $`V_{fin}`$ satisfying the conditions
$`\mathrm{}`$ is invariant with respect to permutations of coordinates
$`\mathrm{}`$ is invariant with respect to the transformations $`(x_1,x_2,\mathrm{})(\sigma _1x_1,\sigma _2x_2,\mathrm{})`$, where $`\sigma _j=\pm 1`$
– if $`x^{(j)}`$ converges to $`x`$ coordinate-wise and $`\mathrm{}(x^{(j)})`$ converges to $`\mathrm{}(x)`$, then $`\mathrm{}(x^{(j)}x)`$ converges to 0 (an equivalent formulation: the $`\mathrm{}`$-convergence on the sphere $`\mathrm{}(x)=1`$ is equivalent to the coordinate-wise convergence).
Let $`V_{\mathrm{}}`$ is the completion of $`V_{fin}`$ with respect to the norm $`\mathrm{}`$.
A compact operator $`A`$ in a Hilbert space is an element of Arazy class (see ) $`C_{\mathrm{}}`$ if the sequence of its singular values is an element of $`V_{\mathrm{}}`$. Consider a space $`\mathrm{B}_{\mathrm{}}`$ (operator ball) of all compact operators $`T`$ in the Hilbert space satisfying the conditions
$`TC_{\mathrm{}}`$
$`T<1`$, where $`T`$ denotes the standard norm of a operator in a hilbrt space
$`T=T^t`$
We define the angles $`\mathrm{\Psi }_j(T,S)`$ in $`\mathrm{B}_{\mathrm{}}`$ by formulas (11)–(12). We define the distance in $`\mathrm{B}_{\mathrm{}}`$ by
$$d_{\mathrm{}}(T,S)=\mathrm{}(\mathrm{\Psi }_1(T,S),\mathrm{\Psi }_2(T,S),\mathrm{})$$
Proposition 16. a) $`d_{\mathrm{}}(T,S)`$ is a metric.
b) The space $`\mathrm{B}_{\mathrm{}}`$ is complete with respect to the metric $`d_{\mathrm{}}(T,S)`$.
The statement a) can be easyly obtained from Theorem A by a limit considerations. For a proof of the statement B see , 8.6.3.
11. Some references. a) For matrix inequalities see for instance , .
b) The formula for the distance in a symmetric space associated with the invariant Riemannian metrics was obtained in
c) Our Theorem B for the unitary group $`\mathrm{U}(n)=\mathrm{U}(n)\times \mathrm{U}(n)/\mathrm{U}(n)`$ is Nudelman–Shvartsman theorem
d) Generalization of Fan Ky theorem to arbitrary simple Lie algebras was obtained in .
e) Let $`G`$ be a simple Lie group, $`K`$ be its maximal compact subgroup and $`KG/K`$ be the hypergroup of $`K`$-biinvariant subsets in $`G`$ with convolution product. The problem about triangle inequality in Grassmannians is related to a classical problem on structure of the hypergroup $`KG/K`$, see for instance , , .
f) Some nonstandard geometries on groups are discussed in .
g) Some applications of geometry of angles are contained in , 6.3.
h) Conjecture. I think that complete triangle inequality for angles coincides with Horn–Klyachko inequalities , see also and some comments in .
Independent University of Moscow, Bolshoj Vlas’evskij per., 11, Moscow, 121002
neretin@main.mccme.rssi.ru
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# Charge and Transition Form Factors of Light Mesons with Light-Front Quark Model
## I Introduction
Form factors are very important physical quantities in understanding the internal structure of hadrons. In this paper, we study two types of form factors: charge and transition form factors for some light mesons. The former is occurred in the elastic electron meson scattering, in which one off-shell photon exchanges between the electron and one of the quarks in the meson. The latter, on the other hand, comes from the reactions where the meson is produced by one on-shell and one off-shell photons. It is well known that these form factors must be treated with the non-perturbative calculations. There are many different approaches to do that, such as lattice calculations , vector meson dominance (VMD) , perturbative QCD (pQCD) , QCD sum rules , nonlocal quark-pion dynamics , and light-front quark model (LFQM) .
LFQM is the only relativistic quark model in which a consistent and fully relativistic treatment of quark spins and the center-of-mass motion can be carried out. Thus it has been applied in the past to calculate various form factors . This model has many advantages. For example, the light-front wave function is manifestly boost invariant as it is expressed in terms of the momentum fraction variables (in “+” components) in analog to the parton distributions in the infinite momentum frame. Moreover, hadron spin can also be relativistically constructed by using the so-called Melosh rotation. The kinematic subgroup of the light-front formalism has the maximum number of interaction-free generators including the boost operator which describes the center-of-mass motion of the bound state (for a review of the light-front dynamics and light-front QCD, see ). For charge and transition form factors, we concentrate on the space-like region $`q^20`$ ($`q`$ is the momentum transfer). In this region, the so-called Z graph vanishes and only the valence-quark contributes. We take a consistent treatment with the decay constant and the charge and transition form factors in LFQM. Whatever large or small momentum transfer, it must be emphasized that these derivations are applied to all the space-like region. On the other hand, there are some experimental data whcih are concerning about the charge and transition form factors for some light mesons. They will offer some tests of this approach.
The paper is organized as follows. In Sec. II, the basic theoretical formalism is given and the decay constant and the charge and transition form factors are derived for pseudoscalar meson. In Sec. III, some asymptotic behaviors and the numerical results for some light mesons are present and discussed. Finally, a summary is given in Sec. IV.
## II Framework
We will describe in this section the light-front approach for the calculation of the charge and transition form factors for off-shell photons and light mesons. The hadronic matrix elements is evaluated at space-like momentum transfer, namely the region $`q^20`$.
A meson bound state consisting of a quark $`q_1`$ and an antiquark $`\overline{q}_2`$ with total momentum $`P`$ and spin $`S`$ can be written as
$`|M(P,S,S_z)={\displaystyle }`$ $`\{d^3p_1\}\{d^3p_2\}2(2\pi )^3\delta ^3(\stackrel{~}{P}\stackrel{~}{p}_1\stackrel{~}{p}_2)`$ (2)
$`\times {\displaystyle \underset{\lambda _1,\lambda _2}{}}\mathrm{\Psi }^{SS_z}(\stackrel{~}{p}_1,\stackrel{~}{p}_2,\lambda _1,\lambda _2)|q_1(p_1,\lambda _1)\overline{q}_2(p_2,\lambda _2),`$
where $`p_1`$ and $`p_2`$ are the on-mass-shell light-front momenta,
$$\stackrel{~}{p}=(p^+,p_{}),p_{}=(p^1,p^2),p^{}=\frac{m^2+p_{}^2}{p^+},$$
(3)
and
$`\{d^3p\}{\displaystyle \frac{dp^+d^2p_{}}{2(2\pi )^3}},`$ (4)
$`|q(p_1,\lambda _1)\overline{q}(p_2,\lambda _2)=b_{\lambda _1}^{}(p_1)d_{\lambda _2}^{}(p_2)|0,`$ (5)
$`\{b_\lambda ^{}(p^{}),b_\lambda ^{}(p)\}=\{d_\lambda ^{}(p^{}),d_\lambda ^{}(p)\}=2(2\pi )^3\delta ^3(\stackrel{~}{p}^{}\stackrel{~}{p})\delta _{\lambda ^{}\lambda }.`$ (6)
In terms of the light-front relative momentum variables $`(x,k_{})`$ defined by
$`p_1^+=(1x)P^+,p_2^+=xP^+,`$ (7)
$`p_1=(1x)P_{}+k_{},p_2=xP_{}k_{},`$ (8)
the momentum-space wave-function $`\mathrm{\Psi }^{SS_z}`$ can be expressed as
$$\mathrm{\Psi }^{SS_z}(\stackrel{~}{p}_1,\stackrel{~}{p}_2,\lambda _1,\lambda _2)=R_{\lambda _1\lambda _2}^{SS_z}(x,k_{})\varphi (x,k_{}),$$
(9)
where $`\varphi (x,k_{})`$ describes the momentum distribution of the constituents in the bound state, and $`R_{\lambda _1\lambda _2}^{SS_z}`$ constructs a state of definite spin ($`S,S_z`$) out of light-front helicity ($`\lambda _1,\lambda _2`$) eigenstates. Explicitly,
$$R_{\lambda _1\lambda _2}^{SS_z}(x,k_{})=\underset{s_1,s_2}{}\lambda _1|_M^{}(1x,k_{},m_1)|s_1\lambda _2|_M^{}(x,k_{},m_2)|s_2\frac{1}{2}s_1\frac{1}{2}s_2|SS_z,$$
(10)
where $`|s_i`$ are the usual Pauli spinors, and $`_M`$ is the Melosh transformation operator:
$$_M(x,k_{},m_i)=\frac{m_i+xM_0+i\stackrel{}{\sigma }\stackrel{}{k}_{}\times \stackrel{}{n}}{\sqrt{(m_i+xM_0)^2+k_{}^2}},$$
(11)
with $`\stackrel{}{n}=(0,0,1)`$, a unit vector in the $`z`$-direction, and
$$M_0^2=\frac{m_1^2+k_{}^2}{(1x)}+\frac{m_2^2+k_{}^2}{x}.$$
(12)
In practice it is more convenient to use the covariant form for $`R_{\lambda _1\lambda _2}^{SS_z}`$ :
$$R_{\lambda _1\lambda _2}^{SS_z}(x,k_{})=\frac{\sqrt{p_1^+p_2^+}}{\sqrt{2}\stackrel{~}{M}_0}\overline{u}(p_1,\lambda _1)\mathrm{\Gamma }v(p_2,\lambda _2),$$
(13)
where
$`\stackrel{~}{M}_0\sqrt{M_0^2(m_1m_2)^2},`$ (14)
$`\mathrm{\Gamma }=\gamma _5(\mathrm{pseudoscalar},S=0).`$ (15)
We normalize the meson state as
$$M(P^{},S^{},S_z^{})|M(P,S,S_z)=2(2\pi )^3P^+\delta ^3(\stackrel{~}{P}^{}\stackrel{~}{P})\delta _{S^{}S}\delta _{S_z^{}S_z},$$
(16)
so that
$$\frac{dxd^2k_{}}{2(2\pi )^3}|\varphi (x,k_{})|^2=1.$$
(17)
In principle, the momentum distribution amplitude $`\varphi (x,k_{})`$ can be obtained by solving the light-front QCD bound state equation. However, before such first-principle solutions are available, we would have to be contented with phenomenological amplitudes. One example that has been often used in the literature for heavy mesons is the Gaussian-type wave function,
$$\varphi (x,k_{})_\mathrm{G}=𝒩\sqrt{\frac{dk_z}{dx}}\mathrm{exp}\left(\frac{\stackrel{}{k}^2}{2\omega ^2}\right),$$
(18)
where $`𝒩=4(\pi /\omega ^2)^{3/4}`$ and $`k_z`$ is of the internal momentum $`\stackrel{}{k}=(\stackrel{}{k}_{},k_z)`$, defined through
$$1x=\frac{e_1k_z}{e_1+e_2},x=\frac{e_2+k_z}{e_1+e_2},$$
(19)
with $`e_i=\sqrt{m_i^2+\stackrel{}{k}^2}`$. We then have
$`M_0=e_1+e_2,k_z={\displaystyle \frac{xM_0}{2}}{\displaystyle \frac{m_2^2+k_{}^2}{2xM_0}},`$ (20)
and
$$\frac{dk_z}{dx}=\frac{e_1e_2}{x(1x)M_0}$$
(21)
which is the Jacobian of transformation from $`(x,k_{})`$ to $`\stackrel{}{k}`$. This wave function has been also used in many other studies of hadronic transitions. A variant of the Gaussian-type wave function is
$`\stackrel{~}{\varphi }(x,k_{})_G=𝒩\sqrt{{\displaystyle \frac{dk_z}{dx}}}\mathrm{exp}\left({\displaystyle \frac{M_0^2}{2\omega ^2}}\right),`$ (22)
with $`M_0`$ being given by (12). This wave function is equivalent to $`\varphi (x,k_{})_\mathrm{G}`$ when the constituent quark masses are equal; otherwise, the results will be different. In this paper, we will assume that the $`u`$ and $`d`$ quarks in pion have the same masses. There is another wave function
$`\varphi (x,k_{})_M=𝒩^{}\mathrm{exp}\left({\displaystyle \frac{M_0^2}{8\omega ^2}}\right),`$ (23)
which will be also used in the numerical calculations.
### A Decay Constants
The decay constant of a pseudoscalar meson $`P(q_1\overline{q}_2)`$ is defined by
$`0|A^\mu |P(P)=i\sqrt{2}f_PP^\mu ,`$ (24)
where $`A^\mu `$ is the axial vector current. It can be evaluated by using the light-front wave function given by (2.1) and (2.5)
$`0|\overline{q}_2\gamma ^+\gamma _5q_1|P`$ $`=`$ $`{\displaystyle \{d^3p_1\}\{d^3p_2\}2(2\pi )^3\delta (\stackrel{~}{P}\stackrel{~}{p_1}\stackrel{~}{p_2})\varphi _P(x,k_{})R_{\lambda _1\lambda _2}^{00}(x,k_{})}`$ (26)
$`\times 0|\overline{q}_2\gamma ^+\gamma _5q_1|q_1\overline{q}_2.`$
Since $`\stackrel{~}{M}_0\sqrt{x(1x)}=\sqrt{𝒜^2+k_{}^2}`$, it is straightforward to show that
$`f_P=\mathrm{\hspace{0.17em}2}\sqrt{3}{\displaystyle \frac{dxd^2k_{}}{2(2\pi )^3}\varphi _P(x,k_{})\frac{𝒜}{\sqrt{𝒜^2+k_{}^2}}},`$ (27)
where
$`𝒜=m_1x+m_2(1x).`$ (28)
Note that the factor $`\sqrt{3}`$ in (27) arises from the color factor implicitly in the meson wave function. We illustrate this process in Fig.1 (a). When the decay constant is known, it can be used to constrain the parameters of the light-front wave function.
### B Charge Form Factors
The charge form factor of a pseudoscalar meson $`P`$ is determined by the scattering of one virtual photon and one meson. We illustrate this process in Fig.1 (b). This form factor can be defined by the matrix element
$`P(P^{})|J^\mu |P(P)=F_P(Q^2)(P+P^{})^\mu `$ (29)
where $`J^\mu `$ is the vector current and $`Q^2q^2=(PP^{})^2`$. As discussed in the above subsection, we readily obtain
$`P(P^{})|\overline{q}\gamma ^\mu q|P(P)`$ $`=`$ $`{\displaystyle \underset{\lambda _1,\lambda _1^{},\lambda _2,\lambda _2^{}}{}}{\displaystyle \{d^3p_1\}\{d^3p_2\}2(2\pi )^3\delta (pp_1p_2)}`$ (31)
$`\varphi _P^{}(x^{},k_{}^{})\varphi _P(x,k_{})R_{\lambda _1^{}\lambda _2^{}}^{00}R_{\lambda _1\lambda _2}^{00},`$
where $`k_{}^{}k_{}+xq_{}`$. After comparing (29) with (31), we can get
$`F_P(Q^2)={\displaystyle \frac{dxd^2k_{}}{2(2\pi )^3}\varphi _P^{}(x^{},k_{}^{})\varphi _P(x,k_{})\frac{\stackrel{~}{M}_0}{\stackrel{~}{M^{}}_0}\left(1+\frac{xq_{}k_{}}{𝒜^2+k_{}^2}\right)}.`$ (32)
Thus, if we can determine the parameters in the wave function, we will have the charge form factor $`F_P(Q^2)`$ in terms of (32).
### C Transition Form Factors
There are two types of the transition form factors: $`F_{\pi \gamma }`$ and $`F_{\pi \gamma }`$. The form factor $`F_{P\gamma }`$, in which the meson is produced by one on-shell and one off-shell photon $`(\gamma \gamma ^{}P)`$, is defined by the $`P\gamma \gamma ^{}`$ vertex
$`\mathrm{\Gamma }_\mu =ie^2F_{P\gamma }(Q^2)\epsilon _{\mu \nu \rho \sigma }P^\nu q^\rho ϵ^\sigma ,`$ (33)
where $`q`$ is the momentum of the off-shell photon, $`q^2=q_{}^2=Q^2`$, and $`ϵ`$ is the polarization vector of the on-shell photon. We illustrate this process in Fig.1 (c) and write down the amplitude in the light-front framework
$`\mathrm{\Gamma }_\mu `$ $`=`$ $`{\displaystyle \underset{\lambda _1,\lambda _2,\lambda }{}}e_qe_{\overline{q}^{}}e^2{\displaystyle \{d^3p_1\}\{d^3p_2\}2(2\pi )^3\delta (\stackrel{~}{P}\stackrel{~}{p_1}\stackrel{~}{p_2})\varphi _P(x,k_{})}`$ (34)
$`\times `$ $`[({\displaystyle \frac{q_{}^2}{p^+}}{\displaystyle \frac{m_1^2+(k_{}+q_{})^2}{p_1^+}}{\displaystyle \frac{m_2^2+k_{}^2}{p_2^+}})^1\overline{v}(p_2,\lambda _2)\overline{)}ϵu(p_1^{},\lambda )\overline{u}(p_1^{},\lambda )\gamma _\mu u(p_1,\lambda _1)`$ (36)
$`+(12)]R^{00}_{\lambda _1\lambda _2},`$
where $`e_q`$ and $`e_{\overline{q}^{}}`$ are the electric charge of $`q`$ and $`\overline{q}^{}`$ quarks, respectively. It is straightforward to show that
$`F_{P\gamma }(Q^2)`$ $`=`$ $`4{\displaystyle \frac{\sqrt{3}}{\sqrt{2}}}e_qe_{\overline{q}^{}}{\displaystyle \frac{dxd^2k_{}}{2(2\pi )^3}\varphi _P(x,k_{})\frac{𝒜}{\sqrt{𝒜^2+k_{}^2}}}`$ (37)
$`\times `$ $`\left[{\displaystyle \frac{1}{(1x)\left(q_{}^2\frac{m_1^2+(k_{}+q_{})^2}{1x}\frac{m_2^2+k_{}^2}{x}\right)}}+{\displaystyle \frac{1}{x\left(q_{}^2\frac{m_1^2+k_{}^2}{1x}\frac{m_2^2+(k_{}q_{})^2}{x}\right)}}\right]`$ (38)
The form factor $`F_{P\gamma ^{}}`$ arising from the $`P\gamma ^{}\gamma ^{}`$ vertex, where $`\gamma \gamma `$ represents two off-shell photons is defined by
$`\mathrm{\Gamma }_{\mu \nu }=ie^2F_{P\gamma ^{}}(Q^2,Q^2)\epsilon _{\mu \nu \rho \sigma }𝒬^\rho 𝒫^\sigma ,`$ (39)
where $`𝒬\frac{1}{2}(q^{}q)`$, $`𝒫q^{}+q`$, and $`Q^2=q^2=q_{}^2`$. This process is illustrated in Fig.1 (d). We expect that this amplitude is similar to the one off-shell photon case and we have
$`\mathrm{\Gamma }_{\mu \nu }`$ $`=`$ $`{\displaystyle \underset{\lambda _1,\lambda _2,\lambda }{}}e_qe_{\overline{q}^{}}e^2{\displaystyle \{d^3p_1\}\{d^3p_2\}2(2\pi )^3\delta (pp_1p_2)\varphi _P(x,k_{})}`$ (40)
$`\times `$ $`[({\displaystyle \frac{q_{}^2q_{}^2}{p^+}}{\displaystyle \frac{m_1^2+(k_{}+q_{})^2}{p_1^+}}{\displaystyle \frac{m_2^2+k_{}^2}{p_2^+}})^1\overline{v}(p_2,\lambda _2)\gamma _\nu u(p_1^{},\lambda )\overline{u}(p_1^{},\lambda )\gamma _\mu u(p_1,\lambda _1)`$ (42)
$`+(12)]R^{00}_{\lambda _1\lambda _2}.`$
From (39) and (42), we arrive at
$`F_{P\gamma ^{}}(Q^2,Q^2)`$ $`=`$ $`4{\displaystyle \frac{\sqrt{3}}{\sqrt{2}}}e_qe_{\overline{q}^{}}{\displaystyle \frac{dxd^2k_{}}{2(2\pi )^3}\varphi _P(x,k_{})\frac{𝒜}{\sqrt{𝒜^2+k_{}^2}}}`$ (43)
$`\times `$ $`[{\displaystyle \frac{1}{(1x)\left(q_{}^2q_{}^2\frac{m_1^2+(k_{}+q_{})^2}{1x}\frac{m_2^2+k_{}^2}{x}\right)}}`$ (45)
$`+{\displaystyle \frac{1}{x\left(q_{}^2q_{}^2\frac{m_1^2+k_{}^2}{1x}\frac{m_2^2+(k_{}q_{})^2}{x}\right)}}].`$
## III Numerical Results and Discussions
We now compare our results of the form factors with the experimental data. Before doing that, we must determine the parameters $`m_1`$, $`m_2`$, and $`\omega `$ in the wave function $`\varphi _P(x,k_{})`$. Of course, we assume that this wave function is process-independent.
In the $`\pi \gamma `$ case, the constituent masses of the $`u`$ and $`d`$ quarks are the same, i.e., $`m_1=m_2m_q`$. We can use the experimental value of decay constant $`f_\pi =92.4`$ MeV to determine the parameters $`m_q`$ and $`\omega _\pi `$ by (27). However, there are two parameters with only one experimental value. Therefore, principally, we can get infinite combinations which all satisfy the decay constant value. If we can find another constraint, these parameters wull be deterimined uniquely. From the transition form factor $`F_{\pi \gamma }`$, we have another constraint. In (38), if we consider the limit $`Q^2\mathrm{}`$, there is a simple form
$`F_{\pi \gamma }(Q^2\mathrm{})=4{\displaystyle \frac{\sqrt{3}}{\sqrt{2}}}{\displaystyle \frac{(e_u^2e_d^2)}{\sqrt{2}}}{\displaystyle \frac{dxd^2k_{}}{2(2\pi )^3}\varphi _\pi (x,k_{})\frac{𝒜}{\sqrt{𝒜^2+k_{}^2}}\left(\frac{1}{x(1x)Q^2}\right)}.`$ (46)
From , we have
$`Q^2F_{\pi \gamma }(Q^2)|_{Q^2\mathrm{}}=6(e_u^2e_d^2)f_\pi .`$ (47)
Comparing (46) with (47), we obtain
$`f_\pi ={\displaystyle \frac{\sqrt{3}}{3}}{\displaystyle \frac{dxd^2k_{}}{2(2\pi )^3}\frac{\varphi _\pi (x,k_{})}{x(1x)}\frac{𝒜}{\sqrt{𝒜^2+k_{}^2}}}.`$ (48)
From (27) and (48), we can uniquely determine all the parameters in the wave function by using only one experimental value $`f_\pi `$. Here we show the parameters of two wave functions $`\varphi _G`$ and $`\varphi _M`$ fitted to the decay constants given by (27) and (48)($`\varphi _G=\stackrel{~}{\varphi }_G`$ in pion case) as
$`\varphi _G`$ $`:`$ $`m_g=0.243\text{GeV},\omega _\pi =0.328\text{GeV};`$ (49)
$`\varphi _M`$ $`:`$ $`m_g=0.198\text{GeV},\omega _\pi =0.513\text{GeV}.`$ (50)
We use the Gaussian-type wave function to calculate the form factors because the value $`m_q`$ of the wave function $`\varphi _M`$ seems to be too small. Moeover, we can evaluate $`F_\pi (Q^2)`$ and $`F_{\pi \gamma }(Q^2)`$ in all momentum transfer region by using (32), (38). Because both parameters which we needed have been fixed, we have no degree of freedom to adjust this wave function. Thus, whether these preditions are consistent with the experiments or not are very strict tests for the Gaussian-type model. From Fig.2 and Fig.3, we find that these predictions are in good agreement with the experimental data .
For $`F_{\pi \gamma ^{}}(Q^2,Q^2)`$, since there are no experimental data yet, we must proceed carefully. If we take both limits of $`Q^2,Q^2\mathrm{}`$, (45) becomes
$`F_{\pi \gamma ^{}}(Q^2,Q^2)|_{Q^2,Q^2\mathrm{}}`$ $`=`$ $`{\displaystyle \frac{2}{\sqrt{3}}}{\displaystyle \frac{dxd^2k_{}}{2(2\pi )^3}\varphi _\pi (x,k_{})\frac{𝒜}{\sqrt{𝒜^2+k_{}^2}}}`$ (52)
$`\times \left({\displaystyle \frac{1}{xQ^2+(1x)Q^2}}+{\displaystyle \frac{1}{xQ^2+(1x)Q^2}}\right).`$
Noting that, the wave function $`\varphi _\pi (x,k_{})`$ is symmetric in $`x`$ and $`1x`$, we get the asymptoyic behavior of the transition form factor as
$`F_{\pi \gamma ^{}}(Q^2,Q^2)|_{Q^2,Q^2\mathrm{}}={\displaystyle \frac{4}{\sqrt{3}}}{\displaystyle \frac{dxd^2k_{}}{2(2\pi )^3}\varphi _\pi (x,k_{})\frac{𝒜}{\sqrt{𝒜^2+k_{}^2}}\left(\frac{1}{xQ^2+(1x)Q^2}\right)},`$ (53)
which is consistent with the assumption in . Thus we have the confidence to make the prediction about the values of $`F_{\pi \gamma ^{}}(Q^2,Q^2)`$ in terms of (45).
We also consider the $`\eta \eta ^{}`$ system. Due to the mixing in this system, $`\eta `$ and $`\eta ^{}`$ both have $`\eta _8`$ and $`\eta _0`$ components. These two states
$`|\eta _8`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{6}}}|\overline{u}u+\overline{d}d2\overline{s}s`$ (54)
and
$`|\eta _0`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{3}}}|\overline{u}u+\overline{d}d+\overline{s}s`$ (55)
are the $`SU(3)`$ octet and singlet, respectively. The decay constants of octet and singlet, $`f_P^a`$ are defined as
$`0|J_{\mu 5}^a|P(p)=\sqrt{2}if_P^ap_\mu ,(a=8,1;P=\eta ,\eta ^{}),`$ (56)
where $`J_{\mu 5}^a`$ denotes axial-vector current. Recent investigations have shown that the decay constants of the $`\eta \eta ^{}`$ system are not adequately described by only one mixing angle. A two-mixing-angle parametrization is given by
$`f_\eta ^8=f_8\mathrm{cos}\theta _8,f_\eta ^0=f_0\mathrm{sin}\theta _0,`$ (57)
$`f_\eta ^{}^8=f_8\mathrm{sin}\theta _8,f_\eta ^{}^0=f_0\mathrm{cos}\theta _0,`$ (58)
where $`\theta _8\theta _0`$. From the phenomenological analysis , we get the values
$`f_81.26f_\pi ,\theta _821.2^{},`$ (59)
$`f_01.17f_\pi ,\theta _09.2^{}.`$ (60)
Using the values of $`f_8`$ and $`f_0`$, we can deterimine the parameters by (27) and (48)
$`m_8=0.306\text{GeV}\omega _{\eta _8}=0.414\text{GeV},`$ (61)
$`m_0=0.285\text{GeV}\omega _{\eta _0}=0.384\text{GeV},`$ (62)
where $`m_8`$ and $`m_0`$ are the parameters of quark masses in $`\eta _8`$ and $`\eta _0`$, respectively. With these parameters in (62) and the values of mixing angles in (60), we could calculate $`F_{\eta \gamma }`$ and $`F_{\eta ^{}\gamma }`$ and the results are ploted in Fig.4.
## IV Summary
The charge and transition form factors of $`\pi `$ have been studied in the present paper. In the light-front relativistic quark model, these form factors have been evaluated in a frame where $`q^+=0`$ and $`q^20`$ and there is no need to calculate the contribution from the so-called $`Z`$ graph .
We have only used one experimental value, the pion decay constant, to fix the two parameters in the pion wave function. This point is in contrast to pQCD which treats the decay constant as one part of the wave function. Thus, we would emphasize that the wave function contains no more degree of freedoms to adjust. When the parameters are fixed, we evaluate the charge as well as one and two virtual photon transition form factors in $`8\text{G}eV^2q^2<0`$. Our calculations are based on the important assumption that: the wave functions are independent of processes. We compare the results of calculation with the experimental data and find that this assumption is valid for $`8\text{G}eV^2q^2<0`$ region. Basing on these consistency of $`F_\pi `$ and $`F_{\pi \gamma }`$, we have the confidence to make the prediction of the $`F_{\pi \gamma ^{}}`$. The decay constants $`f_{\eta _8}`$, $`f_{\eta _0}`$, and the mixing angles $`\theta _8`$, $`\theta _0`$ have been obtained by using the phenomenological analysis. With the same approach, we have gotten $`F_{\eta \gamma }`$ and $`F_{\eta ^{}\gamma }`$ which agree well with the experimental data.
ACKNOWLEDGMENTS
We are grateful to T. Feldmann and H.Y. Cheng for the helpful discussions on the $`\eta \eta ^{}`$ mixing system. This work was supported in part by the National Science Council of ROC under Contract Nos. NSC89-2112-M-009-035.
FIGURE CAPTIONS
Fig. 1 The diagram of (a) one pseudoscalar meson decay to vacuum, (b) the scattering of one virtual photon and one meson, (c) a meson is produced by one on-shell and one off-shell photons, and (d) a meson is produced by two off-shell photons.
Fig. 2 The charge form factor of pion in small and large momentum transfer. Data are taken from and , respectively.
Fig. 3 The one off-shell photon transition form factor of pion. The solid line represents the results obtained with this approach. The dotted line represents the limiting behavior $`2f_\pi `$ ($`0.185`$ GeV). Data are taken from .
Fig. 4 The one off-shell photon transition form factor of $`\eta `$ and $`\eta ^{}`$. The dotted line represents the limiting behavior $`0.182`$ GeV and $`0.300`$ GeV, respectively. Data are both taken from .
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# General Relativistic Treatment of the Thermal, Magnetic and Rotational Evolution of Isolated Neutron Stars with Crustal Magnetic Fields
## 1 Introduction
The question whether and, if so, how the magnetic field of isolated neutron stars (NSs) decays is a controversial issue and a subject of hot scientific debates. The observed rotational periods $`P`$ and time derivative $`\dot{P}`$ of about 700 pulsars (PSRs) (see PSR catalogue, Taylor et al. 1993TML93 ) and studies of their inferred surface magnetic field strength versus their active age ($`\tau _a`$) provide evidence that the magnetic fields of NSs is subject to decay. This evidence is rather strong for the old population of PSRs, i.e. $`\tau _a100`$ Myrs, while for the younger population evidence for magnetic field decay is much weaker. As investigated by population synthesis methods (Bhattacharya et al. 1992BWHV92 , Hartman et al. 1997HBWV97 ), observations and models are in harmony provided one accepts the hypothesis that NS magnetic field decays very little during the first $`\tau _a10`$ Myrs. The NS magnetic field decay depends strongly on where is the field located in the NS, what is its structure and strength. It is also related to the equation of state (EOS) of NS matter and the conductive properties of dense matter which are moreover affected by its thermal history too. In case the field is of fossil origin, or alternatively has been generated via a dynamo action during the proto–NS phase (Thompson & Duncan 1993TD93 ), it is expected to thread most of the star. Early estimates of the electric conductivity of NS matter lead to the conclusion that magnetic field of NSs will not be dissipated during a Hubble time (Baym, Pethick & Pines 1969BPP69 ). However, more recent investigations regarding the conductive properties of nuclear matter in the presence of a strong B-field yield conductivities leading to much shorter field decay times (see e.g. Haensel, Urpin & Yakovlev 1990HUY90 , Goldreich & Reissenegger 1992GR92 , Urpin & Shalybkov 1995US95 , Shalybkov & Urpin 1997SU97 ). Besides the conductive properties of nuclear matter, another distinct mechanism leading to a B-decay process is based on the idea of magnetic flux expulsion from NS core driven by rotation or/and buoyancy (Muslimov & Tsygan 1985MT85 ; Srinivasan et al 1990SBMT90 ). Once the expelled field reaches the star’s crust it subsequently suffers Ohmic decay. Ding et al. (1993)DCC93 considered in detailed this process and have found typical decay times of the order of $`100`$ Myrs. Recent investigations, which take into account the effect of the NS crust onto the process of flux expulsion (Konenkov & Geppert 2000KG00 ) estimated even longer decay times for a field anchored in fluxoids in the superfluid NS core.
Although the issue of whether the field is penetrating the entire star or part of it is an open one, there exist, however, good reasons to believe that the NS magnetic field is maintained by currents in the crust. One possibility is the generation of the field via thermomagnetic effects during the first years of the NS life (Blandford et al. 1983BAH83 , Urpin, Levshakov & Yakovlev 1986ULY86 , Wiebicke & Geppert 1996WG96 ). The most recent and detailed investigation of the magnetic and spin evolution of isolated NSs with a crustal dipolar field has been performed by Urpin & Konenkov (1997UK97 , UK97 thereafter). They found a good agreement with observational data provided the EOS is not too soft, the initial surface magnetic field strength lies in the range of $`10^{12}`$ to $`3\times 10^{13}`$ G and is initially confined to densities of $`10^{12}`$ to $`10^{13}`$ g cm<sup>-3</sup> within the crust. Miralles, Urpin & Konenkov (1998)MUK98 considered the effect of a crustal field decay onto the thermal evolution of a NS and they have concluded that a considerable amount of heating takes place in the the crust after about $`3`$ to $`10`$ Myrs , a period during which the NS has cooled down to $`10^5`$ K. Consequently, Joule heating can maintain a warm ($`510^4`$ K) NS surface for a period of hundreds of Myrs.
All of the above described works ignore the curvature of spacetime. To date there exist no calculation treating self consistently the thermal and magnetic evolution of NSs which incorporate General Relativistic (GR) effects. As it has been argued elsewhere by the present authors (Geppert, Page & Zannias 2000GPZ00 , referred to as GPZ00 hereafter), relativistic effects on the B-field evolution must be included in detailed investigations. A previous attempt to incorporate GR effects has been presented by Sengupta (1997S97 , 1998S98 ). However this author failed to take into account the proper boundary condition associated with relativistic treatment and furthermore his formalism applies only to a Schwarzschild geometry. Moreover, he claims that the decay rate of the field is decreased by several orders of magnitude, a conclusion being in variance with the one obtain by GPZ00GPZ00 and the present work.
In the present paper we investigate in details and in a self consistent manner, GR effects at first on both: the thermal and magnetic evolution of NSs. We consider three different EOSs and via numerical integration of Einstein’s equations, neutron star models are constructed characterized by the compactness ratio varying in a large range. We solve simultaneously the relevant evolution equations and thus our approach naturally reveals the mutual dependencies of EOS, mass–to–radius relation, initial field structure and strength, and thermal history of NSs. In particularly, as far as the cooling is concerned we consider two scenarios: the ‘standard’ slow neutrino emission scenario and also the so called enhanced neutrino emission which results in a much lower temperatures in young NSs. In order to avoid dealing with the uncertainties related to the behavior of the field within the superconducting core, as a first step in the present paper, we consider only magnetic fields not penetrating the NS core. Thus, the present analysis is within the framework of the crustal magnetic field hypothesis. We use our results from the GR treatment of thermal and magnetic field evolution implemented by a semi-relativistic treatment of the spin evolution to confront the crustal field hypothesis to observational data.
The paper is organized as follows: In the subsequent section we remind the reader of the GR formulation of the equations of stellar structure, the heat transport and conservation equations as well as the induction equation on a static spherically symmetric background geometry. We also present the evolution equations for an axisymmetric dipolar magnetic field as well as relativistic expressions for the Joule heating. Section 3 describes the microphysics used in our models. In section 4 we present the results of our model calculations and section 5 is devoted to the discussion of the results.
## 2 General Relativistic formalism of static spherical stars
The spacetime geometry will be assumed to be spherically symmetric, which means that our results may not be accurate for fast rotating neutron stars. Employing the familiar Schwarzschild coordinates $`(t,r,\theta ,\varphi )`$, the interior and exterior spacetime geometry takes the form (Misner et al. 1973MTW73 ; Wald 1984W84 )
$$ds^2=e^{2\mathrm{\Phi }}c^2dt^2+\frac{dr^2}{12Gm/c^2r}+r^2d\mathrm{\Omega }^2.$$
(1)
The radial proper length is thus $`dl=dr/\sqrt{12Gm/c^2r}`$ and the proper time $`d\tau =e^\mathrm{\Phi }dt`$.
Einstein’s equations coupled to a perfect fluid energy-momentum tensor give us the standard equations (Misner et al. 1973MTW73 ; Wald 1984W84 ):
$$\frac{dm}{dr}=4\pi r^2\rho ,$$
(2)
which determines the so called mass function $`m=m(r)`$,
$$\frac{d\mathrm{\Phi }}{dr}=\frac{Gmc^2+4\pi Gr^3P}{c^4r^2(12Gm/c^2r)},$$
(3)
for the ‘gravitational potential’ $`\mathrm{\Phi }=\mathrm{\Phi }(r)`$ and the TOV equation of hydrostatic equilibrium
$$\frac{dP}{dr}=(\rho c^2+P)\frac{d\mathrm{\Phi }}{dr}=$$
$$\frac{(\rho +P/c^2)(Gm+4\pi Gr^3P/c^2)}{r^2(12Gm/c^2r)}.$$
(4)
Regularity of the geometry at $`r=0`$ implies the inner boundary condition for Eq. 2
$$m(r=0)=0$$
(5)
while for Eq. 4 the central pressure
$$P(r=0)=P_c$$
(6)
is specified, through the EOS of the form $`P=P(\rho )`$, with the central density $`\rho _c`$ treated as a free parameter. Due to the linear nature of Eq.(3), $`\mathrm{\Phi }`$ can be scaled so that it always can be arranged to fulfill:
$$e^{\mathrm{\Phi }(R)}=\sqrt{1\frac{2GM}{c^2R}},$$
(7)
Once the interior spacetime geometry has been so specified it is joined smoothly across the ”surface” of the star to an exterior Schwarzschild field characterized by $`M=m(R)`$.
It should be mentioned that the stellar surface in our computation is fixed by
$$R=R_{\mathrm{star}}=r(\rho =\rho _b)$$
(8)
where $`\rho _b=10^{10}`$ g cm<sup>-3</sup>. This guarantees that the EOS is temperature independent. The layers at densities below $`\rho _b`$, called the envelope, are treated separately (see § 2.4).
### 2.1 Thermal evolution equations
Besides the above equations of stellar structure we shall need the equations describing the thermal evolution of the star. At the temperatures we are interested in, the neutrinos have a mean free path much larger than the radius of the star (Shapiro & Teukolsky 1983ST83 ) and thus leave the star once they are produced. Energy balance arguments (see for instance Thorne (1966T66 ) then imply
$$\frac{d(Le^{2\mathrm{\Phi }})}{dr}=\frac{4\pi r^2ne^\mathrm{\Phi }}{\sqrt{12Gm/c^2r}}\left(\frac{dϵ}{dt}+e^\mathrm{\Phi }(q_\nu q_h)\right)$$
(9)
where $`L`$ is the internal luminosity, $`ϵ`$ the internal energy per baryon, $`q_\nu ,q_h`$ is the neutrino emissivity and heating rate per baryon while $`n`$ stands for the baryon number density. The corresponding inner boundary condition for $`L`$ is
$$L(r=0)=0$$
(10)
The time derivative of $`ϵ`$ can be written in the form
$$\frac{dϵ}{dt}=\frac{dϵ}{dT}\frac{dT}{dt}=c_v\frac{dT}{dt}$$
(11)
through the specific heat at constant volume $`c_v`$ (which, for degenerate matter, is the same as the specific heat at constant pressure $`c_P`$).
The energy transport equation is :
$$\frac{d(Te^\mathrm{\Phi })}{dr}=\frac{3}{16\sigma _{SB}}\frac{\kappa \rho }{T^3}\frac{Le^\mathrm{\Phi }}{4\pi r^2\sqrt{12Gm/c^2r}}$$
(12)
where $`T`$ is the local temperature, $`\sigma _{SB}`$ the Stefan-Boltzmann constant and $`\kappa `$ the ‘opacity’. (Notice that within the relativistic framework an ‘isothermal’ configuration is defined by $`e^\varphi T=\mathrm{constant}`$ instead of $`T`$ = constant.) The associated boundary condition is
$$T_b=T_b(L_b)$$
(13)
which relates the temperature at the outer boundary (defined more precisely further bellow,), $`T_b`$, to the luminosity, $`L_b`$, in this layer. The location of this outer boundary layer is chosen such that $`L_b`$ be equal to the total photon luminosity of the star, $`L_{}L(r=R)`$, which in turn is related to the effective temperature $`T_e`$ by $`L_{}4\pi R^2\sigma _{SB}T_e^4`$. We can thus write Eq. 13 as $`T_b=T_b(T_e)`$ and this ‘$`T_b`$ \- $`T_e`$ relationship’ is discussed in Sect. 2.4.
The opacity is related to the total (thermal) conductivity by:
$$\lambda =\frac{16\sigma _{SB}T^3}{3\kappa \rho }.$$
(14)
If one neglects the red-shift $`e^\mathrm{\Phi }`$ and defines the energy flux $`F`$ as $`L/(4\pi r^2)`$ one can write Eq. 12 as
$$F=\lambda T$$
(15)
where $``$ is the radial gradient calculated with the proper length, which is the usual form of the heat conduction equation. The total conductivity is the sum of the electron and photon conductivities
$$\lambda =\lambda _e+\lambda _\gamma $$
(16)
since these two processes of heat conduction work independently and in parallel.
We will present our results of thermal evolution by using the ‘effective temperature at infinity’ $`T_e^{\mathrm{}}T_ee^{\mathrm{\Phi }(R)}`$ related to the ‘luminosity at infinity’ $`L_{}^{\mathrm{}}L_{}e^{2\mathrm{\Phi }(R)}`$ through the ‘radius at infinity’ $`R^{\mathrm{}}Re^{\mathrm{\Phi }(R)}`$ by
$$L_{}^{\mathrm{}}=4\pi (R^{\mathrm{}})^2\sigma _{SB}(T_e^{\mathrm{}})^4.$$
(17)
These three quantities ‘at infinity’ are, in principle, measurable and, in particular, $`R^{\mathrm{}}`$ would be the areal radius of the star that an observer ‘at infinity’ would measure with an extremely high angular resolution instrument (Page 1995P95 ).
The solution of the complete set of equations of stellar structure, Eqs. (2-7) and thermal evolution, Eqs. (9-13), requires knowledge of the equation of state, the opacity $`\kappa `$, the neutrino emissivity $`q_\nu `$ and also the specific heat $`c_v`$. We shall devote Sect. 3 to the detailed specification of those variables.
### 2.2 Magnetic evolution equations
Besides the equations of stellar structure and thermal evolution we shall also need the equations describing the magnetic field evolution. In this paper consideration will be restricted to dipolar (poloidal) magnetic fields. The GR formulation of the evolution equation of such fields has been discussed in detail by GPZ00GPZ00 while for a more general set up see Rädler et al. (2000RFGZ00 ). As has been shown in the first reference, such a field can be expressed in terms of a relativistic generalization of the familiar Stoke’s stream function by:
$$B^r(t,r,\theta )=\frac{2F(t,r)}{r^2}\mathrm{cos}\theta ,$$
(18)
$$B^\theta (t,r,\theta )=\frac{1}{r}\left(1\frac{2Gm}{c^2r}\right)^{\frac{1}{2}}\frac{F(t,r)}{r}\mathrm{sin}\theta .$$
(19)
while the relevant induction equation (see Appendix) yields the following equation for the relativistic Stoke’s function
$$\frac{4\pi \sigma }{c^2}e^\mathrm{\Phi }\frac{F}{t}=\left(1\frac{2Gm}{c^2r}\right)\frac{^2F}{r^2}$$
$$+\frac{1}{r^2}\frac{F}{r}\left[\frac{2Gm}{c^2}+\frac{4\pi G}{c^2}r^3\left(\frac{P}{c^2}\rho \right)\right]\frac{2}{r^2}F.$$
(20)
The appropriate boundary conditions, as $`r0`$, is the same as in the flat space case: a regular field at the star’s center requires $`\frac{F(t,r)}{r^2}`$ to be finite. The outer boundary condition, however, differs from that valid in the flat space, and its GR form is as follows (see: GPZ00GPZ00 ):
$$R\frac{F(t,r)}{r}|_R=G(y)F(t,R)$$
(21)
where:
$$G(y)=y\frac{2y\mathrm{ln}(1y^1)+\frac{2y1}{y1}}{y^2\mathrm{ln}(1y^1)+y+\frac{1}{2}}\mathrm{with}y=R/R_S$$
(22)
($`R_S2GM/c^2`$ being the star’s Schwarzschild radius). Since $`G(y)<0`$ (in particular, in the flat space-time case, $`G(\mathrm{})=1`$) the boundary condition forces a bending of $`F`$ in the upper layers.
As an initial profile for the Stoke function we use the same formula as UK97, for later comparison with the work of these authors:
$$F(r,0)=B_0R^2(1r^2/r_0^2)/(1R^2/r_0^2)\mathrm{at}r>r_0$$
(23)
$$F(r,0)=0\mathrm{at}r<r_0$$
Notice that this initial $`F(r,0)`$ does not satisfy the outer boundary condition, but will immediately be forced to do it at the first numerical time step. There will thus be a rapid relaxation of $`F`$ in its early evolution due to the enforcement of the boundary condition and the propagation of the resulting curvature of $`F`$ toward higher densities.
### 2.3 Joule heating
Due to the finite conductivity $`\sigma `$, magnetic energy is dissipated into heat (= Joule heating). The heat production per unit (proper) time and unit (proper) volume is given by
$$Q_h=nq_h=\frac{j^2}{\sigma }$$
(24)
where $`j`$ is the current. As it is shown in the Appendix, for a dipole poloidal magnetic field $`B`$, $`j`$ is given by $`j=j_\varphi e_\varphi `$ with
$$j_\varphi =\frac{c}{4\pi }\frac{\mathrm{sin}\theta }{r}\times $$
$$\{e^\mathrm{\Phi }\left(1\frac{2Gm}{c^2r}\right)^{\frac{1}{2}}\frac{}{r}\left[e^\mathrm{\Phi }\left(1\frac{2Gm}{c^2r}\right)^{\frac{1}{2}}\frac{F}{r}\right]\frac{2F}{r^2}\}.$$
(25)
It is seen immediately that in the limit $`e^\mathrm{\Phi }=1`$ and $`m=0`$ the above expression reduces to its flat space-time counterpart (Miralles et al. 1998MUK98 ). Since our numerical calculations assume spherical symmetry we use, in Eq. 25, a spherical average of $`\mathrm{sin}^2\theta `$ ($`<\mathrm{sin}^2\theta >=\frac{2}{3}`$) obtaining
$$<Q_h>=\frac{c^2}{24\pi ^2}\frac{1}{\sigma }\frac{1}{r^2}\times $$
$$\{e^\mathrm{\Phi }\left(1\frac{2Gm}{c^2r}\right)^{\frac{1}{2}}\frac{}{r}\left[e^\mathrm{\Phi }\left(1\frac{2Gm}{c^2r}\right)^{\frac{1}{2}}\frac{F}{r}\right]\frac{2F}{r^2}\}^2.$$
(26)
### 2.4 Outer boundary: magnetized envelopes
The layers at densities below $`\rho _b=10^{10}`$ g cm<sup>-3</sup> are defined as the envelope, and extend up to the atmosphere, where the photosphere is located, while by interior we mean the whole star where $`\rho >\rho _b`$. The presence of the magnetic field affects strongly the heat transport in the envelope (but not in the deeper layer where $`\rho >10^{10}`$ g cm<sup>-3</sup>, which motivates the above choice of $`\rho _b`$), and results in a non uniform distribution of the surface temperature. The corresponding ‘$`T_bT_e`$ relationships’ have been calculated by Page & Sarmiento (1996PS96 ) for dipolar (and quadrupolar) fields. We use these results which hence give us a field dependent ‘$`T_bT_e`$ relationship’ that adjust itself to the evolution of the magnetic field. Notice however that, for fields much stronger than 10<sup>12</sup> G, this relationship is not reliable when $`T_e`$ is much lower than 10<sup>6</sup> K, and it is most probably very inaccurate when $`T_e`$ is below 10<sup>5</sup> K.
Our outer boundary condition assumes that the envelope is made of catalyzed matter. If an upper layer of light elements were present, the heat transport is strongly enhanced when no magnetic field is present (Potekhin, Chabrier, & Yakovlev 1997PCY97 ), but there is, to date, no published model of magnetized envelopes with light elements.
One should finally emphasize that when the thermal evolution is controlled by the Joule heating the star’s luminosity, and $`T_e`$, is given by the heating rate and is independent of the outer boundary condition as discussed at the end of § 4.2.
This closes the system of equations and boundary conditions to be solved.
### 2.5 Numerical method
The thermal evolution equations are solved by a Henyey-type code (e.g., Page 1989P89 ) while the induction equation for the Stoke function is solved with a Crank-Nicholson method (Press et al., 1986PFTV86 ). The whole set of equations for the thermo-magnetic evolution should be solved simultaneously at each time step but we have decided to solve the thermal and then the magnetic equations alternatively, i.e., the thermal equations are solved at a given time step using the field of the previous step (which appears in the Joule heating term) and once the new temperature is obtained the induction equation is solved to obtain the new magnetic field. This method is much faster than a full simultaneous solution and gives results which, as we have verified explicitly in a few case, are practically indistinguishable from the full simultaneous solution.
## 3 Input microphysics
### 3.1 The equation of state
The first ingredient needed to build NS models is the equation of state (EOS). In principle the EOS should give us not only the relationship between pressure and density, i.e., $`P=P(\rho )`$, but also the chemical composition of matter.
We separate the crust from the core at the density $`\rho =\rho _{cr}1.6\times 10^{14}`$ g cm<sup>-3</sup> (Lorenz, Ravenhall & Pethick 1993LRP93 ) and use the EOSs of Negele & Vautherin (1973NV73 ) for the inner crust, at $`\rho >\rho _{drip}4.4\times 10^{11}`$ g cm<sup>-3</sup>, and Haensel, Zdunik & Dobaczewski (1989HZD89 ) for the outer crust, i.e., we assume that the chemical composition is that of cold catalyzed matter. The EOS, and its associated chemical composition, in the crust is well determined under the assumption that matter is in its (catalyzed) ground state. There is however still the possibility that a strong phase of hypercritical accretion occurred after the supernova explosion, which may, or may not, alter the chemical composition of the crust. We will not consider this possibility here but only mention that it does have an enormous effect on the magnetic field and its subsequent evolution (Geppert, Page & Zannias 1999GPZ99 )
In the core, the EOS is relatively well constrained up to densities around 2 - 3 $`\times \rho _{cr}`$ while its behavior at higher densities is still a mystery (Prakash 1998Pr98 ). We will thus consider three different cases which hopefully illustrate the whole range of possibilities. We take as a ‘Medium EOS’ the one calculated by Wiringa, Fiks & Fabrocini (1988WFF88 ), using their model called av14+UVII, a ‘Stiff EOS’ from Pandharipande, Pines & Smith (1976PPS76 ) and a ‘Soft EOS’ from Pandharipande (1971P71 ). Notice that the Stiff EOS is based on the presence of a lattice of neutrons in the inner core and is not anymore considered as realistic but we still use it since it has been used by many authors and represents the case of extreme stiffness. Moreover, none of the consequences of the presence of such a lattice phase is taken into account in our calculations, e.g., in $`C_v`$, $`ϵ_\nu `$ and the transport coefficients: we will assume that neutrons and protons form a quantum liquid and also boldly consider them as superfluid and superconductor. Our opinion is that this EOS should be abandoned but we consider it for comparison with previous works of other authors. The Soft EOS on the other hand, despite of its age, is still representative of modern soft EOSs and includes hyperons.
We will consider models of 1.4 $`M_{}`$stars whose overall properties are listed in Table 1.
### 3.2 Neutrino processes
Neutrino emission drives the cooling as long as the internal temperature is higher than about $`10^8`$ K. For processes in the crust we consider the two dominant ones which are the plasmon process and the electron-ion bremsstrahlung (Page 1989P89 ). The former is only relevant during the first few years of the life of the NS but is very strong and brings down the crust temperature to about $`10^9`$ K, relaxing it from the arbitrary initial conditions. The latter has only a very small effect, mostly when the surface temperature is around $`10^6`$ K, i.e., for young stars.
The crucial neutrino emission processes occur in the core and we consider two scenarios (see Page 1998P98 for more details). In the ‘Standard Cooling’, or ‘slow cooling’, scenario we include the modified Urca processes and their associated, and weaker, bremsstrahlung processes following Yakovlev & Levenfish (1995YL95 ). This scenario applies when the NS core contains only neutrons, protons, as well as electrons and muons which maintain charge neutrality, and the proton fraction is low enough that the direct Urca process is forbidden by momentum conservation. In the ‘Enhanced Cooling’, or ‘fast cooling’, cases we add a strong neutrino emission at densities larger to $`\rho _{fast}=4\times 10^{14}`$ g cm<sup>-3</sup> with a rate
$$ϵ_\nu ^{FAST}=10^{26}(\rho /\rho _{cr})^{2/3}T_9^6\mathrm{erg}\mathrm{g}^1\mathrm{s}^1$$
(27)
where $`T_9T/10^9`$ K. This rate is representative of many of the possible enhanced ones as the direct Urca from nucleons or hyperons, but is stronger than what produced by a pion or kaon condensate (Prakash 1998Pr98 ). For, comparison the inefficient modified Urca process gives approximately $`ϵ_\nu ^{MURCA}10^{21}T_9^8\mathrm{erg}\mathrm{g}^1\mathrm{s}^1`$.
In our three model 1.4 $`M_{}`$stars, the masses of the inner cores where the fast neutrino emission is allowed for the ‘Enhanced Cooling’ cases are, 0.11, 1.29 and 1.39 $`M_{}`$for the Soft, Medium and Stiff EOS, respectively.
### 3.3 Electrical conductivities
In general, the electrical conductivity $`\sigma `$ can be expressed in terms of the electron relaxation time $`\tau `$ as
$$\sigma =\frac{e^2n_e\tau }{m_e^{}},$$
(28)
where
$$m_e^{}=\mu _e/c^2=m_e[1+1.018(\rho _6Z/A)^{2/3}]^{1/2}$$
(29)
is the electron effective mass and $`n_e`$ the electron number density ($`\mu _e`$ being the electron chemical potential, $`\rho _6=\rho /10^6`$ g cm<sup>-3</sup>, Z and A the charge and mass number of the ions). In the liquid phase we use the calculation of $`\tau =\tau _{\mathrm{e}\mathrm{i}}`$ for electron-ion scattering by Itoh et al (1983IMII83 ). In the solid phase $`\tau `$ is given by
$$\frac{1}{\tau }=\frac{1}{\tau _{\mathrm{e}\mathrm{ph}}}+\frac{1}{\tau _{\mathrm{e}\mathrm{imp}}}$$
(30)
where we use the results of Itoh et al (1984IKMS84 ) for the electron-phonon scattering $`\tau _{\mathrm{e}\mathrm{ph}}`$ and of Yakovlev & Urpin (1980YU80 ) for the electron-impurity scattering $`\tau _{\mathrm{e}\mathrm{imp}}`$. In general, $`\tau _{\mathrm{e}\mathrm{i}}`$ is a function of $`\rho `$, A and Z, and $`\tau _{\mathrm{e}\mathrm{ph}}`$ depends also on the temperature $`T`$ while $`\tau _{\mathrm{e}\mathrm{imp}}`$ depends on $`\rho `$, $`Z`$ and the impurity concentration $`Q_{\mathrm{imp}}`$. We show in Fig. 1 the value of $`\sigma `$ for typical values of $`T`$ and $`Q_{\mathrm{imp}}`$.
### 3.4 Pairing
The occurrence of pairing, neutron superfluidity and proton superconductivity, strongly affects both the thermal and magnetic evolution of neutron stars. The thermal effects are very strong during the neutrino cooling phase, which last about 10<sup>5</sup> to almost 10<sup>7</sup> yrs depending on the model, and the subsequent photon cooling phase (see, e.g., Page 1998P98 for a review). As a result, pairing will in a large part determine the time at which Joule heating starts to control the thermal evolution and then during this Joule heating phase the effect of pairing becomes negligible. We treat the suppressive effect of pairing on $`C_v`$ and $`ϵ_\nu `$ according to the treatment of Levenfish & Yakovlev (1994aLY94a , 1994bLY94b ) and Yakovlev & Levenfish (1995YL95 ).
The superconductive phase in the core has the effect of producing an impenetrable barrier for the magnetic field, which is initially confined to the crust in the models of the present work. This guarantees the confinement of the magnetic field, and the currents, to the crust since the superconductive phase transition happens, at the crust-core boundary, well before the magnetic field had time to diffuse to this layer. For definiteness we plot in Fig. 2 the pairing critical temperatures that we adopt in this paper. Notice that the values we adopt for <sup>1</sup>S<sub>0</sub> pairing of both neutrons and protons are typical of modern calculations. In the case of neutron <sup>3</sup>P<sub>2</sub> pairing we explicitly adopt values which vanish a densities above 10<sup>15</sup> g cm<sup>-3</sup> to ensure that pairing will not suppress the strong neutrino emission in our Enhanced Cooling scenarios for the Soft and Medium EOSs.
In the case of the Stiff EOS any published calculation of neutron or proton pairing shows a non vanishing value of $`T_c`$ at the density in the center of the star, given the low value of this central density. To avoid suppression of the neutrino emission by pairing we will assume that the fast neutrino emission is not affected by neutron and proton pairing in this case of Stiff EOS in the Enhanced Cooling scenario .
## 4 Results
We will present here our results for the thermal, magnetic and rotational evolution of isolated neutron stars considering 1.4 $`M_{}`$stars built with the three EOSs described above within both the ‘standard’ and the ‘fast’ cooling scenarios. In order to investigate the influence of the (a priori unknown) initial structure and strength of the magnetic field onto NS’s evolution, we considered for each EOS and cooling scenario three classes of qualitatively different field models, existent at the beginning of the evolution. They are characterized by the initial surface field strength $`B_0`$, depth of penetration of the current $`\rho _0`$ and impurity concentration $`Q`$ as listed in Table 2.
### 4.1 A detailed example
We first discuss here an example of the internal evolution of the Stoke function, magnetic field, currents and local heating rate, as shown in Fig. 3, which will help for the general discussion presented below. We choose the 1.4 $`M_{}`$star with the medium EOS and ‘standard’ cooling; the initial surface field strength is 10<sup>12</sup> G and the currents are initially located at $`\rho _0=10^{13}`$ g cm<sup>-3</sup>. The impurity content for $`\sigma _{\mathrm{imp}}`$ is $`Q_{\mathrm{imp}}=0.01`$. This is the model 2 of the central panel of Fig. 5 but it is reproduced in Fig. 4 along with two similar models with different impurity contents (upper panel) and the evolution of the internal temperature (lower panel). GR effects are included but will be discussed in the next subsections.
The four chosen times, labeled as a, b, c & d are marked in the upper panel of Fig. 4 and correspond, respectively, to a) the initial field decay during the period when the neutron star is still hot and $`\sigma `$ is temperature dependent, being dominated by electron-phonon scattering, b) the plateau where the cooling lead to an enormous increase of $`\sigma `$, and consequently a stagnation of the field decay, which will eventually be controlled by impurity scattering, c) second phase of decay of the field when time becomes comparable to the impurity scattering decay time scale and, finally, d) the late exponential decay.
The four panels of Fig. 3 directly illustrate the diffusion of the field toward regions of higher conductivity as time runs. Once $`\sigma `$ becomes $`T`$ independent the diffusion equation formulates an eigenvalue problem whose solution can be formally written as $`F(\rho ,t)=_{n=1}^{\mathrm{}}\mathrm{exp}(t/\tau _n)a_nF_n(\rho )`$, in terms of eigenmodes $`F_n`$ with decay times $`\tau _n`$, and expansion coefficients $`a_n`$. Since the n-th mode, $`F_n(\rho )`$, has n nodes, in the crust, initially a very large number of modes must contributes significantly, i.e., with large positive and negative coefficients $`a_n`$, to produce a mutual cancellation resulting in a vanishing $`F(\rho ,t)`$ in the high density region. The diffusion of the stoke function into the high density region is simply due to the faster decay of the modes with $`n>1`$ compared to the nodeless fundamental mode. When the fundamental mode is dominating, i.e., when the stoke function is non zero in the whole crust (it has reached the crust-core boundary through diffusion), the field evolution is a power-law like decay, phase c (Urpin, Chanmugan & Sang 1994UCS94 ). Finally, when only the fundamental mode $`n=1`$ is left the decay becomes purely exponential, phase d.
Notice that the field strength inside the crust is at least one order of magnitude higher than at the surface: this is due to a very large $`B_\theta `$ forced by the presence of the derivative $`dF/dr`$ in Eq. 19, while at the surface the boundary condition ensures that $`B_r`$ and $`B_\theta `$ are comparable.
With respect to the currents, noticeable are the negative currents in the layers where the field is growing which are due to induction, i.e., Lenz law. Once the field has reached the crust-core boundary the currents are positive in the whole crust and the only negative currents left are the supercurrents induced in the skin layer of the proton superconductor. Instead of imposing an ad-hoc boundary condition at the crust-core interface to simulate the effect of the proton superconductor we have preferred to keep the central boundary condition and introduce an enormous value for $`\sigma `$, 10<sup>200</sup> s<sup>-1</sup>, once protons become superconductor, i.e., when $`T<T_c`$. This allows us to see explicitly the induced supercurrents which, however, because of the finite radial resolution of the numerical scheme, are located in the whole last zone of the crust instead of the physical skin layer of the superconductor which is a few tens of fermis thick. We have checked explicitly that an ad-hoc crust-core boundary condition for $`F`$ gives the same results as our boundary condition with enormous $`\sigma `$.
As a last remark, from Fig. 4, and comparing with Fig. 1, we see that impurity scattering starts to dominate $`\sigma `$ after the field has reached the plateau (phase b): this stagnation value of the field in independent of $`Q_{\mathrm{imp}}`$ and is only a result of the enormous increase of $`\sigma `$ due to the cooling. After this, the length of the plateau (phase b) is controlled by $`Q_{\mathrm{imp}}`$, higher impurity contents leading naturally to an earlier onset of the second phase of field decay (phase c).
### 4.2 Magneto-thermal evolution within the ‘standard’ cooling scenario
In Figs. 5 and 6 we present the modeled evolution of the, more or less, observable quantities, i.e. the surface temperature (measured by a far distant observer) and the surface magnetic field. While the surface temperatures, or, at least, its upper limits, of isolated NSs can be inferred from X–ray spectra, the magnetic field of isolated pulsars is mostly estimated by the precise measurement of the rotational period and its time derivative.
Within the ‘standard’ cooling scenario the surface temperature stays quite high ($`10^6`$ K) during the neutrino dominated cooling era, which lasts about $`10^6`$ yrs (for our stiff EOS models) to almost $`10^7`$ yrs (soft EOS models). Later on the cooling is driven by photon emission from NS’s surface which appears in the cooling curves by a strong increase in the slope. During the neutrino cooling era, the surface temperature drops down by a factor of about $`2`$ (stiff EOS) to $`4`$ (soft EOS), i.e. the approximately isothermal crust has at the end of the neutrino cooling epoch a temperature of about $`510^7`$ (soft EOS) to $`10^8`$K (stiff EOS). As seen from Fig. 1, for such and higher temperatures, the electrical conductivity in the crust is determined by electron–phonon collisions except at the highest densities. With the following temperature drop during the photon cooling era, the electron-phonon relaxation time $`\tau _{\mathrm{e}\mathrm{ph}}`$ increase dramatically and thus impurity scattering will begin to control $`\sigma `$ (see Eq. 30), at a temperature, and thus an age, depending on the impurity concentration. From that stage on, $`\sigma `$ becomes temperature independent.
Notice that the differences in the thermal evolution for the different EOSs, before Joule heating becomes efficient, are mostly due to the differences in the fraction of the core which is paired. A larger paired region implies a lower neutrino emission, and thus a higher temperature during the neutrino cooling era, and also a lower specific heat, and thus an earlier transition to the photon cooling era and a faster temperature drop during that era (see, e.g., Page 1998P98 ). We have taken here the choice of well defined density dependences of $`T_c`$, for both the neutrons and protons, independently of the EOS, as shown in Fig. 2. Different density dependences of the $`T_c`$’s would obviously give different results. For example, assuming high values of $`T_c`$ down to center of the star for the soft EOS would results in a cooling history practically indistinguishable form the cooling history of the stiff EOS model. Given the present uncertainty on the value and density dependence of $`T_c`$ for <sup>3</sup>P<sub>2</sub> neutron pairing (Baldo et al. 1998BEEHS98 ) any choice has, unfortunately, some arbitrariness. Any effect of the resulting cooling history on the field evolution should thus be specifically formulated in terms of cooling history and not in terms of the stiffness of the EOS.
We have, very roughly, for the field decay time-scale, $`\tau _{\mathrm{decay}}l^2`$, $`l`$ being a typical length scale of the crustal field structure. This immediately implies that $`\tau _{\mathrm{decay}}`$ increases when $`\rho _0`$ increases and also when the stiffness of the EOS is increased, since both increase $`l`$. Moreover, since also $`\tau _{\mathrm{decay}}\sigma `$, a higher $`\rho _0`$ locates the currents in a region of higher $`\sigma `$ and increases more $`\tau _{\mathrm{decay}}`$. Thus, the models 1 give fast decay, models 2 intermediate decay and models 3 slow decay, as is clear form Fig. 5.
The cooling influences the field decay by rising $`\sigma `$ till it becomes temperature independent when dominated by impurity scattering. This happens during the photon cooling era as mentioned above, and happens earlier for stiffer EOSs given our choice of EOS independent $`T_c`$’s. Consequently, we obtain that a stiffer EOS results in a slower field decay because of its cooling behavior and also because of the larger length scale $`l`$. Our choice of $`T_c`$ thus maximalize the effect of the EOS’s stiffness on the field decay.
We now turn our discussion to the analysis of GR effects upon magnetic field evolution. For each model we simulated the field evolution with and without GR effects, but note that GR effects on the star structure and cooling have always included. (Had we for instance, have turned off the GR effects on the star structure, the resulting models would make no sense at all since the differences in the size, the central density etc, of the model would tender them unrealistic.) On the other hand, GR effects on on the cooling have already been discussed in the literature long ago (Nomoto & Tsuruta 1987NT87 , Gudmundsson et al. 1983GPE83 ). The models with GR effects included are marked as ‘GR’ in Fig. 5 and drawn with thick lines. It is seen from those plots that the decay of the field is faster in NSs built on a soft EOS than in the case of a medium or stiff EOS but the decelerating GR effects are more pronounced. This is most remarkable for long living fields: the difference of the surface field strength for model 3, with the soft EOS, after $`10^{10}`$yrs of evolution is about two orders of magnitude. Comparing the field evolution for the soft and stiff EOS cases, while the final surface field strength for model 3 in the stiff case is larger than those of the soft case by a factor of 300 when GR effects are neglected, that factor reduces to about 7 when relativistic effects are taken into account correctly.
In the late photon cooling era, when most of the initial thermal energy of the NS has been irradiated away, the Joule heating by the decay of the crustal field completely determines the cooling behavior. Notice, however, that most of the magnetic energy has been dissipated earlier when it had no effect on the cooling. The amount of heat released in the process of Joule heating is determined by the strength of the field at that time and its decay rate. Therefore, an initial field configuration as given by models of class 3 will result in a significant Joule heating while in models of class 1 the effect is practically nil. The field decay of the class 2 models yields some noticeable Joule heating only in case of the stiff EOS since its strength is large enough for that until about $`10^8`$yrs. For long periods Joule heating is especially effective in NSs with a stiff EOS: their crust is quite thick, they keep a larger field and there is a lot of magnetic energy to dissipate. The effect of GR on the Joule heating is not spectacular but still significant in the case of a soft EOS at the latest stages. For the different EOSs it can result in stronger or weaker heating: the rate of dissipation of magnetic energy is lower with GR but the magnetic energy to be dissipated is larger due to the previous slower evolution and the net effect can be an enhancement or a reduction of the heating.
Notice, finally, that at these late ages when the thermal evolution is entirely controlled by the Joule heating the outer boundary condition, Eq. 13, has no influence on the thermal evolution: this is due to the fact that the heating mechanism is temperature independent since $`\sigma `$ is totally controlled by impurity scattering. The star’s luminosity $`L_{}`$ is thus simply given by the total heating rate integrated over the whole star
$$L_{}^{\mathrm{}}=Hq_he^{2\mathrm{\Phi }}𝑑V$$
(31)
since, given the low temperature at this times, $`q_\nu q_h`$. This is very fortunate for our study since the ‘$`T_bT_e`$ relationship’ is poorly known at the low temperatures reached in this late phase of evolution. It would of course not be the case if the heating mechanism were temperature dependent.
Notice finally that, in the case of magnetars, the Joule heating dominates the cooling from the very beginning and thus it has a strong effect onto the field evolution (Geppert, Page, Colpi & Zannias 1999GPCZ99 ).
### 4.3 Magneto-thermal evolution within the ‘fast’ cooling scenario
Enhanced neutrino emission is caused by the presence of exotic states of matter in the NS core. Kaon or pion condensation, hyperons, quarks as well as possible direct Urca processes enlarge the neutrino emissivity considerably. This results in a very fast cooling of the core, so that during the so called isothermalization phase the heat of the hotter crust is transported into the core. Depending on the EOS and the assumptions about pairing that phase can last of the order of 1 to 100 yrs (see, e.g., Page 1998P98 ). During this phase the surface temperature drops very rapidly and after this the crustal temperature is so low that the conductivity is almost completely controlled by electron–impurity collisions. Thus, for a period of time given by the impurity decay time scale $`\tau _{\mathrm{imp}}`$, the field suffers practically no decay. For $`t>\tau _{\mathrm{imp}}`$ the field decays according to a power law and, when it has diffused down to the crust–core boundary, the decay becomes exponential. In class 3 models the power law like decay is missed because the field is already initially located close to the crust–core boundary so that it reaches this depth during $`\tau _{\mathrm{imp}}`$. Notice that $`\tau _{\mathrm{imp}}Q^1`$ and that it is also reduced with the softening of the EOS as an effect of the decreasing crustal thickness. Generally, the field decay is slower than in the ‘standard’ cooling scenario.
Due to the accelerated cooling, the Joule heating, even for the model 1 and 2 fields, dominates the cooling earlier compared to the ‘standard’ cooling scenario. At ages above $`10^7`$$`10^8`$ yrs, however, the star’s temperature is comparable, or even higher, to that predicted by the analogous models within the ‘standard’ cooling scenario. In class 2 and particularly class 3 models with the Stiff EOS, the stars’ temperatures in this range of ages are noticeably higher, compared to the analogous ‘standard’ cooling cases, since the magnetic field is higher and thus the joule heating more efficient. After $`10^{10}`$yrs, however, the difference in the surface temperatures between the ‘standard’ and the accelerated cooling scenario vanishes.
### 4.4 Rotational evolution
Given the temporal evolution of the magnetic field for the various EOSs and cooling scenarios, the rotational period of the NSs can be estimated by integrating the equation which relates the loss of rotational energy to the radiation of electromagnetic energy by magneto–dipole radiation
$$P\dot{P}=\frac{2\pi ^2}{3}\frac{R^6\stackrel{~}{B}^2(t)}{cI},$$
(32)
where $`\stackrel{~}{B}`$ is surface magnetic field at the magnetic pole, and $`I`$ is the star’s moment of inertia. A very accurately relativistic expression for the moment of inertia, is given by $`I=0.21e^{2\mathrm{\Phi }}MR^2`$ (Ravenhall & Pethick 1994RP94 ). However, this classical magneto–dipole radiation formula assumes flat space-time. Ideally we would like to have the exact GR version of this equation. To our knowledge, the GR version of it remains to be calculated and in the absence of the exact formula we shall only use an approximate expression. We may recall that Eq. 32 implicitly considers the energy loss at the light cylinder and then extrapolates the field strength back to the star’s surface assuming the flat space-time $`1/r^3`$ radial dependence of the field. The full GR solution of a slowly rotating dipolar field, in vacuum, shows however (Muslimov & Tsygan 1990MT90 , 1992MT92 ) that the actual field at the star’s magnetic pole is amplified by a factor $`f=3y^2[y\mathrm{ln}(1y^1)+\frac{1}{2}(2+y^1)]`$, where $`yR/R_S`$, compared to the flat space-time one. For instance, for our soft EOS star of 1.4 $`M_{}`$, the resulting value of the amplification factor is $`f=1.86`$ while for the stiff one, $`f=1.27`$ (Page & Sarmiento 1996PS96 ). In view of the Muslimov-Tsygan results, we will henceforth relate the surface magnetic field $`B`$ at magnetic pole as we numerically calculated, it to the corresponding field $`\stackrel{~}{B}`$ used in Eq. 32 for the rotational evolution by $`B=f\stackrel{~}{B}`$. For each of our three EOS (and a fixed stellar mass of 1.4 $`M_{}`$) and the two cooling scenarios, we specify the initial field strength and current depths $`\rho _0`$ and finally the impurity content $`Q_{\mathrm{imp}}`$. The specification of those parameters yields a unique evolution for the surface $`B`$-field and via Eq. 32 the corresponding $`P`$. The resulting magneto-rotational evolution of isolated NSs is presented in Fig. 7.
We have also plotted, as dots, the PSRs data of the Princeton catalogue (Taylor, Manchester, & Lyne 1993TML93 ): their $`B`$ field is calculated with Eq. 32, using $`P`$ and $`\dot{P}`$ from that catalogue and the appropriate GR ‘$`f`$ factor’ for each stellar model. The star’s radius $`R`$ and moment of inertia $`I`$ are correspond to values of the corresponding models having mass of 1.4 $`M_{}`$. The set of $`B`$$`P`$ diagrams of Fig. 7 is the arena where our theoretical models confront the reality. There is a major effect on the data interpretation arising from varying the EOS: we see that given the observed $`P`$$`\dot{P}`$ data the “observed” surface $`B`$-field of pulsars differ almost one order of magnitude when we consider models where the equation of state is varied from the extreme soft one up to the extreme stiff one (see also UK97UK97 ). This effect is a straightforward consequence of the fast growth of the magnetic field at small distance from the star which is moreover amplified by GR effects, combined with the much smaller stellar radius obtained by employing the soft EOS . It should emphasized however, that although a complete and accurate model of pulsar spin-down will probably alter the simple Eq. 32 (expected at least to modify the overall numerical coefficient) it certainly will not change the radial dependence of the field in the near zone.
A basic feature emerging from the Fig.(7), is the following: a given model is viable if it manages to maintain a strong enough field for a long enough time such that its rotational period can increase up to values compatible with the bulk of the observed pulsar. It is also clear that if a ‘standard’ cooling scenario applies and irrespectively of the EOS, the field decays much faster than in the case of the accelerated cooling: this results in smaller saturation values of the rotational periods compared to the accelerated cooling model, with the same initial field structure and the same EOS.
The evolutionary tracks offer the possibility to check the whether a given EOS can be naturally compatible with the observational constraints without extreme “fine tuning” of the initial parameters $`\rho _0`$ and $`B_0`$, of the impurity parameter $`Q_{\mathrm{imp}}`$ and its preference to standard versus enhanced cooling. Thus for instance: assuming a soft EOS, a comparison of the evolutionary tracks with the observed pulsar population strongly favors the enhanced cooling scenario. Within the ‘standard’ cooling scenario, only the strongest initial fields $`10^{14}`$ G are acceptable. In this case, even $`B_010^{13}`$ G requires very high $`\rho _0`$ and very low $`Q_{\mathrm{imp}}`$, definitely making the ‘standard’ cooling scenario with a soft EOS very unappealing. Notice that the fast cooling scenario is very natural for a soft EOS since the central densities reached in NSs are above ten times nuclear matter density. Even if the pulsar fields (for a given $`P`$) were much lower than estimated we would still be in the same situation since the high field is required to be able to spin down the pulsars to the range of observed $`P`$’s before the field decays too much. In the case of a medium EOS the requirements are much less stringent. Within the ‘standard’ cooling scenario, initial fields of the order of 10<sup>12</sup> G require high initial depths $`\rho _0`$ and low pollution in order for the evolutionary tracks to reach the bulk of the pulsar population. Finally, in the case of a stiff EOS, the field decay is slow enough that the whole pulsar population is reachable with initial fields in the range 10<sup>11</sup> \- 10<sup>13</sup> G without basically any significant restrictions on $`\rho _0`$ and $`Q_{\mathrm{imp}}`$ in both the fast and the slow cooling scenarios.
One cannot overemphasize that the above analysis is plagued by the intrinsic uncertainties of Eq. 32 and should be considered as only indicative. Moreover, it is worth stressing here that tracks starting with an initial field $`B_0=10^{14}`$G have to be considered with reserve since for such (and larger) field strengths the decay may not any longer described by a linear diffusion equation. The possible occurrence of a Hall–cascade (Goldreich & Reissenegger 1992GR92 ) would imply that the field evolution would deviate significantly from Eq. 20, which is strictly valid in the limit that the magnetization parameter $`\omega _B\tau <1`$ ($`\omega _B`$ being the Larmor frequency and $`\tau `$ the relaxation time of the electrons as given by Eq. 30).
## 5 Discussion and conclusions
We studied in detail, and for a large variety of possible models, the magnetic, thermal and rotational evolution of isolated NSs, assuming that their magnetic fields, and the currents supporting them, are confined to the stellar crust. Our calculations take into account, for the first time, all mutual effects of the thermal and magnetic evolution self consistently in a wholly GR formalism.
### 5.1 Comparison with the work of Urpin and Konenkov
Urpin & KonenkovUK97 (1997, UK97 thereafter) have presented the most detailed study of the evolution of crustal neutron star magnetic field while Miralles, Urpin & Konenkov (1998MUK98 ) completed the former work by the inclusion of Joule heating in the thermal evolution. It should be mentioned, however, that neither of these works incorporated GR effects on the $`B`$-field evolution.
Our field evolution results, without GR effects incorporated, are very close to those of UK97 (their curves labeled 3, dashed lines, in their figure 2 correspond to our models 2). The main differences are seen in the early evolution of the field (our phase a): this can be easily attributed to the fact that UK97 started the field evolution at an age of 100 yrs while we started it immediately at the star’s birth. Since UK97 used isothermal stars for their modeling of the B-field evolution they could not model the early hot, and non isothermal, phase properly. In our calculations, this time difference allows the initial currents to rapidly relax from their initial distribution (Eq. 23, which does not fulfill the outer condition) during a phase of low conductivity while the models of UK97 must do it in conditions of higher conductivity, i.e., their relaxation is much slower. However, if one contemplates the scenario in which crustal magnetic fields are generated by a thermomagnetic instability during this early hot phase (Blandford et al. 1983BAH83 , Urpin, Levshakov & Yakovlev 1986ULY86 , Wiebicke & Geppert 1996WG96 ) it is equally reasonable to start the B-field evolution by pure ohmic decay at the end of this phase as UK97 did. The slight differences seen in the strength of the field at the onset of the impurity scattering dominated phase (phase b) can be easily explained as due to slight differences in the cooling histories, particularly differences in the time at which photon cooling takes over neutrino cooling (the knee at ages around 10<sup>6</sup> – 10<sup>7</sup> yrs) and also as a result of the differences in the previous phase a.
Finally, the field evolution in the late phases, c and d, agree well with the results of UK97 once we take into account the small differences in the previous phases.
### 5.2 General Relativistic effects
While previous studies considered the influence of GR effects into the star’s structure and on its thermal evolutions (see e.g. Page 1989P89 ), their incorporation into the magnetic field evolution had not been properly accounted for. In the present paper we have been working with the GR version of the diffusion equation, Eq. 20, accompanied by proper boundary conditions, Eq. 21, derived in detail by GPZ00GPZ00 .
The analysis of GPZ00GPZ00 clearly showed that most of the effect of GR on the field decay is due to the presence of the red-shift factor $`e^\mathrm{\Phi }`$ in Eq. 20. This can be intuitively understood by noticing that the red-shift factor relates the proper time in each layer inside the star, i.e., the physical time in the layer where the currents are located and decaying, to the time of an observer who is observing the decay, at infinity (i.e., the coordinate time). As it can be easily seen from Eq. 20, the effects of the red shift factor on the field decay could be approximated by a flat spacetime diffusion equation running however on a slower time scale, given by some kind of averaged red-shift. This is clear from our Figures 5 and 6 where the GR curves have the same shape as the non-GR ones but shifted to the right.
Comparing the field evolution when GR effects are included with the evolution when they are neglected, we see that we obtain quite larger fields already in phase a, and in the late decay (phase c), we obtain much larger fields, by up to two orders of magnitude for very compact stars. GR effects in a field evolution can be viewed as almost equivalent to an evolution without GR effects but with currents initially located to higher densities and, in the late decay phase, with lower impurity content. This means that all constraints previously obtained about the location of the currents and the impurity contents of the crust are significantly weakened when GR effect are taken into account.
### 5.3 Comparison with the work of Sengupta
As far as the previously mentioned work of Sengupta (1998S98 ) is concerned, although we are in qualitative agreement with his results, we find many quantitative as well as interpretational differences. For instance this author considered a soft EOS (similar to the one we used) and field configurations with currents supposedly initially located at densities of $`2\times 10^{11}`$ and $`4\times 10^{11}`$ g cm<sup>-3</sup>. However, we obtain much higher densities at the depths at which he locates these layers: using his values of $`xr/R`$ = 0.979 and 0.9834 we find densities of the order of $`6\times 10^{13}`$ and $`2.6\times 10^{13}`$ g cm<sup>-3</sup> respectively. Consequently, his field decay curves should be compared with our class 2 models: we see then a rough quantitative agreement with respect to the importance of the GR effects with the significant difference that the late exponential decay (phase d) is absent in his GR models. As mentioned above, since the B-field evolutionary curves with and without GR effects should show approximately similar shapes we deduce that the different behavior of the $`B`$ field at late times obtained by Sengupta, must be due to numerical inaccuracies combined with his neglect of appropriate boundary conditions at the stellar surface as well as his employment of the Schwarzschild geometry.
### 5.4 Constraining the neutron star EOS and cooling history within the crustal magnetic field hypothesis
The analysis of the $`BP`$ diagrams of § 4.4 may be a tool to constrain the structure of matter at high density. In a similar analysis, UK97 concluded that a stiff EOS with ‘standard’ cooling is the most promising model for understanding the observed pulsar population properties and that a medium EOS requires currents located at relatively high densities and low impurity contents. These authors prefer the ‘standard’ cooling scenario on the basis that it implies an early decay of the field which may have an observational support coming the fact the young pulsars with an associated supernova remnant have stronger magnetic field than the bulk of the population. Our results, with GR included, show that a medium EOS with ‘standard’ cooling is also compatible with the observed $`P`$ and $`\dot{P}`$ with much weaker constraints on the initial dipole strength $`B_0`$ and penetration density $`\rho _0`$ and on the impurity content $`Q_{\mathrm{imp}}`$ than when GR effects are neglected. A soft EOS with ‘standard’ cooling needs very special conditions to accommodate the observational data: large $`B_0`$, high $`\rho _0`$, and very low $`Q_{\mathrm{imp}}`$.
But, notice that in case fast neutrino cooling is operating any EOS could be compatible with the observed pulsar population with only weak constraints on $`B_0`$, $`\rho _0`$, and $`Q_{\mathrm{imp}}`$ as we argued in § 4.4. However, this cooling scenario is more physical in the case of a soft EOS and probably incompatible with an EOS as stiff as our stiff EOS (for which the central density of a 1.4 $`M_{}`$star does not even reach twice the nuclear density). In summary, it appears rather difficult within the crustal field hypothesis alone to draw any strong conclusion about the EOS of neutron stars and their cooling histories.
### 5.5 The EOS and cooling history of neutron stars and the crustal magnetic field hypothesis
There exist independent arguments in favor of a not too soft neutron star EOS, none of them being, at the present time, compelling. Interpretation of kilohertz quasi-periodic oscillations (KHz QPO’s) in several low mass X-ray binaries indicate these systems contain neutron stars with masses around 2 $`M_{}`$(Kluźniak 1998K98 ) which a very soft EOS is not able to sustain. The strong gravitational light bending around a very compact star would make it almost impossible for such a star to show any modulation in its surface thermal emission, in contradistinction to some observations (Page & Sarmiento 1996PS96 ).
With regard to the thermal evolution of neutron stars, comparison of cooling models with current estimate of surface temperature of young neutron stars shows no clear evidence of occurrence of fast cooling (Ögelman 1995O95 ; Page 1998P98 ). This means that if these neutron stars do contain some ‘exotic’ phase of matter, its strong neutrino emission must be quenched by pairing (Page & Applegate 1992PA92 ; Page 1998P98 ) so that their thermal evolution is very close to the ‘standard’ one.
These two lines of arguments are consistent with the conclusions arising from the crustal field hypothesis, including GR effects: the EOS of neutron stars is away from the very soft regime and their thermal evolution is close to the prediction of the ‘standard’ model.
### 5.6 Joule heating and detectability of <br>Old Isolated Neutron Stars
Our study of the effect of the Joule heating produced by the decaying currents gives results very similar to the ones of Miralles et al (1998MUK98 ) and shows that the GR effects on the field evolution do not introduce any important change on the resulting late time thermal evolution.
Another consequence of the GR effects on the B-field evolution is that we can predict significantly stronger field strength over a Hubble time. This implies that Old Isolated Neutron Stars will be spinning very slowly and are hence more likely to be able to accrete matter from the interstellar medium. As a consequence this significantly increases the chances of detecting them through their thermal radiation, either due to the Joule heating or to the accretion.
## 6 Conclusions
The effect of GR on both the thermal and magnetic evolution is to slow it down, mostly because the proper time inside the star runs more slowly than the observer’s time. This effect is of course stronger the more compact the star is. On the other hand, with increasing compactness the field decay, as any diffusion process, is accelerated because of the reduction of the length scale. The competition of these two opposite tendencies reduces the sensitivity of the field evolution to the softness or stiffness of the EOS. As a result it is difficult to draw conclusions about the nature of the dense matter EOS from magnetic field evolution studies alone. However when taking into account information from other approaches, as the cooling history, gravitational lensing, kHz QPO’s, a consistent picture of neutron star structure and evolution is obtained, in which the EOS is not too soft and the cooling history is close to the so called ‘standard’ cooling scenario.
Reversely, we can consider these results as an argument in favor of the crustal magnetic field hypothesis. It is seen to be compatible with the observed distribution of pulsars in the $`P`$$`\dot{P}`$ diagram without requiring fine tuning of the models’ parameters. Moreover, these model parameters are consistent with values deduced from various other neutron star studies.
## Appendix A Induction Equation and Joule Heating Rate
In this appendix, we shall provide a few of the intermediate calculations leading to the derivation of the relativistic expression for the Joule heating employed in the numerical computations. We begin by recalling that the covariant form of Maxwell’s equations are as follows (Misner et al. 1973MTW73 ; Wald 1984W84 ):
$$^\alpha F_{\alpha \beta }=\frac{4\pi }{c}J_\beta $$
(33)
$$_{[\alpha }F_{\beta \gamma ]}=0$$
(34)
where $`F_{\alpha \beta }=F_{\beta \alpha }`$, $`J_\alpha `$ and $``$ are the coordinate components of the Maxwell tensor, the conserved four current and the covariant derivative operator, respectively. We may first recall that once a solution $`F_{\alpha \beta }`$ of Eqs. 3334 has been specified, an observer with four velocity $`U^\alpha ,U^\alpha U_\alpha =1`$ measures electric and magnetic fields $`(E,B)`$ given respectively by:
$$E_\alpha =F_{\alpha \beta }U^\beta ,B_\alpha =\frac{1}{2}ϵ_{\alpha \beta }^{}{}_{}{}^{\gamma \delta }F_{\gamma \delta }U^\beta $$
(35)
where $`ϵ_{\alpha \beta \gamma \delta }`$ stands for the four-dimensional Levi-Civita tensor density. It follows then easily from Eq. 35 that an inversion yields:
$$F_{\alpha \beta }=U_\alpha E_\beta U_\beta E_\alpha +ϵ_{\alpha \beta \gamma \delta }U^\gamma B^\delta $$
(36)
Thorne et al. (1982TM82 , 1986TPM86 ) have introduced an elegant version of curved spacetime electrodynamics, the absolute space formulation, that is reminiscent of the familiar flat spacetime electrodynamics formulated in terms of the electric and magnetic fields. This can be accomplished by working directly with the physical frame components of the electric and magnetic fields $`(E,B)`$, as determined by the static observers relative to their orthonormal frames. Recalling that such observers have a four velocity field $`U^a=e^\mathrm{\Phi }\delta _o^a`$, (see Eq. 1) then it follows that Maxwell’s Eqs. 3334 can be written in the following equivalent form:
$$E=4\pi \rho ,B=0$$
(37)
$$\times (ZB)=\frac{4\pi }{c}ZJ+\frac{1}{c}\frac{E}{t}$$
(38)
$$\times (ZE)=\frac{1}{c}\frac{B}{t}$$
(39)
where in this appendix $`Z`$ stands for the lapse function related to the red shift factor by $`Z=e^\mathrm{\Phi }`$ and the $``$ operators are formed out of the following three metric:
$$ds_{}^{2}{}_{3}{}^{}=\frac{dr^2}{12Gm/c^2r}+r^2d\mathrm{\Omega }^2.$$
(40)
by the rules of vector calculus as applied to the above orthogonal coordinate system. It follows then rather easily that in the absence of any convective motion Eqs. 3739, combined with Ohm’s law $`J=\sigma E`$, and within the MHD approximation (i.e., neglecting the displacement current in Eq. 38), yield the following induction equation governing the evolution of the neutron star magnetic field:
$$\frac{1}{c}\frac{B}{t}+\times [\frac{c}{4\pi \sigma }\times (ZB)]=0,$$
(41)
The above form of the induction equation combined with Eq. 40 as well the line element Eq. 1 with the metric functions corresponding to a non singular perfect fluid solution of Einstein’s equations yield Eq. 20 of the main text (for a detailed derivation see GPZ00GPZ00 ).
We now consider the GR expression for the Joule heating. To do so, we start from the covariant expression of the energy momentum tensor for an arbitrary electromagnetic field (Misner et al. 1973MTW73 ; Wald 1984W84 ) i.e.:
$$T_{\mu \nu }=\frac{1}{4\pi }\left[F_{\mu \gamma }F_{\nu }^{}{}_{}{}^{\gamma }\frac{1}{4}g_{\mu \nu }F_{\alpha \beta }F^{\alpha \beta }\right].$$
(42)
Using the representation of Eq. 36, one can easily show that the electromagnetic energy density $``$ as seen by the static observers is given by:
$$T_{\mu \nu }U^\mu U^\nu =\frac{1}{8\pi }[E^\alpha E_\alpha +B^\alpha B_\alpha ]=$$
$$\frac{1}{8\pi }[EE+BB]$$
(43)
where $`E`$ and $`B`$ stand for the physical components of the electric and magnetic fields, respectively, satisfying Eqs. 3739. Poynting’s theorem now has the following form
$$\frac{}{\tau }=𝒮+2g𝒮jE$$
(44)
where
$$𝒮\frac{c}{4\pi }(E\times B)$$
(45)
is the Poynting vector and
$$g\frac{Z}{Z}=\mathrm{\Phi }$$
(46)
is the gravitational acceleration, $`d\tau Zdt`$ being the proper time interval. The second term in the r.h.s. of Eq. 44 is a purely relativistic effect which results from the inertia of (electromagnetic) energy (see, e.g., Thorne et al 1986TPM86 ). The only dissipative term in Eq. 44 is the last one, i.e., the Joule heating term, using Ohm’s law:
$$Q_h^{(\mathrm{GR})}=\frac{jj}{\sigma }$$
(47)
By construction $`Q_h^{(\mathrm{GR})}`$ represents energy per unit (proper) time and (proper) volume.
Acknowledgments. This work was supported by a binational grant DFG (grant #444 - MEX - 1131410) - Conacyt (grant #E130.443), Conacyt (grant #2127P - E9507), UNAM - DGAPA (grant #IN105495) and Coordinación Científica - UMSNH. D.P. and T.Z. are thankful to the Astrophysikalisches Institut Potsdam for its kind hospitality and U.G. to the Instituto de Astronomía of UNAM.
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# MULTIPHOTON RESONANT TRANSITIONS OF ELECTRONS IN THE LASER FIELD IN A MEDIUM
## I Introduction
As is known the coherent interaction of electrons with a plane monochromatic wave in a dielectric medium can be described as a resonant scattering of a particle on the ”phase lattice” of a traveling wave similar to the Bregg scattering of the particle on the crystal lattice , .The latter is obvious in the frame of reference (FR) of rest of the wave. Since the index of refraction of a medium $`n>1`$ ( $`n(\omega )`$ $`n`$ as the wave is monochromatic) in this FR there is only a static periodic magnetic field and an elastic scattering of a particle takes place. The law of conservation for Cherenkov process taking into account the quantum recoil translates into the Bragg resonance condition between the de Broglie wave of the particle and this static periodic structure. Hence, in induced Cherenkov process the interaction resonantly connects two states of the particle which are degenerated over the longitudinal momentum with respect to the direction of the wave propagation: the states with longitudinal momenta $`p_x`$ \- of the incident particle and the states with longitudinal momenta $`p_x+\mathrm{}\mathrm{}k`$ \- of scattered ”Bragg” particle, as far as the conservation law of this process is $`\left|p_x\right|=\left|p_x+\mathrm{}\mathrm{}k\right|`$ ( $`\mathrm{}`$\- number absorbed or radiated photons with a wave vector $`k=k_x`$ ). The latter is the same as the Bragg condition of coherent elastic scattering. Therefore, in stimulated Cherenkov process no matter how weak the wave field is the usual perturbation theory is not applicable because of such degeneration of the states. So, the interaction near the resonance is necessary to describe by the secular equation . The latter, in particular, reveals zone structure of the particle states in the field of transverse electromagnetic (EM) wave in a dielectric medium , . Note that the application of the perturbation theory ignoring the mentioned degeneration in this process has reduced to essentially incorrect results which have been elucidated in the paper.
In the present work the case of strong radiation field is considered within the scope of the relativistic quantum theory for electron-laser interaction in a medium. Using the resonance approximation for the above mentioned two degenerated states in a monochromatic radiation field a nonperturbative solution of the Dirac equation (nonlinear over field solution of the Hill type equation) are obtained. The multiphoton probabilities of free electrons coherent scattering on a strong monochromatic wave at the Cherenkov resonance are counted, taking into account the above mentioned specificity of induced Cherenkov process , and spin-laser interaction as well. In the result of this resonant scattering the electron beam quantum modulation at high frequencies occurs that corresponds to the electron energy exchange at the coherent reflection from the ”phase lattice” of slowed wave in a medium. So, we can expect to have, in principle, a coherent X-ray source in induced Cherenkov process, since such-quantum modulated beam is a potential source of coherent radiation itself.
## II NONLINEAR SOLUTION OF THE DIRAC EQUATION FOR ELECTRON IN STRONG EM RADIATION FIELD IN A MEDIUM
In this section we shall solve Dirac equation for a spinor particle in the given radiation field in a medium
$$i\frac{\mathrm{\Psi }}{t}=\left[\widehat{\alpha }(\widehat{𝐩}e𝐀(\tau ))+\widehat{\beta }m\right]\mathrm{\Psi },$$
(1)
where
$$\widehat{\alpha }=\left(\begin{array}{cc}0& \sigma \\ \sigma & 0\end{array}\right);\widehat{\beta }=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)$$
(2)
are the Dirac matrices, with the $`\sigma `$ Pauli matrices, $`m`$ and $`e`$ are the mass and charge of a particle respectively (here we set $`\mathrm{}=c=1`$), $`\widehat{𝐩}=i\mathrm{}`$\- the operator of the generalized momentum $`𝐀=𝐀(tnx/c)`$-is the vector potential of a linearly polarized plane wave propagating in the $`OX`$ direction in a medium:
$$𝐀=\{0,A_0(\tau )\mathrm{cos}\omega \tau ,0\};\tau =tnx/c.$$
(3)
We shall assume that EM wave is adiabatically switched on at $`\tau =\mathrm{}`$ and switched off at $`\tau =+\mathrm{}`$ ($`𝐀(\tau =\mathrm{})=0`$).
To solve the problem it is more convenient to pass to the FR of rest of the wave ($`R`$ frame moving with the velocity $`V=1/n`$). As it is noticed, in this FR there is only static magnetic field which will be described according to (3) by the following vector potential
$$𝐀_R=\{0,A_0(x^{^{}})\mathrm{cos}k^{}x^{^{}},0\},$$
(4)
where
$$k^{}=\omega \sqrt{n^21}.$$
(5)
The wave function of a particle in $`R`$ frame is connected with the wave function in laboratory frame $`\mathrm{\Lambda }`$ by the Lorentz transformation of the bispinors
$$\mathrm{\Psi }=\widehat{S}(\vartheta )\mathrm{\Psi }_R,$$
(6)
where
$$\widehat{S}(\vartheta )=ch\frac{\vartheta }{2}+\alpha _xsh\frac{\vartheta }{2};th\vartheta =V=\frac{1}{n}$$
(7)
is the transformation operator. For $`\mathrm{\Psi }_R`$ we have the following equation
$$i\frac{\mathrm{\Psi }_R}{t^{^{}}}=\left[\widehat{\alpha }(\widehat{𝐩^{}}e𝐀_R(𝐱^{^{}}))+\widehat{\beta }m\right]\mathrm{\Psi }_R.$$
(8)
Since the interaction Hamiltonian does not depend on the time and transverse (to the direction of the wave propagation) coordinates the eigenvalues of the operators $`\widehat{H^{}}`$, $`\widehat{p}_y^{}`$, $`\widehat{p}_z^{}`$ are conserved: $`E^{}=const`$, $`p_y^{}=const`$, $`p_z^{}=const`$ and the solution of Eq.(8) can be represented in the form of a linear combination of free solutions of the Dirac equation with amplitudes $`a_i(x^{})`$ depending only on $`x^{}`$:
$$\mathrm{\Psi }_R(𝐫^{},t^{})=\underset{i=1}{\overset{4}{}}a_i(x^{})\mathrm{\Psi }_i^{(0)}.$$
(9)
Here
$`\mathrm{\Psi }_{1,2}^{(0)}=\left({\displaystyle \frac{E^{}+m}{2E^{}}}\right)^{\frac{1}{2}}\left[\begin{array}{c}\phi _{1,2}\hfill \\ \frac{\sigma _xp_x^{}+\sigma _yp_y^{}}{E^{}+m}\phi _{1,2}\hfill \end{array}\right]\mathrm{exp}\left[i(p_x^{}x^{}+p_y^{}y^{}E^{}t)\right],`$
$$\mathrm{\Psi }_{3,4}^{(0)}=\left(\frac{E^{}+m}{2E^{}}\right)^{\frac{1}{2}}\left[\begin{array}{c}\phi _{1,2}\hfill \\ \frac{\sigma _xp_x^{}+\sigma _yp_y^{}}{E^{}+m}\phi _{1,2}\hfill \end{array}\right]\mathrm{exp}\left[i(p_x^{}x^{}+p_y^{}y^{}E^{}t)\right],$$
(10)
where
$$p_x^{}=(E^{}_{}{}^{}2p_y^{}_{}{}^{}2m^2)^{\frac{1}{2}},\phi _1=\left(\begin{array}{c}1\\ 0\end{array}\right),\phi _2=\left(\begin{array}{c}0\\ 1\end{array}\right).$$
(11)
The solution of Eq. (8) in the form (9) corresponds to an expansion of the wave function in a complete set of the wave functions of a electron with certain energy and transverse momentum $`p_y^{}`$ (with longitudinal momenta $`\pm (E^{}_{}{}^{}2p_y^{}_{}{}^{}2m^2)^{\frac{1}{2}}`$ and spin projections $`S_x=\pm \frac{1}{2}`$). The latter are normalized to one particle per unit volume. Since there is symmetry with respect to the direction $`𝐀_R`$ (the $`OY`$ axis) we have taken, without loss of generality, the vector $`𝐩^{}`$ in the XY plane $`(p_z^{}=0)`$.
Substituting Eq.(9) into Eq.(8) then multiplying by the Hermitian conjugate functions and taking into account (10) and (2) we obtain a set of differential equations for the unknown functions $`a_i(x^{})`$. The equations for $`a_1`$, $`a_3`$ and $`a_2`$, $`a_4`$ are separated and for these amplitudes we have the following set of equations
$`p_x^{}{\displaystyle \frac{da_1(x^{})}{dx^{}}}=iep_yA_y(x^{})a_1(x^{})eA_y(x^{})\left(p_x^{}ip_y^{}\right)a_3(x^{})\mathrm{exp}(2ip_x^{}x^{}),`$
$$p_x^{}\frac{da_3(x^{})}{dx^{}}=iep_yA_y(x^{})a_3(x^{})eA_y(x^{})\left(p_x^{}+ip_y^{}\right)a_1(x^{})\mathrm{exp}(2ip_x^{}x^{}).$$
(12)
A similar set of equations is also obtained for the amplitudes $`a_2(x^{})`$ and $`a_4(x^{})`$. For simplicity we shall assume that before the interaction there are only electrons with specified longitudinal momentum and spin state, i. e.
$$\left|a_1(\mathrm{})\right|^2=1,\left|a_3(+\mathrm{})\right|^2=0,\left|a_2(\mathrm{})\right|^2=0,\left|a_4(+\mathrm{})\right|^2=0.$$
(13)
From the condition of conservation of the norm we have
$$\left|a_1(x^{})\right|^2\left|a_3(x^{})\right|^2=const$$
(14)
and the probability of reflection is$`\left|a_{3,4}(\mathrm{})\right|^2`$.
The application of the following unitarian transformation
$`a_1(x^{})=b_1(x^{})\mathrm{exp}\left(i{\displaystyle \frac{ep_y^{}}{p_x^{}}}{\displaystyle _{\mathrm{}}^x^{}}A_y(\eta )𝑑\eta i{\displaystyle \frac{\vartheta ^{}}{2}}\right)`$
$$a_3(x^{})=b_3(x^{})\mathrm{exp}\left(i\frac{ep_y^{}}{p_x^{}}_{\mathrm{}}^x^{}A_y(\eta )𝑑\eta +i\frac{\vartheta ^{}}{2}\right),$$
(15)
simplifies Eq.(12). Here $`\vartheta ^{}`$ is the angle between the momentum of electron and the direction of the wave propagation in the $`R`$ frame. The new amplitudes $`b_1(x^{})`$ and $`b_3(x^{})`$ satisfy the same initial conditions: $`\left|b_1(\mathrm{})\right|^2=1,`$ $`\left|b_3(+\mathrm{})\right|^2=0,`$ according to (13).
From Eq.(12) and Eq.(15) for the $`b_1(x^{})`$ and $`b_3(x^{})`$ we obtain the following set of equations
$`{\displaystyle \frac{db_1(x^{})}{dx^{}}}=f(x^{})b_3(x^{})`$
$$\frac{db_3(x^{})}{dx^{}}=f^{}(x^{})b_3(x^{})$$
(16)
where
$$f(x^{})=\frac{eA_y(t)p^{}}{p_x^{}}\mathrm{exp}\left(2ip_x^{}x^{}i\frac{2ep_y}{p_x^{}}_{\mathrm{}}^x^{}A_y(\eta )𝑑\eta \right);p^{}=\sqrt{p_y^{}_{}{}^{}2+p_x^{}_{}{}^{}2}$$
(17)
Using the following expansion by the Bessel functions
$`\mathrm{exp}\left(i\alpha \mathrm{sin}kx\right)={\displaystyle \underset{N=\mathrm{}}{\overset{\mathrm{}}{}}}J_N\left(\alpha \right)\mathrm{exp}\left(iNkx\right),`$
we can reduce Eq. (16) to the following form
$`{\displaystyle \frac{db_1(x^{})}{dx^{}}}={\displaystyle \underset{N=\mathrm{}}{\overset{\mathrm{}}{}}}f_N\mathrm{exp}\left[i(2p_x^{}Nk)x^{}\right]b_3(x^{})`$
$$\frac{db_3(x^{})}{dx^{}}=\underset{N=\mathrm{}}{\overset{\mathrm{}}{}}f_N\mathrm{exp}\left[i(2p_x^{}Nk)x^{}\right]b_1(x^{})$$
(18)
where
$$f_N=\frac{p^{}}{2p_y^{}}Nk^{}J_N\left(2\xi \frac{m}{p_x^{}}\frac{p_y^{}}{k}\right);\xi =eA/m$$
(19)
## III Resonant approximation for transition amplitudes
Because of conservation of particle energy and transverse momentum ( in $`R`$ frame) the real transitions in the field will occur from a $`p_x^{}`$ state to the $`p_x^{}`$ one and, consequently, the probabilities of multiphoton scattering will have a maximal values for the resonant transitions
$$2p_x^{}=sk^{},(s=\pm 1,\pm 2\mathrm{})$$
(20)
The latter expresses the condition of exact resonance between the electron de Broglie wave and the incident ”wave lattice”. In the $`\mathrm{\Lambda }`$ frame inelastic scattering takes place and the Eq.(20) corresponds to the well known Cherenkov conservation law
$$\frac{2E(1nv\mathrm{cos}\vartheta )}{(n^21)}=s\omega $$
(21)
where $`\vartheta `$ is the angle between the electron momentum and the wave propagation direction in the $`\mathrm{\Lambda }`$ frame (the Cherenkov angle), $`v`$ and $`E`$ are the electron velocity and energy.
So, we can utilize the resonant approximation keeping only resonant terms in the Eq.(18). Generally, in this approximation, at detuning of resonance $`\left|\delta _s\right|=\left|2p_x^{}sk^{}\right|<<k^{}`$ , we have the following set of equations for the certain $`s`$-photon transition amplitudes $`b_1^{(s)}(x^{})`$ and $`b_3^{(s)}(x^{})`$:
$`{\displaystyle \frac{db_1^{(s)}(x^{})}{dx^{}}}=f_s\mathrm{exp}\left[i\delta _sx^{}\right]b_3^{(s)}(x^{})`$
$$\frac{db_3^{(s)}(x^{})}{dx^{}}=f_s\mathrm{exp}\left[i\delta _sx^{}\right]b_1^{(s)}(x^{})$$
(22)
This resonant approximation is valid for the slow varying functions $`b_1^{(s)}(x^{})`$ and $`b_3^{(s)}(x^{})`$, i. e. at the following condition
$$\left|\frac{db_{1,3}^{(s)}(x^{})}{dx^{}}\right|<<\left|b_{1,3}^{(s)}(x^{})\right|k^{}.$$
(23)
At first we shall solve the case of exact resonance ($`\delta _s=0`$). According to the boundary conditions (14) we have the following solutions for the amplitudes
$$b_1^{(s)}(x^{})=\frac{ch\left[_x^{}^{\mathrm{}}f_s𝑑\eta \right]}{ch\left[_{\mathrm{}}^{\mathrm{}}f_s𝑑\eta \right]},b_3^{(s)}(x^{})=\frac{sh\left[_x^{}^{\mathrm{}}f_s𝑑\eta \right]}{ch\left[_{\mathrm{}}^{\mathrm{}}f_s𝑑\eta \right]}$$
(24)
and for the reflection coefficient
$$R^{(s)}=\left|b_3^{(s)}(\mathrm{})\right|^2=th^2\left[f_s\mathrm{}x^{}\right]$$
(25)
where $`\mathrm{}x^{}`$ is the coherent interaction length. The reflection coefficient in the laboratory frame of reference is the probability of absorption at $`v<1/n`$ or emission at $`v>1/n`$. The latter can be obtained expressing the quantities $`f_s`$ and $`\mathrm{}x^{}`$ by the quantities in this frame since the reflection coefficient is Lorentz invariant. So
$$R^{(s)}=th^2\left[F_s\mathrm{}\tau \right]$$
(26)
where
$$F_s=\left[\frac{(1nv\mathrm{cos}\vartheta )^2}{n^21}+v^2\mathrm{sin}^2\vartheta \right]^{1/2}\frac{s\omega }{2v\mathrm{sin}\vartheta }J_s\left(\frac{\xi }{n^21}\frac{2mv\mathrm{sin}\vartheta }{\omega (1nv\mathrm{cos}\vartheta )}\right)$$
(27)
and $`\mathrm{}\tau `$ for actual cases is the laser pulse duration in the $`\mathrm{\Lambda }`$ frame. The condition of applicability of this-resonant approximation (23) is equivalent to the condition
$$\left|F_s\right|<<\omega \text{,}$$
(28)
which restricts as the intensity of the wave as well as the Cherenkov angle. Besides, to satisfy the condition (28) we must take into account the very sensitivity of the parameter $`F_s`$ towards the argument of Bessel function, according to Eq.(27). For the wave intensities when $`F_s\mathrm{}\tau 1`$ the reflection coefficient is in the order of unit which can occur for the large number of photons $`s>>1`$ when the argument of Bessel function $`Zs1`$ in Eq.(27) (according to asymptotic behavior of Bessel function $`J_s(Z)`$ at $`Zs1`$).
Let us estimate the reflection coefficient of an electron from the laser pulse or the most probable number of absorbed/emitted photons due to resonance interaction in induced Cherenkov process. For the typical values of experimental parameters of this process in the gaseous medium with the index of refraction $`n110^4`$, at the initial electron energy $`E50Mev`$ and Cherenkov angle $`\vartheta 1mrad`$, during the ”Bragg reflection” from Neodymium laser pulse ($`\omega \mathrm{}\tau 10^2`$ , $`\mathrm{}\omega =1.17eV`$ ) with an intensity $`10^{10}W/cm^2`$ ($`\xi 10^4`$) electron absorbs or emits about $`10^5`$ photons.
For the off resonant solution, when $`\delta _s0`$, but $`f_s^2>\delta _s^2/4`$ from Eq.(22) for $`R^{(s)}`$ the following expression we obtain
$$R^{(s)}=\frac{f_s^2}{\mathrm{\Omega }_s^2}\frac{sh^2[\mathrm{\Omega }_s\mathrm{}x^{}]}{1+\frac{f_s^2}{\mathrm{\Omega }_s^2}sh^2[\mathrm{\Omega }_s\mathrm{}x^{}]}$$
(29)
where $`\mathrm{\Omega }_s=\sqrt{f_s^2\delta _s^2/4}`$, which has the same behavior as in the case of exact resonance . In opposite case when $`f_s^2\delta _s^2/4`$ the reflection coefficient is a oscillating function on interaction length.
During the coherent interaction with EM wave the quantum modulation of particles beam density occurs too which in difference to classic one after the interaction remains unlimitedly long (for the monochromatic beam). This is a result of coherent superposition of particle states with various energy and momentum due to absorbed and emitted photons in the radiation field which is conserved after the interaction. The quantum modulated state of the particle leads to modulation of the beam density after the interaction at the frequency of the stimulating wave and its harmonics . The density modulated particles beam can be used to generate spontaneous superradiation The various radiation mechanisms of quantum modulated beams are investigated in the works -.
In stimulated Cherenkov process the beam quantum modulation occurs if the particles wave packet size ($`\mathrm{\Delta }x`$ ) is enough large : $`\mathrm{\Delta }x>>\lambda `$ ( $`\lambda `$ -is the radiation wavelength ). In the opposite case the classic modulation or bunching of the beam takes place (klystron interaction scheme). From Eq. (9) and Eq. (26) for the electron wave function after the reflection from the wave pulse we have the superposition of incident and reflected electron waves (in the $`R`$ frame)
$$\mathrm{\Psi }_R=a_1(\mathrm{})\mathrm{\Psi }_1^{(0)}+a_3(\mathrm{})\mathrm{\Psi }_3^{(0)}$$
(30)
and in the result the probability density $`\rho _R=\mathrm{\Psi }_R^+\mathrm{\Psi }_R`$ is modulated at the X-ray frequencies
$$\rho _R^{(s)}=1+th^2\left[f_s\mathrm{}x^{}\right]+2\left[1\frac{p_x^{}_{}{}^{}2}{E^{}_{}{}^{}2}\right]th\left[f_s\mathrm{}x^{}\right]\mathrm{cos}(sk^{}x^{}\phi _0)$$
(31)
where
$`\left[1{\displaystyle \frac{p_x^{}_{}{}^{}2}{E^{}_{}{}^{}2}}\right]\mathrm{sin}\phi _0=\mathrm{sin}\vartheta ^{}.`$
In the laboratory frame of the reference from Eq.(7) and Eq.(30) we have
$$\rho ^{(s)}\frac{1}{\sqrt{n^21}}(1+th^2\left[F_s\mathrm{}\tau \right]+2th\left[F_s\mathrm{}\tau \right]\mathrm{cos}(s\omega \tau \vartheta ^{})$$
(32)
where it is taken into account that in actual case $`\left|s\omega \right|<<E`$. As is seen from Eq.(30) the modulation depth is in the order of unit for the intensities when $`F_s\mathrm{}\tau 1`$ which can be satisfied for the moderate intensities of the laser radiation in the order of $`10^{10}W/cm^2.`$
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# Deterministic walks in random media: evidence of generic scale invariance
\[
## Abstract
Deterministic walks over a random set of points in one and two dimensions ($`d=1,2`$) are considered. Points (“cities”) are randomly scattered in $`𝐑^d`$ following a uniform distribution. A walker (a “tourist”), at each time step, goes to the nearest neighbor city that has not been visited in the past $`\tau `$ steps. Each initial city leads to a trajectory composed of a transient part and a final $`p`$-cycle attractor. The distribution of transient times, $`p`$-cycles and number of cities per attractor are studied. It is shown numerically that transient times (for $`d=1,2`$) follow a Poisson law with a $`\tau `$ dependent decay but the density of $`p`$-cycles follows a power law $`D(p)p^{\alpha (\tau )}`$ for $`d=2`$. For large $`\tau `$, the expoent tends to $`\alpha 5/2`$. Some analytical results are given for the $`d=1`$ case. Since the power law is robust and does not depend on free parameters, this system presents “generic scale invariance”. Applications to animal exploratory behavior and other local minimization problems are suggested.
\]
The study of random walks has been very fruitful in physics and mathematics, and the theory of such stochastic processes is a well developed subject. The study of deterministic walks is also an interesting subject but presents the analytical difficulties common to the area of non-linear dynamical systems and has been less investigated . Here we propose a simple and intriguing problem, a deterministic walk over a random graph with $`N`$ nodes that is also an example of a local (“on-line”) optimization dynamics. It may be called the “local traveling salesman problem” or perhaps the “tourist problem” for short. The deterministic dynamics produces a division of the system phase space in several $`𝒪(N)`$ attractor basins which trap the walker (ergodicity is broken). The problem is reminiscent of walks in rugged landscapes but the equivalent of “local minima” are cycles instead of point attractors.
The model is defined as follows: points are randomly distributed with a uniform density $`\rho `$ in $`𝐑^d`$, where $`d`$ is the dimensionality of the space. These points may be thought as “cities” or “safe places” and they may be viewed as vertices of a random graph. At each time step, the “tourist” follows the deterministic rule: Go to the nearest city (or place) that has not been visited in the past $`\tau `$ time steps. Notice that the tourist wants to minimize only the distance to the next city (a local optimization procedure), not the sum of all distances in the trajectory or some other global cost function.
Our model can be of general interest for optimization theory with local constraints, and studies of deterministic dynamical systems with quenched disorder. However, we would like to suggest some specific motivations for considering this class of problems. The local optimization procedure could be naturally related to exploratory, foraging or migratory behaviors of animals. For example, rodents present two competing drives: an exploratory drive (“curiosity”) and a defensive behavior called thigmotaxis. The later refers to rodent aversion to open spaces and preference for places where its whiskers can touch vertical surfaces or objects, which provide protection . It is arguable that, for biological agents (and biologically inspired robots) it could be sometimes more important to minimize the distance traveled in each movement between two safe places instead of to optimize some global cost function. In another scale, the model could describe migratory or nomadic behaviors of humans, elephants, flamingos and other animals with well developed spatial memory. Cycles could be related to stable migratory routes on environments with localized resources, for example, oceanic islands, oasis and water holes. Local procedures are also usual in optimal foraging theory . The need for local optimization emerges due to short range sensorial capacities. Long distances may also imply non-additive costs: animals (and tourists) need safe places to stay during night, which puts a maximal distance that can be traveled at each time step.
Starting from a random city, the tourist performs a trajectory composed of a transient part and a final $`p`$-cycle attractor. In this letter we report the statistics for some relevant quantities similar to those measured in Kauffman networks : a) the probability $`P_\tau (t)`$ for obtaining a transient of size $`t`$, defined as the number of steps before the walker enters in some attractor (irrespective to the cycle period); for large $`t`$, it is Poisson-like $`P_\tau (t)\mathrm{exp}(t/\xi (\tau ))`$ with the decay time $`\xi (\tau )`$ growing exponentially for $`d=1`$ and linearly for $`d=2`$; b) the total density of attractors $`𝒟(\tau )`$, which decays exponentially for $`d=1`$ and as $`\tau ^1`$ for $`d=2`$; c) for $`d=2`$, the density of $`p`$-cycles $`D_\tau (p)`$ which follows a power law $`D_\tau (p)p^{\alpha (\tau )}`$ for $`\tau >0`$, with $`\alpha 5/2`$ for large $`\tau `$ (generic scale invariance); and d) the average number $`n_\tau (p)`$ of cities present in a $`p`$-cycle. Also some analytical results are obtained for the $`d=1`$ case.
The discrete time step is simply a label: it does not measure the actual physical time spent when the walker travels between the points. This independence makes irrelevant the density $`\rho `$ of cities because only relative distances are important (which city is the nearest) and not absolute distances, contrasting to standard random walks where the mean length step defines an intrinsic length scale. The only parameter is the memory window $`\tau `$. Self-avoidance is limited to this window and trajectories can intersect outside this range. If $`\tau =0`$ (no memory) the tourist goes simply to the nearest city until it finds two cities that are reciprocally nearest neighbors, entering in a $`2`$-cycle. In this simple case, attractors may be identified with geometrical (cluster) properties. For $`\tau =N1`$ the trajectory is totally self-avoiding and one has a kind of TSP nearest-neighbor algorithm . The interesting cases are the intermediate ones. For example, if $`\tau =1`$, the last visited city cannot be revisited, and only $`p`$-cycles with $`p3`$ can exist. For generic $`\tau `$, the relation $`p\tau +2`$ holds.
In the numerical experiments, $`N`$ points are randomly scattered following a uniform distribution in the interval $`[0,1]^d`$. Each point has $`d`$ spatial coordinates $`(x^1=x,x^2=y,\mathrm{},x^d)`$. The cities receive arbitrarily labeled as $`i=1,\mathrm{},N`$ and one constructs the Euclidean distance matrix $`𝐃`$ (for example, $`D_{ij}=\sqrt{(x_ix_j)^2+(y_iy_j)^2}`$ in $`d=2`$). The dynamics can be performed over the entries of this matrix instead of on real space. Starting from some city, one gets a transient trajectory until the walker enters some periodic attractor and is trapped. The number of steps before the walker enters the cycle defines the transient length $`t`$. The period $`p`$ and the number of different cities $`n`$ that pertain to the attractor are also determined. The same city can be visited more than once, thus $`np`$.
A finite size study showed that the behavior of the system is smooth as a function of $`1/N`$, so we have used $`N=3000`$ as a reasonable number for our simulations. Since each city is used as a starting point, a landscape with $`N`$ cities produces $`N`$ different transients. The statistics over $`N_R`$ realizations of sets of cities (“landscapes”) are collected. Unless stated otherwise, the results have been obtained by using $`N_R=100`$.
A natural question is if there is some critical $`\tau `$ that produces a phase transition, for example the emergence of an untraped (percolating) transient state. The distribution of transient times does not suggest this possibility because it is Poisson, $`P_\tau (t)\mathrm{exp}(t/\xi (\tau ))`$. This is shown for $`d=1`$ in Fig. 1 and for $`d=2`$ in Fig. 2. The exponential decay may be understood as follows. One finds numerically that the total number of attractors $`N_A(\tau )`$ is proportional to $`N`$, that is, the total density of attractors $`𝒟(\tau )=N_A/N`$ is constant and depends only on $`\tau `$. These attractors are scattered in phase space in a random uniform manner. Like points scattered randomly in space, one expects that the distribution of distances between attractors follows a Poisson law when these distances are larger than the attractor size $`p`$. By supposing that transient times are proportional to these distances, a Poisson law follows also for the transient times.
For $`d=1`$, the characteristic times grow as $`\xi (\tau )\mathrm{exp}(\gamma \tau )`$ (inset Fig. 1) and the total density decays as $`𝒟(\tau )\mathrm{exp}(\gamma ^{}\tau )`$ (Fig. 3a). For $`d=2`$, one observes the linear dependence $`\xi ((\tau )\tau ^\delta ,\delta =1.0`$ (inset Fig. 2); the total density decays as a power law $`𝒟(\tau )\tau ^\delta ^{}`$ (Fig. 3b). The average transient is proportional to the average distance between attractors, (which are inversely proportional to the attractor density). This means that $`𝒟(\tau )\xi (\tau )`$ should be constant, that is, $`\gamma =\gamma ^{}`$ and $`\delta =\delta ^{}`$. This is indeed the case (see Fig. 3c), but for $`d=1`$, although the exponential terms cancel, a linear residue remains. A better expression for the $`1D`$ decay time is $`\xi (\tau )=c\tau \mathrm{exp}(\gamma \tau )`$. We conjecture that the linear prefactor arises from the transient time spent in the cities of the attractor before it stabilizes (this time is larger for $`d=1`$ systems).
A property of natural interest is the density $`D_\tau (p)`$ of $`p`$-cycles, estimated as the number of different $`p`$-cycles divided by $`N`$, in the limit of very large systems. Evaluating this quantity requires careful enumeration because, when starting from all the possible initial states, one must not count the same attractor twice. Notice that $`𝒟(\tau )=_pD_\tau (p)`$. For $`d=1`$, $`D_\tau (p)`$ is certainly non Poisson, although the evidence for a power law is weak (Fig. 4a). For $`d=2`$, one observes clear power laws $`D_\tau (p)p^{\alpha (\tau )}`$ (Fig. 4b). The exponent $`\alpha (\tau )`$ stabilizes around $`\alpha =5/2`$ for large $`\tau `$ (inset Fig. 4b). Since there is no fine tuning of any explicit parameter in our system (such as $`\tau `$), the scale invariance is “generic” , that is, intrinsic to the problem. This is the most surprising result of our study. We conjecture that this scale invariance is related to two factors: a) a uniform distribution of points has a single length scale, $`\lambda =\rho ^{1/d}`$ but this length is irrelevant to the dynamics since only relative distances are considered when making a move; and b) the window $`\tau `$ defines a minimal length $`p_{min}=\tau +2`$ but not a maximal one. A clear explanation of this power law is still lacking.
We stress that this problem is not related only to geometrical properties since the cycles are appear only due to the introduced dynamics. Naively, one could think that a $`p`$-cycle is a geometrical object, for example a cluster where the distances between the points are smaller than any distance outside the cluster. This indeed is a sufficient but not necessary condition to obtain a $`p`$-attractor. For example, for $`d=2`$ (Fig. 5a), a walker with memory $`\tau =1`$ starts from city $`A`$ and finds the $`4`$-cycle $`ABCD`$. Although city $`E`$ is close to the cluster (since $`BE<AB`$), it is never visited because $`BC<BE`$ and $`CD<CE`$. However, if the tourist starts from city $`C`$, one gets a $`3`$-cycle that includes city $`E`$! This degeneracy and superposition of attractors can be understood observing that Fig. 5 shows trajectories in configuration space, not in phase space. In phase space, points corresponds to $`\tau +1`$-uples $`(𝐗_t,\mathrm{},𝐗_{t\tau })`$ where $`𝐗_t`$ is the position $`(x,y)`$ of the tourist at time $`t`$ and trajectories never intersect. Only for $`\tau =0`$ the configuration space is equivalent to the phase space.
Finally, we present the average number of cities $`n_\tau (p)`$ pertaining to cycles of period $`p`$ (Fig. 6). For $`d=1`$ there is almost no dispersion in the number of cities per attractor. A $`p`$-cycle has $`n(p)`$ cities. We also found that, for $`n>2(\tau +2)`$, the following relation holds for even cycles:
$$n_\tau (p)=p/2+\tau +1.$$
(1)
To see how this relation emerges, notice that for each $`\tau `$ there is a minimum cycle of period $`p_\tau =\tau +2`$, which we call a base block (Fig. 5b). A base block is composed of $`n_\tau =\tau +2`$ cities. The next cycles follow specific constructions (Fig. 5c). But when $`n`$ is large, geometrical constraints impose that the most common $`p`$-cycles are made of two base blocks (one in each attractor extremity) joined by $`n_I`$ intermediate cities, see Fig. 5d. An attractor with $`n`$ cities thus have $`n_I=n2n_\tau `$ intermediate points. These intermediate cities contribute to the total period with $`p_I=2n_I+2`$ steps (since for $`n_I=0`$, the joining of the base blocks contributes with two steps). Thus, the total period is $`p=2\times p_\tau +p_I=2(n\tau 1)`$, which leads to Eq. (1). This relation holds for cycles with $`n2n_\tau =2(\tau +2)`$, because only these cycles can incorporate two independent base blocks. For $`\tau =1`$, this is the unique conceivable manner of constructing cycles, meaning that odd cycles are prohibited (and also $`p=6`$ cycles, see Fig. 4a). For $`\tau >1`$, it is possible to construct odd cycles by using internal loops (an example with $`\tau =2`$ is given in Fig. 5e).
In $`d=2`$, the attractors are polygons with different forms and shapes so that this strict relation between periods and cities does not hold, although $`n`$ also scales linearly with $`p`$ (not shown). For $`\tau =1`$, one finds that odd cycles are less probable than even cycles (Fig. 4b), which is reminiscent of the $`d=1`$ behavior. Indeed, this occurs because elongated odd attractors in two-dimensional space are prohibited by the same geometrical constraints present in the one-dimensional case.
Another analytical result for the $`d=1`$ case can be obtained. Consider points $`x_i`$ randomly scattered along the real line, defining segments of size $`s_i=x_ix_{i1}`$. Without loss of generality, we assume that $`\rho =1`$, which means that $`s=1`$. It is easy to see that the distribution of interval sizes $`P(s)`$ follows a Poisson distribution $`P(s)=\mathrm{\Theta }(s)\mathrm{exp}(s)`$, where $`\mathrm{\Theta }(s)`$ is the Heaviside step function. For $`\tau =0`$, there exist only $`2`$-cycles attractors, which correspond to pairs of reciprocal nearest neighbors. The probability $`P_2`$ for this configuration is equal to the probability that $`s_{i1}>s_i`$ and $`s_{i+1}>s_i`$. Since $`s_{i1}`$ and $`s_{i+1}`$ are drawn independently, one gets: $`P_2=_{\mathrm{}}^{\mathrm{}}𝑑s_iP(s_i)P(s_{i1}>s_i|s_i)P(s_{i+1}>s_i|s_i)=_0^{\mathrm{}}𝑑s_i\text{e}^{s_i}\left(_{s_i}^{\mathrm{}}𝑑s\text{e}^s\right)^2=1/3`$, that is, on average, one third of the sequences of four points leads to reciprocal nearest neighbors and so to $`2`$-cycles. Since the number of sequences of four points is, in the large $`N`$ limit, equal to the number of points, one obtains $`D_0(2)=1/3`$. This has been fully confirmed by our numerical simulations (see caption Fig. 4).
The model may be generalized by introducing a stochastic component (a “temperature” $`T=1/\beta `$). For example, the probability for the tourist to travel from its present city $`j`$ to some city $`i`$ may be a function of the distance, say, $`P(ji)\mathrm{exp}(\beta D_{ij}/\lambda )`$, where $`\lambda =\rho ^{1/d}`$ normalizes the distances. In this case we expect, for $`T=\beta ^11`$, a punctuated-equilibrium behavior with sporadic transitions between attractor basins. It is an open question to determine if there is a critical temperature $`T_c`$ where full ergodicity is recovered.
Acknowledgments: O. Kinouchi thanks N. Caticha and the participants of the Statistical Physics of Neural Networks seminar at the Max Plank Institute for the Physics of Complex Systems, Dresden (march 1999) for useful discussions when this problem was being formulated. R. Vicente, N. Alves and U. P. C. Neves made useful suggestions for the final version of the manuscript. We also thank P. Stadler for discussing the problem and for the suggestion of the name “tourist problem”, made during his visit in Brazil. O. Kinouchi acknowledges support from FAPESP.
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# Enhanced four-wave mixing via elimination of inhomogeneous broadening by coherent driving of quantum transition with control fields
## Abstract
We show that atoms from wide velocity interval can be concurrently involved in Doppler-free two-photon resonant far from frequency degenerate four-wave mixing with the aid of auxiliary electromagnetic field. This gives rise to substantial enhancement of the output radiation generated in optically thick medium. Numerical illustrations addressed to typical experimental conditions are given.
PACS number(s): 42.50.Gy, 42.65.Dr, 42.65.Ky
Strong optical resonances inherent to free atoms and molecules are negated by the fact that only small fraction of the species can be concurrently resonance coupled in warm bulk gases. This is because of Maxwell distribution of the Doppler shifts of their resonances. Doppler-free (DF) coupling is usually achievable under equal frequency counter-propagating weak waves in two-photon-resonant ladder schemes. Because of the phase matching requirements, such schematic can not be implemented for far from frequency degenerated four-wave mixing (FWM). Second, the conditions of intermediate one-photon quasi-resonance can be very seldom satisfied in this case. Third, DF coupling vanishes with the growth of not only frequency difference but with strengths of the coupled fields too due to ac-Stark effect. Forth, DF coupling can not be routinely achieved in Raman schemes, like considered in this paper, because of the inherent frequency difference. Fifth, DF absorption does not indicate readily achievable DF FWM polarization in general case. Moreover, it is not obvious that increased FWM polarization would result in enhanced output of generated radiation, because of accompanying increased absorption.
This paper is aimed at demonstrating that the outlined limitations in optical physics can be removed by making use quantum coherence processes induced by an auxiliary intense control field. Considered effects lead to concurrent contribution of atoms from a wide velocity interval to the induced resonance and to eliminating it’s Doppler broadening under moderate light intensities. We investigate two- photon-resonant Raman-type FWM, controlled by an auxiliary driving field. The explicit formulae for power- dependent absorption/gain indices and for nonlinear FWM susceptibilities, accounting for interplay of power and Doppler shifts of the resonances and illustrating the major idea of the proposed method are derived. In order to avoid accompanying population transfer, which would complicate the formulae and would mask the major effect under consideration, the field coupled to the ground state is assumed to be weak.
A possible achievement of sub-Doppler resolution using intense control field was shown in . We propose a novel scheme, which enables to control FWM coupling with the aid of auxiliary electromagnetic (EM) field, taking no part in the FWM process itself. Accompanying increase of absorption of the fundamental radiation is considered too. As the outcome, substantial enhancement in quantum conversion efficiency in optically thick Doppler-broadened medium is shown. Numerical illustrations are given for the model, relevant to the FWM experiments with sodium dimer vapors . As a matter of fact that detuning from the intermediate resonance are larger than the Doppler width of the transition and the populated ground level is coupled to weak fields only, none of the CPT or EIT effects, usually employed for the enhancement of resonant FWM , are involved in the proposed technique.
Consider FWM process $`\omega _1\omega _2+\omega _3^+=\omega _S\omega _4`$ and transition configuration (fig.1a), similar to those studied in the experiments . However, in our case the EM radiation $`E_3(t,z)`$ consists of two components: weak $`E_3^+(\omega _3^+)`$ and counter propagating strong one $`E_3^{}(\omega _3^{})`$. Their frequencies can be the same or different. Strong radiation $`E_2`$ and weak $`E_1`$ co-propagate in the same direction as $`E_3^+`$. Only lower level remains populated, because $`E_1`$ is assumed so weak, that the populations can not be driven. Density matrix equations in the interaction representation are:
$`L_{01}\rho _{01}=i\{\rho _{00}V_{01}+\rho _{02}V_{21}\},`$ (1)
$`L_{03}\rho _{03}=i\{\rho _{00}V_{03}+\rho _{02}(V_{23}^++V_{23}^{})\},`$ (2)
$`L_{02}\rho _{02}=i\{\rho _{01}V_{12}+\rho _{03}(V_{32}^++V_{32}^{})\},`$ (3)
where $`L_{ij}=/t+𝐯+\mathrm{\Gamma }_{ij}`$, $`V_{ij}=G_{ij}\mathrm{exp}\{i(\mathrm{\Omega }_itk_iz)\}`$, $`V_{23}^\pm =G_{23}^\pm \mathrm{exp}\{i(\mathrm{\Omega }_3^\pm tk_3^\pm z)\}`$, $`G_{ij}=E_jd_{ij}/2\mathrm{}`$, $`G_{23}^\pm =E_3^\pm d_{23}/2\mathrm{}`$ \- are coupling Rabi frequencies, $`\mathrm{\Omega }_i`$ \- are corresponding resonance detunings (e.g., $`\mathrm{\Omega }_1=\omega _1\omega _{01}`$), $`\mathrm{\Gamma }_{ij}`$ \- homogeneous half-widths of the transitions.
As follows from (1) – (2), induced atomic coherence $`\rho _{02}`$ gives rise to the components in polarizations, responsible for novel effects in absorption and generation of the radiations under consideration. In the lowest order on the strength of the weak fields solution of the equations (1) – (3) can be found in the form:
$`\rho _{02}=r_{02}^{(1)}\mathrm{exp}\{i[(\mathrm{\Omega }_1\mathrm{\Omega }_2)t(k_1k_2)z]\}+r_{02}^{(4)}\mathrm{exp}\{i[(\mathrm{\Omega }_4\mathrm{\Omega }_3^+)t(k_4k_3^+)z]\}+\stackrel{~}{r}_{02}^{}\mathrm{exp}\{i[(\mathrm{\Omega }_4\mathrm{\Omega }_3^{})t(k_S+k_3^{})z]\},`$ $`\rho _{03}=r_{03}\mathrm{exp}\{i(\mathrm{\Omega }_4tk_4z)+\stackrel{~}{r}_{03}\mathrm{exp}\{i(\mathrm{\Omega }_4tk_Sz)\}+\stackrel{~}{r}_{03}^{}\mathrm{exp}\{i(\mathrm{\Omega }_4^{}tk_S^{}z)\}`$, $`\rho _{01}=r_{01}\mathrm{exp}\{i(\mathrm{\Omega }_1k_1z)\}+\stackrel{~}{r}_{01}\mathrm{exp}\{i(\mathrm{\Omega }_1t\stackrel{~}{k}_1z)\}+\stackrel{~}{r}_{01}^{}\mathrm{exp}\{i(\mathrm{\Omega }_1^{}t\stackrel{~}{k}_1^{}z)\},`$ where $`\mathrm{\Omega }_4^{}=\mathrm{\Omega }_1\mathrm{\Omega }_2+\mathrm{\Omega }_3^{},\mathrm{\Omega }_1^{}=\mathrm{\Omega }_4\mathrm{\Omega }_3^{}+\mathrm{\Omega }_2,k_S=k_1k_2+k_3^+,k_S^{}=k_1k_2k_3^{},\stackrel{~}{k}_1=k_4k_3^++k_2,\stackrel{~}{k}_1^{}=k_4+k_3^{}+k_2`$.
Equations for the density-matrix amplitudes become algebraic. With aid of solution for $`\stackrel{~}{r}_{03},r_{01}`$ and $`r_{03}`$ expressions for the susceptibilities, dressed by the strong fields $`E_2`$ and $`E_3^{}`$, can be routinely obtained and presented as:
$`\stackrel{~}{\chi }_4^{(3)}=i{\displaystyle \frac{d_{01}d_{12}d_{23}d_{30}/8\mathrm{}}{P_{01}\stackrel{~}{P}_{02}(P_{03}^++|G_{23}^{}|^2/P_{02}^{})}},`$ (4)
$`{\displaystyle \frac{\chi _1(\mathrm{\Omega }_1)}{\chi _{10}}}={\displaystyle \frac{\mathrm{\Gamma }_{01}}{P_{01}}}{\displaystyle \frac{P_{03}^{}P_{02}+|G_{23}^{}|^2}{P_{03}^{}\stackrel{~}{P}_{02}}},`$ (5)
$`{\displaystyle \frac{\chi _4(\mathrm{\Omega }_4)}{\chi _{30}}}={\displaystyle \frac{\mathrm{\Gamma }_{03}}{P_{03}}}{\displaystyle \frac{P_{01}^{}P_{02}^{}+|G_{12}|^2}{P_{01}^{}\{P_{02}^{}+|G_{23}^{}|^2/P_{03}+|G_{12}|^2/P_{01}^{}\}}},`$ (6)
where $`\chi _{01}`$ and $`\chi _{04}`$ are corresponding resonant values under the strong fields being turned off, $`P_{01}=\mathrm{\Gamma }_{01}+i(\mathrm{\Omega }_1k_1v),P_{01}^{}=\mathrm{\Gamma }_{01}+i(\mathrm{\Omega }_1^{}k_1^{}v),P_{03}=\mathrm{\Gamma }_{03}+i(\mathrm{\Omega }_4k_4v),P_{03}^{}=\mathrm{\Gamma }_{03}+i(\mathrm{\Omega }_4^{}k_4^{}v),P_{02}=\mathrm{\Gamma }_{02}+i[\mathrm{\Omega }_1\mathrm{\Omega }_2(k_1k_2)v],P_{02}^{}=\mathrm{\Gamma }_{02}+i[\mathrm{\Omega }_4\mathrm{\Omega }_3^{}(k_4+k_3^{})v],\stackrel{~}{P}_{02}=P_{02}+|G_{12}|^2/P_{01}+|G_{23}^{}|^2/P_{03}^{}`$, $`v`$ is projection of atom velocity on $`z`$. Difference between $`k_1`$ and $`\stackrel{~}{k}_1`$ as well as between $`k_4`$ and $`k_S`$ is neglected here.
With account of absorption but neglecting depletion of fundamental radiations due to FWM conversion, reduced equation for $`E_4`$ can be written as:
$`dE_4(z)/dz=\mathrm{i2}\pi k_4\stackrel{~}{\chi }_4^{(3)}E_1(0)E_2^{}E_3\mathrm{exp}(\mathrm{i}\mathrm{\Delta }Kz),`$ (7)
where $`\mathrm{\Delta }K=K_4K_1+K_2^{}K_3^+`$, $`K_j=k_j\mathrm{i}\alpha _j/2`$ – are complex wave numbers, $`\alpha _j`$ – power-dependent absorption indices. Quantum conversion efficiency ($`QCE`$) of $`E_1`$ into $`E_4`$ along the medium $`\eta _\mathrm{q}(z)`$ is given by the expression: $`\eta _\mathrm{q}(z)=(\omega _1/\omega _4)|E_4(z)/E_1(0)|^2\mathrm{exp}(\alpha _4z).`$
From (7) one obtains:
$`\eta _\mathrm{q}={\displaystyle \frac{\omega _1}{\omega _4}}{\displaystyle \frac{\left|2\pi \chi _4^{(3)}E_2E_3\right|^2}{|\mathrm{\Delta }K|^2}}\mathrm{exp}(\alpha _4z)\left|exp(\mathrm{i}\mathrm{\Delta }Kz)1\right|^2.`$ (8)
Pre-exponential factor can be expressed over Rabi frequencies, reduced nonlinear susceptibility and absorption indices considered below, ratios of the transition widths and $`|d_{03}|^2/|d_{01}|^2`$. The last factor is proportional to the ratio of the spontaneous relaxation rates. Thus $`QCE`$ can be found as absolute value, dependent on the optical thickness of the medium.
The major physics underlying the proposed technique is as follows. Modulation of the atomic wave-functions by the driving fields gives rise to the Autler-Townes splitting, which exhibits itself in our case as resonance shift. Besides intensities, the later depends on detunings of the driving fields and consequently – on the atomic velocities. It turns out that under appropriate intensities the resonances of atoms at different velocities can be shifted to the approximately same position. To illustrate that, consider one-photon detunings, substantially greater than corresponding Doppler HWHM. Then the resonance power-shift factors in (4) - (6) can be presented as: $`|G|^2/P(1+ikv/p)|G|^2/p`$, where $`p`$ is corresponding factor $`P`$ at $`v=0`$. This shows possible control of the resonance Doppler shifts through the power shifts. More details can be found in . In the same way a factor in the denominators of (4), (5), indicating dressed two-photon resonance, can be presented as:
$`\stackrel{~}{P}_{02}\stackrel{~}{\mathrm{\Gamma }}_{02}+i\stackrel{~}{\mathrm{\Omega }}_{02}i\{(|G_{12}|^2/\mathrm{\Omega }_1^2)k_1+`$
$`+(1+{\displaystyle \frac{|G_{23}^{}|^2}{(\mathrm{\Omega }_4^{})^2}})(k_1k_2){\displaystyle \frac{|G_{23}^{}|^2}{(\mathrm{\Omega }_4^{})^2}}k_3^{}\}v,`$ (9)
where $`\stackrel{~}{\mathrm{\Gamma }}_{02}`$ and $`\stackrel{~}{\mathrm{\Omega }}_{02}`$ give half-width and position of the induced resonance. As follows from (9), under proper choice of detuning, relative propagation direction and intensity of the control field $`E_3^{}`$, all Doppler shifts can be compensated by the power shifts in a such way, that dependence on $`v`$ vanishes in the given linear on $`v`$ approximation. This indicates trapping of all atoms, independent of their velocities in DF dressed two-photon resonance. It is seen, that as a matter of fact that $`k_2<k_1`$, elimination of Doppler broadening is not possible in the schematic under consideration with the aid of only driving field $`E_2`$. However it becomes possible with an auxiliary counter-propagating control field $`E_3^{}`$, which does not contribute directly in FWM because of phase mismatch. The equations (4), (5) show similar behavior of absorption index and nonlinear susceptibility near induced resonance.
While approaching closer to the intermediate resonances, required intensities become lower, but relative contribution of the neglected terms, proportional to the higher orders on $`k_iv/\mathrm{\Omega }_i`$, grows. This leads to decrease of the coherently coupled velocity interval.
The discussed outcomes can be illustrated with the numerical model of sodium dimer transitions : $`\lambda _{01}=661`$ nm, $`\lambda _{12}=746`$ nm, $`\lambda _{23}=514`$ nm and $`\lambda _{03}=473`$ nm. Corresponding homogeneous half-widths of the transitions are 20.69, 23.08, 18.30 and 15.92 MHz, Doppler $`HWHM`$ – 0.678, 0.601, 0.873 and 0.948 GHz.
Figure 2 depicts contribution of molecules at different velocities to the absorption index $`\alpha _1(\omega _1)Re\{\chi _1/\chi _{01}\}`$ and to nonlinear susceptibility (trivial Maxwell envelopes are removed), while conditions of elimination of Doppler broadening are fulfilled. The figure shows potentials of coherent coupling of molecules from wide velocity interval compared with the width of the Maxwell distribution, unlike the case in the absence of the control field. This gives rise to strong sub-Doppler resonances. Figure 3 shows modification of the effects while tuning closer to the intermediate resonances. Despite the growth of absorption $`\alpha _1`$, the proposed manipulating results in substantial increase in output of generated radiation at $`\omega _S`$ (fig.1 (b) and (c)).
In conclusion, we show that substantial enhancement in nonlinear-optical response of a Doppler broadened medium can be achieved by coherent driving of quantum transitions so that molecules from wide velocity interval become trapped to one and the same dressed two-photon Doppler-free resonance. The required intensities can be decreased by tuning driving frequencies closer to one-photon resonances, while the coupled velocity interval decreases too. For ladder-type schemes, where Doppler-broadening of two-photon resonances is much larger compared to Raman-like schemes, the considered effects are even more pronounced.
The authors thank B.Wellegehausen for encouraging discussions. This work was supported in part by the Krasnoyarsk Regional Science Foundation, by the Grant 97-5.2-61 in Fundamental Natural Sciences and by the Russian Foundation for Basic Research (Grant 99-02-39003).
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# A bright X-ray transient towards NGC 5128 (Centaurus A)Based in part on observations collected at the European Southern Observatory, La Silla, Chile
## 1 Introduction
The elliptical galaxy NCG 5128 is the stellar body of the giant double radio source Centaurus A (Cen A). It is one example of the family of elliptical galaxies that have an absorbing band of gas and dust projected across their stellar body, obscuring the nucleus at optical wavelengths. The dust lane is thought to be the remnant of a recent ($`10^710^8`$ years ago) merger of a giant elliptical galaxy with a smaller spiral galaxy (Thomson thomson92 (1992)).
Cen A as an active galaxy is usually classified as a FR I type radio galaxy, as a Seyfert 2 object in the optical (Dermer & Gehrels dermer95 (1995)), and as a “misdirected” BL Lac type AGN at higher energies (Morganti et al. morganti92 (1992)). It is one of the best examples of a radio-loud AGN viewed from the side of the jet axis (Graham graham79 (1979); Dufour & van den Bergh dufour79 (1979); Jones et al. jones96 (1996)).
Its proximity of $`<`$ 4 Mpc (Hui et al. hui93 (1993)) makes it uniquely observable among such objects, even though its bolometric luminosity is not large by AGN standards. Therefore, NGC 5128 is a very well studied and frequently observed galaxy in all wavelength bands. Its emission is detected from radio to high-energy gamma-rays (Johnson et al. johnson97 (1997); Israel israel98 (1998)) and it is often used as one of the first targets when new instruments or telescopes are to be tested (e.g. ESO VLT); even space based telescopes are tested on this active galaxy (e.g. HST (Schreier et al. schreier96 (1996)), Chandra (Kraft et al. kraft2000 (2000))) to demonstrate their resolution.
Variability of Centaurus A is reported in many wavelength regimes. In hard and in soft X-rays, observations of the Cen A region have revealed intensity variability greater than an order of magnitude (Bond et al. bond96 (1996); Baity et al. baity81 (1981); Turner et al. turner97 (1997)). Many of the observations referred to in the above publications, however, were made by instruments with a spatial resolution much less than required to resolve the inner parts of Cen A.
We here report the detection of a strong source only 2$`\stackrel{}{.}`$5 off from the nucleus of Cen A in ROSAT High Resolution Imager (HRI) X-ray data. Our data, obtained in 1995, have been recently re-analyzed, and were compared to archival data of all other ROSAT observations of NGC 5128. It turned out that the source is only present in the 1995 observations, implying strong variability. Recently published Chandra X-ray images of Cen A (on the Chandra X-ray observatory public WWW pages) show a weak source at the ROSAT position.
One important motivation for this *Letter* is to make the existence of a transient source so close to the direction of NGC 5128 known to all observers, inasmuch as all observations of variability (at least in soft X-rays) of Cen A, with instruments of low spatial resolution, must take the presence of this object into account.
## 2 Observations
In July 1995 a multiwavelength campaign took place to observe Cen A (NGC 5128) from radio to gamma-rays. One of the goals of this campaign was to look for correlated variability in different wavelength bands. Over a time interval of 14 days simultaneous measurements with various instruments were made (Steinle et al. steinle99 (1999)). The soft X-ray regime (0.1–2.4 keV) was observed five times with the ROSAT HRI during this time interval. The exposure times were in the order of 5 000 s each. Before and after these 1995 observations four other observations of Cen A were made with ROSAT (three with the HRI and one with the Position Sensitive Proportional Counter (PSPC)): 1990 (HRI, principal investigator (PI) E. Schreier), 1992 (PSPC, PI E. Schreier), 1994 (HRI, PI S. Döbereiner), and 1998 (HRI, PI S. Wagner). All other observations had substantially longer exposure times than the 1995 observations (see Table 1).
## 3 The transient X-ray source 1RXH J132519.8-430312
During the five ROSAT HRI observations in 1995, spanning a time of 10 days (see Table 1), in each observation a bright X-ray source was detected 2$`\stackrel{}{.}`$5 south-west of the nucleus of NGC 5128 at the outer regions of the elliptical galaxy. The coordinates derived from the ROSAT images are R.A. $`=13^\mathrm{h}25^\mathrm{m}19\stackrel{}{.}`$8 and Dec $`=43\mathrm{°}03\mathrm{}12\mathrm{}`$ (J2000; uncertainty $`5\mathrm{}`$), and the source has been assigned the ROSAT catalogue name 1RXH J132519.8-430312. Throughout the paper, however, we will refer to the transient as hcs113 as it is number 113 in a list of ROSAT HRI / Chandra sources which will be published in a future paper.
In all five observations hcs113 was the brightest point source in the HRI field-of-view with an average count rate of $`(0.033\pm 0.003)`$ counts sec<sup>-1</sup>, which is about a factor of 4 brighter than any other point source and about 30 % of the flux of the combined Cen A nucleus and jet sources. In Fig. 1 we show the sum of all five 1995 observations where hcs113 and the object H13 are marked. H13 is a source from Turner et al. (turner97 (1997))<sup>1</sup><sup>1</sup>1The paper of Turner et al. (turner97 (1997)) lists all ROSAT HRI count rates about 3 orders of magnitude too high, and Figure 4 of that paper shows the light-curve probably in units erg cm<sup>-2</sup> s<sup>-1</sup> and not in photons cm<sup>-2</sup> s<sup>-1</sup>. and is used for comparison, as it is a constant source in all ROSAT observations at a count rate of ($`0.008\pm 0.002`$) counts sec<sup>-1</sup>. H13 is probably a 14th mag M star (Feigelson et al. feigelson81 (1981)) or a distant early type galaxy (Wagner et al. wagner96 (1996)).
All other ROSAT observations from the years 1990, 1992, 1994, and 1998 show no trace of a source at the position of hcs113, even if combined. In Fig. 2, which is the sum of all ROSAT HRI observations without the 1995 multiwavelength data, with a total exposure of 101 862 s, the cuts are set as low as possible to detect any object at the transient’s position, but no object is detected. The derived (2 $`\sigma `$) upper limit for a detection at the position in question is 0.0003 counts sec<sup>-1</sup>. Due to the long exposure time and the low cuts, a large number of additional sources show up in this image, some of which are visible with different intensity or not at all in Fig. 1. These are candidates for other variable sources (which will be investigated in a future paper), but hcs113 is by far the object with the largest variation.
### 3.1 Simultaneous observations at other wavelengths
As the ROSAT observations in July 1995 were part of an extensive multiwavelength campaign to observe Cen A from radio to gamma-rays, simultaneous observations at other wavelengths exist (Steinle et al. steinle99 (1999)). For the measurements with sufficient spatial resolution (radio, optical, soft X-rays) the observations and analysis concentrated on the nucleus of the active galaxy (this is the reason why the transient went unnoticed until recently). The hard X-ray and gamma-ray results obtained with less spatial resolution (BATSE, OSSE, COMPTEL, and EGRET; all instruments on board the Compton Gamma Ray Observatory (CGRO)) were attributed to the nucleus by the assumption that it is the major source for the high-energy emission.
Besides ROSAT, only the optical monitoring with the ESO 2.2 m telescope at La Silla (Chile), had the position of hcs113 in the field-of-view and had enough spatial resolution. No candidate object in the ROSAT error box down to a limiting magnitude of $`18`$ mag is visible in the B and V exposures which were obtined during the same time interval in 1995 (see for example Fig. 3).
### 3.2 Time variability
As listed in Table 1, the measured flux from hcs113 is constant over the 10-day observation period in 1995. The only significant deviation occurs in the observation of July 22, where the flux drops by 30 %. All other ROSAT HRI observations made in 1990, 1992, 1994, and 1998 show no trace of the source. Turner et al. (turner97 (1997)), who analyzed the 1990 ROSAT HRI observation in detail, did not detect any source at the transient’s position, nor did they find any object in the ASCA and EINSTEIN data they used for comparison, which would be consistent with the position of hcs113.
The BATSE instrument on board CGRO has monitored Cen A continuously since its launch in April 1991 in the energy band 20–100 keV, with a very coarse (few degrees) spatial resolution (Wheaton et al. wheaton96 (1996)) giving a long baseline for monitoring flux variations. Around the time of the 1995 observations, an increase of the flux is present in the data, but such variations are very common for this source (see Fig. 1 in Steinle et al. (steinle99 (1999))), and similar flux increases occurred during some of the other ROSAT observations, when hcs113 was not detected.
Recently released public images of one of the first Chandra observations, which was imaging the Cen A region with unprecedented arc-second resolution in the energy range 0.1–10 keV, show a weak source at the position of hcs113. This is very probably the transient source either in its quiescent state or in a new outburst phase. However, hcs113 is not one of the highly variable or transient sources listed in the recent paper by Kraft et al. (kraft2000 (2000)).
### 3.3 Spectral information
As the ROSAT HRI has no energy resolution, no spectral information is available from those observations. Unfortunately the only PSPC observation and the data from the ROSAT All-Sky Survey, which would have had spectral information, were made when the transient was not active. Therefore only limited indirect information can be derived from the simultaneous measurements with the BATSE and OSSE instruments onboard CGRO which observed Cen A in the adjacent higher energy bands.
From the BATSE Cen A monitoring data, no strong enhancement in the 20–100 keV flux is probable (see above). The spectral photon index derived from the data during the 1995 campaign is between $`1.5`$ and $`1.7`$. A consistent spectral index of $`1.6`$ is derived from the OSSE data between 10 keV and 100 keV. These spectral index values are in agreement with the index of $`1.5`$ measured by Baity et al. (baity81 (1981)). The conclusion is that the emission of the transient is mainly at soft X-rays.
## 4 Summary and conclusions
During five ROSAT HRI observations in July 1995, a bright X-ray transient source 1RXH J132519.8-430312 ($`=`$ hcs113) was detected 2$`\stackrel{}{.}`$5 south-west of the nucleus of NGC 5128. When compared with all other ROSAT observations of Cen A made in 1990, 1992, 1994, and 1998 it turns out, that the source is only present during the 10 day period of the 1995 observations and no trace of it can be detected down to 1 % of the 1995 flux level in the other (deep) observations. Chandra observations of Cen A made in September 1999, however, show a source at the ROSAT position, which may be either the persistent conterpart of hcs113, or the transient at a recent active state.
If at 3.0 Mpc (the distance of Cen A), the luminosity of hcs113 in the energy band 0.1–2.4 keV is $`310^{39}`$ erg s<sup>-1</sup>, assuming a power-law spectrum with photon index $`1.5`$, N$`{}_{H}{}^{}=810^{20}`$ cm<sup>-2</sup> , and no additional intrinsic absorption in NGC 5128. This is above the Eddington luminosity of most X-ray binaries, but still within the observed range of luminosities for such objects. Therefore it cannot be ruled out that the transient is located in NGC 5128. On the other hand, if hcs113 is a Galactic object, then its soft X-ray luminosity would be less than $`310^{34}`$ erg s<sup>-1</sup> (assumed maximum distance 10 kpc). This is low for a typical neutron star transient, but is high for a typical cataclysmic variable. If it is closer than 10 kpc, its luminosity may be consistent with that of typical cataclysmic variables, whose optical counterparts are usually brighter than 19 or 20 mag, which then should be detectable in (deep) optical images. As the emission of the object seems to be mainly in the soft X-ray range, a distant AGN or Blazar is ruled out. A preliminary search for counterparts of hcs113 has been carried out using the NED (NASA/IPAC Extragalactic Database), the GSC (Hubble Guide Star Catalogue) and SIMBAD without any obvious object being found. It is, however, not the intent of this *Letter* to investigate the nature of this object in more detail.
If hcs113 is indeed an X-ray binary or other variable object, it would be very probable that the observed outburst was not a single event. Therefore, all Cen A observations in the past with low resolution X-ray instruments sensitive in soft X-rays (several keV), may have attributed a variation in the X-ray flux of the transient source to a variation in the flux from the nucleus (and/or jet) of Cen A.
Regardless of the nature and distance of the object, the fact that it is separated by only 2$`\stackrel{}{.}`$5 from the nucleus of Cen A, poses a strong problem for all spatially unresolved (on this scale) soft X-ray observations. Further observations and the interpretation of previous results (like the Cen A light-curves shown e.g. by Bond et al. (bond96 (1996)) and Turner et al. (turner97 (1997))) as well as the interpretation of the results of current Cen A monitoring programs in X-rays, as e.g. carried out by RXTE ASM and BATSE (CGRO) have to take the existence of the transient into account.
###### Acknowledgements.
This research has made use of the SIMBAD database, operated at CDS, Strasbourg, France, the NASA/IPAC Extragalactic Database (NED), which is operated by the Jet Propulsion Laboratory, Caltech, under contract with NASA, and the ROSAT Data Archive of the Max-Planck-Institut für extraterrestrische Physik, Garching, Germany. We thank the anomymous referee for the helpful comments.
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# Minimal Irreversible Quantum Mechanics: An Axiomatic Formalism.
## I Introduction.
Let us consider the function $`y=f(x)=x^2`$. Can we say if this function is an even function or an odd function? The primary (but incorrect) answer would be that it is an even function. This answer is wrong because, in order to define a function properly, we must also define its domain of definition $`D`$ and its range $`R`$, namely:
$`f(x)=y`$
$$f:DR$$
(1)
Then if $`y=f(x)=x^2`$ is defined as $`f:_{+\text{ }}`$it is an even function. But if it is defined as $`f:_+_+`$ the function is neither even nor odd. The morale of this story is that when we speak about a symmetry of a function necessarily we must define its domain of definition and its range, if not what we may say could be meaningless.
Let us now consider the Schroedinger equation:
$$i\frac{d|\psi >}{dt}=H|\psi >$$
(2)
and its solution:
$$|\psi (t)>=e^{iHt}|\psi (0)>$$
(3)
From what we have just learnt the question: ”is the set of the time evolutions obtained from Schroedinger equation time-reversible or invertible? , has no meaning if we do not define the domain of definition and the range of the states $`|\psi >`$ in eq. (3), i. e. the space where this vectors live. If we choose a Hilbert space $``$ the set of time-evolutions of eq. (3) is time-symmetric and each evolution is invertible. But if we make a different choice the set can become time asymmetric and each time-evolution can become non-invertible. Let us explain why it is so.
Let $`K`$ be the Wigner time-inversion operator. Hilbert space is invariant under time-inversion, namely:
$$K:$$
(4)
But we can choose a non-time-reversal-invariant space as the space of physical admissible state, let be $`\varphi _\text{ },`$ such that:
$$K:\varphi _{}\varphi _+\varphi _{}$$
(5)
and then within this space the set of time evolutions will turn out to be time-asymmetric and each evolution non invertible, as it is shown in the literature (precisely in almost all bibliographical references of this paper) and as we will also demonstrate below.
In this minimal way we can obtain a natural irreversible quantum mechanics. The aim of this paper is to sketch, using the results of many authors quoted in the references and our own results, an axiomatic formalism for this theory, which may have two possible advantages over ordinary quantum mechanics:
i.- Universe is clearly time asymmetric. Then the new theory may describe the real universe better than the usual one. We shall further discuss this possibility in section 16.
ii.- The new theory has more powerful spectral decompositions that makes the study of decaying processes easier
Let us rephrase what we have said using physical language: If we forget the time-asymmetric weak interaction (as it is usual in this kind of research, since weak interaction is so weak that it is difficult to see how it can explain the macroscopic time-asymmetry ), the time-asymmetry problem can be stated in the following question:
How can we explain the obvious time-asymmetry of the universe and most of its subsystems if the fundamental laws of physic are time-symmetric?
There are only two causes for asymmetry in nature: either the laws of nature are asymmetric or the solutions of the equations of the theory are asymmetric. As time asymmetry is not an exception the answer is contained in the question itself: If the laws of nature are time-symmetric essentially the only way we have to explain the time-asymmetry of the universe is to postulate that the state of the universe, or more generally, the space of physical possible solutions of the universe evolution equations is not time-reversal invariant, namely to use the second cause of asymmetry , . In this paper we explore this possibility using an axiomatic framework.
Moreover, certainly the best way to explain a physical idea is to construct an axiomatic structure because, having this structure, somehow we can see the whole idea, even if we cannot foresee all its consequences. Analogously it is easier to criticize an idea when it is presented in an axiomatic language. So we believe that this paper can clarify some issues of the problem of time-asymmetry.
The paper is organized as follows: Section 2: we define the space and the notation we will use. Section 3: the analytic continuation of the solutions is studied. Section 4: density matrices and Liouville space are introduced. Section 5: the space for the observables is chosen. Section 6: the axioms of the theory are stated. Sections 7 and 8: the main consequences of the axioms are obtained. Section 9: time-asymmetry and irreversibility are studied. Section 10: it is shown how Schrödinger and Heisenberg pictures work in the new formalism. Section 11: quantum equilibrium and decoherence are obtained. Section 12: it is shown that the norm and the energy are conserved and how Lyapunov variables appear. Section 13: entropy is defined. Section 14: the thermalization phenomenon is studied. Section 15: the global nature of the time-asymmetry is considered. Section 16: the Reichenbach diagram is presented. Section 17: other results are listed. Section 18: we draw our main conclusions. An appendix completes the paper.
## II Definition space.
Let us consider a quantum system, with a free hamiltonian $`H_0`$, endowed with a discrete plus a continuous spectrum, namely such that:
$`H_0|E_n^{(0)}>=E_n^{(0)}|E_n^{(0)}>`$
$$H_0|E^{(0)}>=E^{(0)}|E^{(0)}>$$
(6)
where, $`n=(0,1,2,\mathrm{}N_0),`$ $`0E^{(0)}\mathrm{},`$ $`E_n^{(0)}0.`$ The total hamiltonian will be $`H=H_0+V,`$ and the perturbation will be such that some bound states of the discrete spectrum become complex poles, namely we will have:
$`H|E_n>=E_n|E_n>`$
$$H|E>=E|E>$$
(7)
where, $`n=(0,1,2,\mathrm{}N),`$ $`0E\mathrm{},`$ and $`N<N_0`$ (in almost all the cases, for simplicity and in order to fix the ideas, we will consider that $`N=0`$, and therefore there is only one discrete ground state; a more general case will be considered in section 14)$`.\{|E_n>,E\pm >\}`$ is a basis of the corresponding Hilbert space $`.`$ (e.g. $`|E\pm >`$ can be the Lippmann-Schwinger retarded or advanced bases $`\{|\omega _\pm >\}`$ of ref. ), and:
$$<E_n|E_m>=\delta _{nm},\text{ }<E\pm |E^{}\pm >=\delta (EE^{}),\text{ }<E\pm |E_n>=0$$
(8)
$$I=\underset{n=0}{\overset{n=N}{}}|E_n><E_n|+_0^{\mathrm{}}|E\pm ><E\pm |dE$$
(9)
$$H=\underset{n=0}{\overset{n=N}{}}E_n|E_n><E_n|+_0^{\mathrm{}}E|E\pm ><E\pm |dE$$
(10)
Let $`\mathrm{\Xi }`$ be the vector space of all possible linear combinations of the basis $`\{|E_n>,|E\pm >\}`$ vectors so if $`|\psi >\mathrm{\Xi }:`$
$$|\psi >=\underset{n=1}{\overset{n=N}{}}\psi _n|E_n>+_0^{\mathrm{}}\psi _\pm (E)|E\pm >𝑑E$$
(11)
where neither $`\psi _n`$ nor $`\psi (E)`$ have any peculiar property.
Let $`K`$ be the Wigner time inversion operator therefore:
$$K|E_n>=|E_n>,\text{ }K|E\pm >=|E>$$
(12)
(for the continuous spectrum the Lippmann-Schwinger advanced and retarded bases has this property). Then, as $`K`$ is antilinear:
$$K|\psi >=\underset{n=1}{\overset{n=N}{}}\psi _n^{}|E_n>+_0^{\mathrm{}}\psi _\pm ^{}(E)|E>𝑑E$$
(13)
To find an irreversible quantum mechanics we must define a subspace $`\varphi _{}`$ of $`\mathrm{\Xi }`$ such that:
$$K:\varphi _{}\varphi _+\varphi _{}$$
(14)
namely a subspace which is not invariant under time inversions.
In our opinion nowadays there is a unique way to define space $`\varphi _{}`$ (see , , , ). In fact it is completely reasonable to ask that $`\varphi _{}`$ would have some logical properties namely that $`\varphi _{}`$ (i.e. $`\psi _nl^2,`$ $`\psi _\pm (E)L^2[0,\mathrm{})`$), $`\varphi _{}`$ must be dense in $`_{}`$ (the outgoing state subspace of $``$) and its topology must be a nuclear one. Precisely we will define the spaces $`_{}`$ and $`\varphi _{}`$ as:
$`|\psi >_{}\psi _+(E)H_{}^2|__+`$
$$|\psi >\varphi _{}\psi _+(E)𝒮H_{}^2|__+=𝒮_{}$$
(15)
where $`𝒮`$ is the Schwarz class function space (this choice allows to perform derivative to any order) and $`H_{}^2`$ is the space of Hardy class function from below (this choice introduces causality in our theory ). Nevertheless other choices have being used: , .
As $`\varphi _{}_{}`$ we have the Gel’fand triplet:
$$\varphi _{}_{}\varphi _{}^\times $$
(16)
where $`\varphi _{}^\times `$ is the space of antilinear functionals $`F`$ over $`\varphi _{}`$, such that :
$$F[\psi ]=<\psi |F>=<F|\psi >^{}$$
(17)
This will be the main arena of all our calculations. But, as we will see, we must also use the time inverted objects. Precisely, the spaces $`_+`$ and $`\varphi _+`$ defined as:
$`|\psi >_+\psi _{}(E)H_+^2|__+`$
$$|\psi >\varphi _+\psi _{}(E)𝒮H_+^2|__+=𝒮_+$$
(18)
where $`H_+^2`$ is the Hardy class from above, and the Gel’fand triplet is:
$$\varphi _+_+\varphi _+^\times $$
(19)
It is easy to see that the spaces $`\varphi _{}`$ and $`\varphi _+`$ satisfy eq. (14).
We close the section with three observations.
i.- Hamiltonian $`H`$ must be time-independent, since our aim is to define an arrow of time in a closed system. In fact, a realistic arrow of time must be defined using the whole universe as the system (open systems will be considered in section 14).
ii.- At first sight we can think that with the method we are about to propose we can define time asymmetry in non-interacting system like a free particle. It is not so since, even if the resulting free particle theory would formally be time-asymmetric, the entropy will not grow. In fact, the entropy will only grow, as we will see, if we have a non-trivial S-matrix with complex poles, which is not the case of a trivial free particle.
iii.- $`\varphi _{}`$ is dense in $`_{}`$ so, if someone would say that the ”real” physical states are those of $`_{},`$ we can answer that any one of these states can be approximated, as close as we wish, with a state of $`\varphi _{}.`$ So, on physical measurement grounds, the states of both spaces are indistinguishable. Nevertheless the two spaces have different kinds of topologies.
## III Analytic continuations.
Let us consider an scattering experiment using hamiltonian $`H`$ and let $`\{|\omega _+>\}`$ be the Lippmann-Schwinger basis (all objects related with this basis will be labelled with $`\omega ,`$ instead of the $`E`$ used in the equations of the last section), then we know that:
$$\underset{n=0}{\overset{N}{}}|\omega _n><\omega _n|+_0^{\mathrm{}}|\omega _+><\omega _+|d\omega =I$$
(20)
where the $`|\omega _n>`$ are the eventual stable bound states. Then:
$$<\phi |\psi >=\underset{n=0}{\overset{N}{}}<\phi |\omega _n><\omega _n|\psi >+_0^{\mathrm{}}<\phi |\omega _+><\omega _+|\psi >𝑑\omega $$
(21)
Let $`z_{n\text{ }}`$be the real and complex poles of the corresponding S-matrix. Then, using a simple analytic continuation of eq. (21) it can be demonstrated that if $`|\psi >\varphi _{}`$ and $`|\phi >\varphi _+,`$ the inner product $`<\phi |\psi >`$ (which is well defined since both vectors belong to $``$ ) reads:
$$<\phi |\psi >=\underset{n=0}{\overset{N_0}{}}<\phi |\overline{f_n}><\stackrel{~}{f_n}|\psi >+_\mathrm{\Gamma }<\phi |\overline{f_z}><\stackrel{~}{f_z}|\psi >𝑑z$$
(22)
where $`\mathrm{\Gamma }`$ is a curve that begins at $`O`$ and goes to $`+\mathrm{}`$ of the real axis under all the poles of the lower halfplane. Also, making the Nakanishi trick , we can obtain:
$$<\phi |\psi >=\underset{n=0}{\overset{N_0}{}}<\phi |\overline{f_n}><\stackrel{~}{f_n}|\psi >+_0^{\mathrm{}}<\phi |\overline{f_\omega }><\stackrel{~}{f_\omega }|\psi >𝑑\omega $$
(23)
where there is a term in the sum for each pole of the S-matrix, precisely: for each complex poles and each real pole corresponding to the bound states of the sum of eq. (20). Analogously it can be demonstrated that:
$$<\phi |H|\psi >=\underset{n=0}{\overset{N_0}{}}z_n<\phi |\overline{f_n}><\stackrel{~}{f_n}|\psi >+_0^{\mathrm{}}\omega <\phi |\overline{f_\omega }><\stackrel{~}{f_\omega }|\psi >𝑑\omega $$
(24)
(se also ) where $`|\overline{f_n}>,|\overline{f_\omega }>\varphi _+^\times ,`$ $`|\stackrel{~}{f_n}>,|\stackrel{~}{f_\omega }>\varphi _{}^\times `$ and in particular $`|\stackrel{~}{f_n}>=|\omega _n>`$ (if $`0nN)`$ and $`|\stackrel{~}{f}_\omega >=|\omega _+>`$ (eq. (44) ref. ).
Also if:
$`z_n`$ is real it is the eigenenergy of a bound state and if,
$`z_n`$ is complex it is a pole of the S-matrix
Therefore any $`|\psi >\varphi _{}`$ reads
$$|\psi >=\underset{n=0}{\overset{N_0}{}}|\overline{f_n}><\stackrel{~}{f_n}|\psi >+_0^{\mathrm{}}|\overline{f_\omega }><\stackrel{~}{f_\omega }|\psi >𝑑\omega $$
(25)
in a weak sense (namely premultiplied by any $`<\phi |\varphi _+)`$ and $`<\stackrel{~}{f_\omega }|\psi >𝒮_{}`$
Analogously:
$$H|\psi >=\underset{n=0}{\overset{N_0}{}}z|\overline{f_n}><\stackrel{~}{f_n}|\psi >+_0^{\mathrm{}}\omega |\overline{f_\omega }><\stackrel{~}{f_\omega }|\psi >𝑑\omega $$
(26)
Then, in an even weaker sense the two last equations can be written as:
$$I=\underset{n=0}{\overset{N_0}{}}|\overline{f_n}><\stackrel{~}{f_n}|+_0^{\mathrm{}}|\overline{f_\omega }><\stackrel{~}{f_\omega }|d\omega $$
(27)
$$H=\underset{n=0}{\overset{N_0}{}}z|\overline{f_n}><\stackrel{~}{f_n}|+_0^{\mathrm{}}\omega |\overline{f_\omega }><\stackrel{~}{f_\omega }|d\omega $$
(28)
The bases $`\{|\overline{f_n}>,|\overline{f_\omega }>\}`$ and $`\{|\stackrel{~}{f_n}>,|\stackrel{~}{f_\omega }>\}`$ are a biorthonormal system , namely:
$$<\stackrel{~}{f_n}|\overline{f_m}>=\delta _{nm},\text{ }<\stackrel{~}{f_n}|\overline{f_\omega }>=0,\text{ }<\stackrel{~}{f_\omega }|\overline{f_n}>=0,\text{ }<\stackrel{~}{f_\omega }|\overline{f_\omega ^{}}>=\delta (\omega \omega ^{})$$
(29)
Also it can be proved that:
$$<\overline{f_n}|\overline{f_m}>=\delta _{nm}\epsilon _n$$
(30)
$$<\stackrel{~}{f_n}|\stackrel{~}{f_m}>=\delta _{nm}\epsilon _n$$
(31)
where $`\epsilon _n=1`$ if $`Imz_n=0`$ and $`\epsilon _n=0`$ otherwise , , Namely the states with $`Imz_n0`$ are ”ghosts” with vanishing norm. This fact is evident since if $`|\overline{f_n}>`$ is one of these ghosts, from eq. (20), we have:
$$<\overline{f_n}|\overline{f_n}>=<\overline{f_n}|(\underset{n=0}{\overset{N}{}}|\omega _n><\omega _n|+_0^{\mathrm{}}|\omega _+><\omega _+|d\omega )|\overline{f_n}>=$$
(32)
$`=<\overline{f_n}|({\displaystyle \underset{n=0}{\overset{N}{}}}|\stackrel{~}{f_n}><\stackrel{~}{f_n}|+{\displaystyle _0^{\mathrm{}}}|\stackrel{~}{f}_\omega ><\stackrel{~}{f}_\omega |d\omega )|\overline{f_n}>=0`$
where we have used eq. (27) and that $`|\stackrel{~}{f_n}>=|\omega _n>`$and $`|\stackrel{~}{f}_\omega >=|\omega _+>`$ (eq. (44) of ref. ).
We will, sometimes, find useful write all these equation using a shorthand notation where we will call the basis $`\{|\omega _0>,|\omega +>\}`$ just $`\{|i>\},`$ the basis $`\{|\overline{f_n}>,|\overline{f_\omega }>\}`$ just $`\{|\overline{i}>\},`$ and the basis $`\{|\stackrel{~}{f_n}>,|\stackrel{~}{f_\omega }>\}`$ just$`|\{|\stackrel{~}{i}>\}.`$ Then eq. (25) reads:
$$|\psi >=\underset{i}{}|\overline{i}><\stackrel{~}{i}|\psi >=\underset{i}{}\psi _i|\overline{i}>$$
(33)
and also we will conventionally say that $`<\stackrel{~}{i}|\psi >𝒮_{}.`$ Eq. (26) reads:
$$H|\psi >=\underset{i}{}z_i|\overline{i}><\stackrel{~}{i}|\psi >$$
(34)
In all these equations $`|\overline{i}>\varphi _+^\times ,|\stackrel{~}{i}>\varphi _{}^\times .`$ The biorthonormality of the system $`\{|\overline{i}>\}`$ and $`\{|\stackrel{~}{i}>\}`$ will be symbolized as:
$$<\stackrel{~}{i}|\overline{j}>=\delta _{ij}$$
(35)
$$\underset{i}{}|\overline{i}><\stackrel{~}{i}|=I$$
(36)
where the symbols have an obvious meaning (e. g. eq. (36) is a shorthand-notation weak version of eq. (27), etc.).
Also:
$$<\overline{i}|\overline{j}>=\delta _{ij}\epsilon _i$$
(37)
$$<\stackrel{~}{i}|\stackrel{~}{j}>=\delta _{ij}\epsilon _i$$
(38)
where $`\epsilon _i=1`$ if $`Imz_i=0`$ and $`\epsilon _i=0`$ in all the other cases.
## IV Density matrices.
Up to now we have just introduced pure states, but we can rephrase everything using mixed states $`\rho `$. In general $`\rho \mathrm{\Xi }\mathrm{\Xi }`$ but usually it is considered that it belongs to a Liouville space $`^{}=.`$ The time evolution of the mixed states can be obtained solving Liouville equation:
$$i\frac{d\rho }{dt}=[H,\rho (t)]=L\rho (t)$$
(39)
where $`L`$ is the Liouville operator.
From this equation we see that any $`\rho _{}`$ that commutes with $`H`$ is an stationary state. This state is diagonal in the same basis than $`H`$ and therefore it can be written as (cf. eq. (10)):
$$\rho _{}=\rho _0|\omega _0><\omega _0|+_0^{\mathrm{}}\rho _\omega |\omega +><\omega +|d\omega $$
(40)
where we have made $`N=0`$ for simplicity as we have announced. The second term of the r. h. s. of the last equation implies the existence of a singular structure in the stationary state, that was studied at large in the paper . So, if we want to develop a rigorous treatment of this singular structure we are forced to consider that $`=()`$, where the first $``$ contains the singular structure and the second factor $``$ is the usual Liouville space $`^{},`$ that now it will be only considered as a regular structure, and therefore we introduce the following eigenbasis of $`L`$:
$`\rho (0)=|\omega _0><\omega _0|,\text{ }\rho (0,\omega )=|\omega _0><\omega +|,\text{ }\rho (\omega ,0)=|\omega +><\omega _0|`$
$$\beta (\omega )=|\omega +><\omega +|,\text{ }\rho (\omega ,\omega ^{})=|\omega +><\omega ^{}+|$$
(41)
This is an orthonormal basis in an inner product that we will define below (cf. eq. (46)).
We can now compute the eigenvalues of the eigenvalues of $`L:`$
$`L\rho (0)=0,\text{ }L\rho (0,\omega )=(\omega _0\omega )\rho (0,\omega ),\text{ }L\rho (\omega ,0)=(\omega \omega _0)\rho (\omega ,0)`$
$$L\beta (\omega )=0,\text{ }L\rho (\omega ,\omega ^{})=(\omega \omega )\rho (\omega ,\omega ^{})$$
(42)
But for $`\rho (\omega ,\omega ^{})`$ it is better to use Riesz quantum numbers:
$`\sigma ={\displaystyle \frac{1}{2}}(\omega +\omega ^{}),\text{ }0\sigma <\mathrm{}`$
$$\nu =\omega \omega ^{},\text{ }2\sigma \nu 2\sigma $$
(43)
So we will write the matrices $`\rho (\omega ,\omega ^{})`$ as:
$$\rho (\omega ,\omega ^{})=\beta (\sigma ,\nu )$$
(44)
So any $`\rho `$ can be written as:
$$\rho =\rho _0\rho (0)+_0^{\mathrm{}}[\rho _{0\omega }\rho (0,\omega )+\rho _{\omega 0}\rho (\omega ,0)+\rho _\omega \beta (\omega )]𝑑\omega +_0^{\mathrm{}}𝑑\sigma _{2\sigma }^{2\sigma }𝑑\nu \rho _{\sigma \nu }\beta (\sigma ,\nu )$$
(45)
The inner product among these $`\rho `$ is naturally defined as:
$$(\rho |\rho ^{})=\rho _0^{}\rho _0^{}+_0^{\mathrm{}}[\rho _{0\omega }^{}\rho _{0\omega }^{}+\rho _{\omega 0}^{}\rho _{\omega 0}^{}+\rho _\omega ^{}\rho _\omega ^{}]𝑑\omega +_0^{\mathrm{}}𝑑\sigma _{2\sigma }^{2\sigma }𝑑\nu \rho _{\sigma \nu }^{}\rho _{\sigma \nu }^{}$$
(46)
From eq. (42), we have.
$$L=_0^{\mathrm{}}(\omega _0\omega )[\rho (0,\omega )\rho ^{}(0,\omega )\rho (\omega ,0)\rho ^{}(\omega ,0)]𝑑\omega +_0^{\mathrm{}}𝑑\sigma _{2\sigma }^{2\sigma }\nu 𝑑\nu \beta (\sigma ,\nu )\beta ^{}(\sigma ,\nu )$$
(47)
Now let us make the analytic continuation.:
The diagonal elements $`\rho (0)`$ and $`\beta (\omega )`$ will remain untouched, since they correspond to the stationary state, but we will require that :
$$\rho _\omega 𝒮.$$
(48)
The terms $`\rho (0,\omega )`$ and $`\rho (\omega ,0)`$ can be treated as in the last section, so they have only one variable $`\omega ,`$ so we will ask that:
$$\rho _{\omega 0}𝒮H_{}^2|__+=𝒮_{},\text{ }\rho _{0\omega }𝒮H_+^2|__+=𝒮_+$$
(49)
Finally let us consider the term $`\beta (\sigma ,\nu ).`$ We could promote both real variables $`\sigma `$ and $`\nu `$ to complex variables but, as $`\nu `$ is the eigenvalue of the Liouville operator, it is only necessary to promote $`\nu z`$ and leave $`\sigma `$ real. Precisely as:
$$\rho _{\sigma \nu }=(\beta (\sigma ,\nu )|\rho )$$
(50)
we can consider the complex value function of $`z:`$
$$\rho _{\sigma z}=(\beta (\sigma ,z)|\rho )$$
(51)
and to ask that:
$$\rho _{\sigma \nu }𝒮H_{}^2|_{2\sigma }^{2\sigma }=𝒮_{}^{(\sigma )}$$
(52)
for any $`\sigma 0,`$ thus $`\rho _{\sigma z}`$ with be an analytical function of $`z`$ in the lower halfplane. Then we will say that $`\rho \mathrm{\Phi }_{}`$ if eqs. (48), (49), and (52) are $`\rho _{\omega 0}𝒮H_{}^2|__+=𝒮_{},`$ $`\rho _{0\omega }𝒮H_+^2|__+=𝒮_+`$. We also define a space $`_{}`$ such that if $`\rho _{}`$ we simply have :
$`\rho _{\omega 0}H_{}^2|__+,\text{ }\rho _{0\omega }H_+^2|__+`$
$`\rho _{\sigma \nu }H_{}^2|_{2\sigma }^{2\sigma }`$
Let us now define the time-inverted spaces $`\mathrm{\Phi }_+`$ and $`_+.`$ If eq. (48) is satisfied and :
$$\rho _{\omega 0}𝒮H_+^2|__+=𝒮_+,\text{ }\rho _{0\omega }𝒮H_{}^2|__+=𝒮_{}$$
(53)
$$\rho _{\sigma \nu }𝒮H_+^2|_{2\sigma }^{2\sigma }=𝒮_+^{(\sigma )}$$
(54)
for any $`\sigma 0,`$ thus in this case $`\rho _{\sigma z}`$ will be an analytical function of $`z`$ in the upper halfplane, then we will say that $`\rho \mathrm{\Phi }_+`$ (but in the definition of this space the basis $`|\omega +>`$ of eq. (41) must be changed by the basis $`|\omega >).`$
We also define a space $`_+`$ such that if $`\rho _+`$ we have:
$`\rho _{\omega 0}H_+^2|__+,\text{ }\rho _{0\omega }H_{}^2|__+`$
$`\rho _{\sigma \nu }H_+^2|_{2\sigma }^{2\sigma }`$
where we have also changed the basis as before.
Let us now consider the poles:
For the terms $`\rho (0,\omega )`$ and $`\rho (\omega ,0)`$ we will find those of the last section, and we can repeat the analytic continuation up to the curve $`\mathrm{\Gamma }`$ of section 3.
For the terms $`\beta (\omega )`$ and $`\beta (\sigma ,\nu )`$, for some fixed $`\sigma `$ and for every pole $`z_n`$ of the S-matrix we will find at the $`\nu `$ or $`z`$plane two poles $`\pm 2(z_n\sigma )`$ (and also a pole at $`\nu =z=0`$ coming from the singular structure of the continuous field, the $`\beta (\omega )`$ term), that we will call $`\zeta _l.`$ Also it may happen that for some $`\sigma _j`$ extra poles $`\zeta _l^j`$ may appear . So introducing a curve $`C,`$ in the lower halfplane, that goes, under all the poles, from $`2\sigma `$ to $`2\sigma `$ of the real axis (fig. 1) and using, as in the pure states case, the Cauchy theorem, if $`\rho \mathrm{\Phi }_+`$ and $`\rho ^{}\mathrm{\Phi }_{},`$ we obtain that:
$`(\rho |\rho ^{})=(\rho |\{\rho (0)\rho ^{}(0)+`$
$`+{\displaystyle \underset{n}{}}[\rho (\overline{z_{n,}0})\rho ^{}(\stackrel{~}{z_n,0)}+\rho (\overline{0,z_n})\rho ^{}(\stackrel{~}{0,z_n})]+{\displaystyle _\mathrm{\Gamma }}[\rho (\overline{z,0})\rho ^{}(\stackrel{~}{z,0})+\rho (\overline{0,z})\rho ^{}(\stackrel{~}{0,z})]dz+`$
$$\underset{jl}{}\overline{\beta (\sigma _{j,}\zeta _l^j)}\stackrel{~}{\beta ^{}(\sigma _{j,}\zeta _l^j)}+_0^{\mathrm{}}d\sigma [\beta (\sigma )\beta ^{}(\sigma )+\underset{l}{}\beta (\overline{\sigma ,\zeta _l})\beta ^{}\stackrel{~}{(\sigma ,\zeta _l)}+_C\beta (\overline{\sigma z})\beta ^{}(\stackrel{~}{\sigma z})dz]\}|\rho ^{})$$
(55)
Namely in weak sense:
$`I=\rho (0)\rho ^{}(0)+{\displaystyle \underset{n}{}}[\rho (\overline{z_{n,}0})\rho ^{}(\stackrel{~}{z_n,0})+\rho (\overline{0,z_n})\rho ^{}(\stackrel{~}{0,z_n})+{\displaystyle _\mathrm{\Gamma }}[\rho (\overline{z,0})\rho ^{}(\stackrel{~}{z,0})+\rho (\overline{0,z})\rho ^{}(\stackrel{~}{0,z})]dz+`$
$$\underset{jl}{}\overline{\beta (\sigma _{j,}\zeta _l^j)}\stackrel{~}{\beta ^{}(\sigma _{j,}\zeta _l^j)}+_0^{\mathrm{}}d\sigma [\beta (\sigma )\beta ^{}(\sigma )+\underset{l}{}\beta (\overline{\sigma ,\zeta _l})\beta ^{}\stackrel{~}{(\sigma ,\zeta _l})+_C\beta (\overline{\sigma z})\beta ^{}(\stackrel{~}{\sigma z})dz]$$
(56)
In these equations the presence of the poles coming from the singular structure ( in each $`\sigma =const.`$ plane) is represented by the terms $`\beta (\sigma )\beta ^{}(\sigma ).`$
Then we can write any $`\rho \mathrm{\Phi }_{}`$ as:
$`\rho =\rho _0\rho (0)+{\displaystyle \underset{n}{}}[\rho _{n0}\rho (\overline{z_{n,}0})+\rho _{0n}\rho (\overline{0,z_n})+{\displaystyle _\mathrm{\Gamma }}[\rho _{z0}\rho (\overline{z,0})+\rho _{0z}\rho (\overline{0,z})]dz`$
$$\underset{jl}{}\overline{\rho _{jl}\beta (\sigma _{j,}\zeta _l^j)}+_0^{\mathrm{}}𝑑\sigma [\rho _\sigma \beta (\sigma )+\underset{l}{}\rho _{\sigma l}\beta (\overline{\sigma ,\zeta _l})+_C\rho _{\sigma z}\beta (\overline{\sigma z})𝑑z]$$
(57)
Analogously, from the analytic continuation of the Liouville operator we obtain:
$`(\rho |L|\rho ^{})=(\rho |\{{\displaystyle \underset{n}{}}[(z_n\omega _0)\rho (\overline{z_{n,}0})\rho ^{}(\stackrel{~}{z_n,0)}+(\omega _0z_{n)}^{}\rho (\overline{0,z_n})\rho ^{}(0,z_n)]+`$
$`+{\displaystyle _\mathrm{\Gamma }}[(z\omega _0)\rho (\overline{z,0})\rho ^{}(\stackrel{~}{z,0})+(\omega _0z^{})\rho (\overline{0,z})\rho ^{}(\stackrel{~}{0,z})]𝑑z+`$
$$\underset{jl}{}\overline{\beta (\sigma _{j,}\zeta _l^j)}\stackrel{~}{\beta ^{}(\sigma _{j,}\zeta _l^j)}+_0^{\mathrm{}}d\sigma [\underset{l}{}\zeta _l\beta (\overline{\sigma ,\zeta _l})\beta ^{}(\stackrel{~}{\sigma ,\zeta _l})|\rho ^{})+_Cz\beta (\overline{\sigma ,z})\beta ^{}(\stackrel{~}{\sigma ,z})|\rho ^{})dz]\}|\rho ^{})$$
(58)
where $`\zeta _l=\overline{\nu }_li\gamma _l,`$ $`\gamma _l0`$ and $`\zeta _l^j=\overline{\nu }_l^ji\gamma _l^j,`$ $`\gamma _l^j0`$, and as in the pure states case, $`\mathrm{\Gamma }`$ is a curve that goes from $`0`$ to $`+\mathrm{}`$ of the real axis under all the poles of the lower half plane. Namely, in weak sense, the Liouville operator reads:
$`L=\{{\displaystyle \underset{n}{}}[(z_n\omega _0)\rho (\overline{z_{n,}0})\rho ^{}(\stackrel{~}{z_{n,}0})+(\omega _0z_{n)}^{}\rho (\overline{0,z_n})\rho ^{}(0,z_n)]+`$
$`+{\displaystyle _\mathrm{\Gamma }}[(z\omega _0)\rho (\overline{z,0})\rho ^{}(\stackrel{~}{z,0})+(\omega _0z^{})\rho (\overline{0,z})\rho ^{}(\stackrel{~}{0,z})]𝑑z+`$
$$\underset{jl}{}\zeta _l^j\overline{\beta (\sigma _{j,}\zeta _l^j)}\stackrel{~}{\beta ^{}(\sigma _{j,}\zeta _l^j)}+_0^{\mathrm{}}d\sigma [\underset{l}{}\zeta _l\overline{\beta (\sigma ,\zeta _l})\beta ^{}(\stackrel{~}{\sigma ,\zeta _l})+_{2\sigma }^{2\sigma }\nu \beta (\overline{\sigma ,\nu })\beta ^{}(\stackrel{~}{\sigma .\nu })d\nu ]$$
(59)
As in the pure states case the bases:
$`\{\rho (0),\rho (\overline{0,z_n}),\rho (\overline{z_n,0}),\rho (\overline{0,z}),\rho (\overline{z,0}),\beta (\omega ),\overline{\beta (\sigma _{j,}\zeta _l^j)},\beta (\overline{\sigma ,\zeta _l}),\beta (\overline{\sigma ,z})\}`$ and
$`\{\rho (0),\rho (\stackrel{~}{0,z_n}),\rho (\stackrel{~}{z_n,0}),\rho (\stackrel{~}{0,z}),\rho (\stackrel{~}{z,0}),\beta (\omega ),\stackrel{~}{\beta (\sigma _j,\zeta _l^j),}\beta (\stackrel{~}{\sigma ,\zeta _l}),\beta (\stackrel{~}{\sigma ,z})\}`$
are a biorthonormal system under the inner product (46) <sup>*</sup><sup>*</sup>*From now on we will consider that the discrete index $`\sigma _j`$ is included in the continuous index $`\sigma ,`$ and also $`\zeta _l^j`$ is included in $`\zeta _l.`$ Nevertheless we will conserve the terms $`\sigma _j`$ in all the spectral decompositions.. Also, as in the pure case:
$`(\beta (\overline{\sigma ,\zeta _l}|\beta (\overline{\sigma ^{},\zeta _l^{}}))=\delta _{\sigma \sigma ^{}}\delta _{ll^{}}\epsilon _l`$
$$(\beta (\stackrel{~}{\sigma ,\zeta _l}|\beta (\stackrel{~}{\sigma ^{}},\zeta _l^{}))=\delta _{\sigma \sigma ^{}}\delta _{ll^{}}\epsilon _l$$
(60)
where $`\epsilon _l=0`$ if $`Im\zeta _l0`$ and $`\epsilon _l=1`$ if $`Im\zeta _l=0`$, so the states corresponding to complex poles are ghosts as before. The same thing happens with $`\rho (\overline{0,z_n}),`$ $`\rho (\overline{z_n,0}),`$ $`\rho (\stackrel{~}{0,z_n}),`$ and $`\rho (\stackrel{~}{z_n,0}).`$
We can generalize the definition of trace as:
$$Tr\rho =\underset{i}{}<i|\rho |i>=(\rho |\underset{i}{}|i><i|)=(\rho |\underset{i}{}\beta (i0))=(\rho |I)$$
(61)
where $`\{|i>\}`$ is any basis of $``$. Now using the inner product (46) it can be easily proved that all the trace of all the off diagonal terms vanish. Therefore the trace of all the ghosts vanish.
Finally equations similar to eqs. (3.1) of paper can be obtained and using this equation it can be proved that, if $`\zeta _l,\zeta _l^{},\zeta _l^{\prime \prime },\mathrm{}`$ are complex, then:
$$Tr\beta (\overline{\sigma ,\zeta _l})\beta (\overline{\sigma ^{},\zeta _l^{}})=0,\text{ }\beta (\overline{\sigma ,\zeta _l})\beta (\overline{\sigma ^{},\zeta _l^{}})\beta (\overline{\sigma ^{\prime \prime },\zeta _{l^{\prime \prime }}})\mathrm{}.=0$$
(62)
This equation follows also from eq. (119) and says that the trace of the product of two ghosts and that the product of three or more ghosts vanish.
Since the Liouville space is a Hilbert space $`=()`$ we will have the Gel’fand triplets:
$$\mathrm{\Phi }_{}_{}\mathrm{\Phi }_{}^\times $$
(63)
$$\mathrm{\Phi }_+_+\mathrm{\Phi }_+^\times $$
(64)
Let us observe that, in order to satisfy eq. (41<sub>5</sub>) it is sufficient that the regular part of $`\rho 𝒮_{}𝒮_+`$ (since really in the second factor of eq.(41) there is a bra not a ket). Thus, as we will see in more detail in section 12 (cf. eq. (114))$`𝒮(`$ $`\varphi _{}\varphi _{})\mathrm{\Phi }_{}`$ and $`𝒮(\varphi _+\varphi _+)\mathrm{\Phi }_+.`$
As we will see $`\mathrm{\Phi }_{}`$ will be the space of physically admissible states, precisely the space of states such that they evolve with a non-decreasing of entropy accordingly to the Second Law of Thermodynamics. $`\mathrm{\Xi }(\mathrm{\Xi }\mathrm{\Xi })\backslash \mathrm{\Phi }_{}`$ is the set of physically non-admissible states. $`\mathrm{\Phi }_+`$ is the space of states such that they evolve with a non-growing entropy and, therefore they are clearly non physical. Macroscopically the physically admissible evolution are those that appear in nature, namely those that begin in an unstable state and go towards equilibrium (Gibbs ink drop spreading in the glass of water, a sugar lump solving in a cup of coffee, etc.). We will consider that everything the same in the microscopical case, namely that $`\mathrm{\Phi }_{}`$ is the space of physically realizable states. The physically non admissible evolutions of space $`\mathrm{\Phi }_+`$ can be obtained by the time inversion of the admissible ones, therefore they begin in an equilibrium state and evolve towards an unstable state (the ink or the sugar concentrating spontaneously and creating the drop or the lump). This kind of evolutions does not appear in nature, because the spontaneous appearance of an unstable state by a fluctuation, even if no completely impossible (remember we are developing quantum mechanics, an essentially statistical theory ), is highly improbable.
For all these reasons we will consider $`\mathrm{\Phi }_{}`$ the space of physical states.
## V Linear operators and observables.
We will now consider the linear operators $`A,`$ which are (anti)-linear functional over $`\mathrm{\Phi }_{}`$ and therefore belong to $`\mathrm{\Phi }_{}^\times ,`$ e. g. a derivative operator belongs to this space.
But these linear operators are merely theoretical or mathematical ”observables”. Real physical observables are less subtle, e. g.: there are not real apparatuses to measure mathematical derivatives. Real physical devices only measure ratios of small but finite quantities. Therefore we can consider that real physical observables live in an space endowed at least with the properties of $`𝒮(𝒮𝒮).`$ As in the case of the state it is also useful that these observable would have some analytic properties. There are two natural subspaces of $`𝒮(𝒮𝒮)`$ with definite analytic properties $`\mathrm{\Phi }_{}`$ and $`\mathrm{\Phi }_+`$. As we will see, in order to reproduce the relation between the Schroedinger and the Heisenberg pictures in the new theory, we must choose $`\mathrm{\Phi }_+`$ as the space of regular physical observables. As
$$\mathrm{\Phi }_{},\mathrm{\Phi }_+$$
(65)
the products between vectors of these two spaces are well defined. Then the mean values of all the observables of $`\mathrm{\Phi }_+`$ in the states of $`\mathrm{\Phi }_{}`$ are well defined. This property is sufficient to develop quantum mechanics formalism.
As we will see in a scattering theory, physical states are related to the preparation, they propagate towards the future in the Schroedinger picture, and they are well represented by states of space $`\mathrm{\Phi }_{\text{}}`$while physical observables are related with measurements, they propagate towards the past in the Heisenberg picture, and they are well represented by states of space $`\mathrm{\Phi }_+`$ .
After all these considerations the mean value of observable $`A\mathrm{\Phi }_+`$ in the states $`\rho \mathrm{\Phi }_{}`$ is:
$$<A>_\rho =A[\rho ]=(\rho |A)$$
(66)
where $`A[\rho ]`$ is an (anti) linear functional and:
$`A=A_0\rho (0)+{\displaystyle \underset{n}{}}[A_{n0}\rho (\stackrel{~}{z_n,0})+A_{0n}\rho (\stackrel{~}{0,z_n})+{\displaystyle _\mathrm{\Gamma }}[A_{z0}\rho (\stackrel{~}{z,0})+A_{0z}\rho (\stackrel{~}{0,z})]dz+`$
$$\underset{jl}{}A_{jl}\stackrel{~}{\beta (\sigma _j,\zeta _l^j)}+_0^{\mathrm{}}𝑑\sigma [A_\sigma \beta (\sigma )+\underset{l}{}A_{\sigma l}\beta (\stackrel{~}{\sigma ,\zeta _l})+_C\rho _{\sigma z}\beta (\stackrel{~}{\sigma ,z})𝑑z]$$
(67)
where the coefficients $`A`$ must satisfy eqs. (53) and (54). To fulfill these conditions, and eqs. (49) and (52) for the coefficients of $`\rho ,`$ it is sufficient that the analytic continuation of eq. (66) would be possible, one variable in the lower half plane and the other in the upper half plane, as we have done in ref. .
From the Gel’fand-Maurin theorem we know that we can diagonalize eq. (67) as:
$$A=\underset{i}{}a_i|a_i><a_i|$$
(68)
where in general $`a_i,`$ $`i`$ is an index such that the whole spectrum of $`A`$ is covered by the sum, and $`|a_i>`$ belongs to some specific rigged Hilbert space. If $`A`$ is selfadjoint then obviously $`a_i`$ and usually $`|a_i>𝒮^\times `$. In particular we can expand the energy operator as:
$$H=\underset{i}{}h_i|h_i><h_i|$$
(69)
where $`h_i`$ and $`|h_i>𝒮^\times `$, that can simply be obtained from eq. (9) or eq. (45) as:
$$H=_0^{\mathrm{}}E|E><E|𝑑E=E_0\beta (0)+_0^{\mathrm{}}𝑑\sigma E_\sigma \beta (\sigma )$$
(70)
Thus $`H`$ has two different spectral expansions, both very useful: one as observable eq. (69) and another as a evolution operator eq. (34), namely, in the short hand notation of section 3:
$$H=\underset{i}{}z_i|\overline{i}><\stackrel{~}{i}|$$
(71)
where $`z_i.`$ The difference between the two expansions comes from the fact that really they are weak equation corresponding to:
$$<\psi _1|H|\psi _2>=\underset{n}{}h_n<\psi _1|h_n><h_n|\psi _2>$$
(72)
where $`|\psi _1>,|\psi _2>𝒮`$,
$$<\phi |H|\psi >=\underset{i}{}z_i<\phi |\overline{i}><\stackrel{~}{i}|\psi >$$
(73)
where $`|\psi >\varphi _{}`$ and $`|\phi >\varphi _+.`$
By now we have all the mathematical objects to formulate our axiomatic theory.
## VI Axioms.
We will follow the main lines of ref. . So we postulate:
Axiom 1. To each dynamical variable $``$ (physical concept) there corresponds a linear operator $`R`$ $`\mathrm{\Phi }_{+\text{ }}`$ <sub>+</sub>(mathematical object) and the possible values of the dynamical variable are the eigenvalues of the operator.
Axiom 2. To each physical state there corresponds a unique state operator $`\rho \mathrm{\Phi }_{}_{}.`$ The average value of the dynamical variable $`,`$ $`(`$e.g.: of position, momentum, energy, etc.) represented by the operator $`R`$, in the virtual ensemble of events that may result from a preparation procedure for the state, represented by the operator $`\rho ,`$ is:
$$<>_\rho =\frac{Tr[\rho R]}{Tr\rho }$$
(74)
From these axioms, if we postulate the invariance of the theory under Galilei transformation, the explicit expression of the operators $`R`$ can be found and Schroedinger and Liouville equations can be deduced, as in book or in paper . Moreover, Planck’s constant $`\mathrm{}`$ appears as a proportionality coefficient between the geometrical generators and the physical magnitudes. In fact, these deductions can be implemented since we have just restricted the domain of definition of the states and of the observables but all the relevant demonstrations of the quoted reference remain valid. We do not reproduce this demonstration here because it is not in our main line of reasoning. So, in order to avoid these demonstrations, even if we maintain the Galilei invariance, we precise the main features of the time evolution by the following axiom:
Axiom 3. The time evolution of a state $`\rho (t)\mathrm{\Phi }_{}`$ is:
$$\rho (t)=e^{iHt}\rho (0)e^{iHt}=e^{iLt}\rho (0)$$
(75)
where $`\rho (0),\rho (t)\mathrm{\Phi }_{}`$ and $`H`$ is the hamiltonian operator of the system.
In this equation we must use eq. (71) if we want to expand $`H.`$ The exponents $`iHt`$ really read $`i\mathrm{}^1Ht,`$ so it is this axiom the one that introduces the universal constant $`\mathrm{}.`$ Of course we make $`\mathrm{}=1`$ below. From this axiom we also can demonstrate that $`\rho (t)`$ satisfies Liouville equation (39).
Since we have restricted the spaces of definition of the observables and the states to two spaces which are contained in $``$ nothing unphysical can really happen. Furthermore we obviously retain the main result of the usual quantum physics, as we will see in the next two sections, but with the new axiomatic structure we will gain new results that are, in fact, confirmed by experimental evidences.
## VII First consequences of the axioms.
We can now obtain the first consequences of the axioms following ref. :
a.- As $`\rho `$ is both in the numerator and in the denominator of eq. (74) we can normalize the state as:
$$Tr\rho =1$$
(76)
b.- If we postulate that projectors like $`P=|u><u|`$ are observables, and since their eigenvalues are 0 and 1, we can see that $`<P>_\rho =<u|\rho |u>`$ is real and positive so (cf. theorem 1 ref. ):
$$\rho =\rho ^{},\text{ }<u|\rho |u>0$$
(77)
c.- As $`\rho =\rho ^{},`$ $`\rho `$ can be expanded as:
$$\rho =\underset{i}{}\rho _i|\rho _i><\rho _i|$$
(78)
where $`\rho _i`$ and $`|\rho _i>`$ belongs to some adequate rigged Hilbert space. Then from the three first equations of this section we can obtain that:
$$\underset{i}{}\rho _i=1,\text{ }\rho _i=\rho _i^{},\text{ }0\rho _i1$$
(79)
d.- If we postulate that the mean value of any dynamical variable must be real, i. e.:
$$<>_\rho $$
(80)
then, for any pure state $`\rho =|v><v|,`$ $`|v>\varphi _{}`$ , we have:
$$<>_\rho =<v|R|v>$$
(81)
so, according to theorem 1 of ref. it is:
$$R=R^{}$$
(82)
so we can expand $`R`$ as in eq. (47), namely:
$$R=\underset{n}{}r_n|r_n><r_n|$$
(83)
where $`r_n`$ and $`|r_n>`$ belongs to some adequate rigged Hilbert space.
## VIII Probabilities.
From axiom 1 we have that:
$$<>_\rho =\underset{n}{}r_np_n(\rho )$$
(84)
where $`p_n(\rho )`$ is the probability to obtain the measurement $`r_n,`$ when we measure the dynamical variable $``$ in the quantum state $`\rho `$. From axiom 2 we also have:
$$<>_\rho =R[\rho ]=Tr[(\underset{n}{}r_n|r_n><r_n|)\rho ]=\underset{n}{}r_n<r_n|\rho |r_n>$$
(85)
In order that the last two equations would be equal it is sufficient that:
$$p_n(\rho )=<r_n|\rho |r_n>$$
(86)
It can be proved that this condition is also necessary if we repeat the corresponding demonstration of ref. .
Then, for every estate $`\rho `$ and every complete set of commuting observable $`\{R^{(\alpha )}\},`$ we can compute the probability to obtain the measurement $`r_n^{(\alpha )}`$ for the observable $`R^{(\alpha )}`$. In fact we can expand the observable as:
$$R^{(\alpha )}=\underset{n}{}r_n^{(\alpha )}|r_n^{(\alpha )}><r_n^{(\alpha )}|$$
(87)
and the probability is:
$$p_n^{(\alpha )}(\rho )=<r_n^{(\alpha )}|\rho |r_n^{(\alpha )}>$$
(88)
So we can see that $`\rho `$ really defines the quantum state of the system since, knowing $`\rho ,`$ we can obtain the probability of any measurement for any observable of the complete set of commuting observables. This is, in fact, the maximal information that we can obtain from a quantum state $`\rho `$ and, in consequence , this information also defines the quantum state of the system.
## IX Time-asymmetry and irreversibility.
In the last two sections we have briefly reviewed some results of ordinary quantum mechanics that turn out to be also valid in the new theory. It would be quite boring to continue this road reobtaining well known results so we will now consider the new features.
We will say that:
Time-asymmetry is the property of some single objects that turn out to be asymmetric, under the action of time-inversion Wigner operator $`K,`$ e. g. non real states $`|\psi >,`$ defined as the states such that $`K|\psi >|\psi >`$. In our case these objects are always statistical objects from the ensemble we are considering, since we are developing a statistical theory. Therefore the time-asymmetry of particular evolution of the members of the ensemble will be never taken into account.
Non time-reversal invariance is the property of some set of objects which are not invariant under $`K`$, e.g. the space $`\varphi _{}`$ which has the property (14).
Irreversibility is the property of some physical time evolutions such that its time inverted evolution turns out to be non-physical, namely it is physically forbidden , .
If we would further precise the term, the just introduced irreversibility would be the dynamical irreversibility (as we will see in a moment this irreversibility stems directly from the axioms). Thermodynamical irreversibility will be defined as the growing of entropy in section 13 (and we will see that more elements must be added to define this notion).
From the just quoted eq. (14) and the definitions at the beginning of section 4 we have that:
$$𝒦:\mathrm{\Phi }_{}\mathrm{\Phi }_+\mathrm{\Phi }_{}$$
(89)
where $`𝒦\rho =K\rho K^{}`$, and we can see that the physically admissible quantum states space of the theory is not time-reversal invariant.
From ref. , eq. (4.2) and eq. (4.3), we know that:
$$e^{iHt}:\varphi _{}\varphi _,\text{ }if\text{ }t>0,\text{ }e^{iHt}:\varphi _+\varphi _+,\text{ }if\text{ }t<0$$
(90)
Therefore, using the same demonstration regarding now the analytic properties of the functions of variable $`\nu ,`$ it can be proved that:
$$e^{iLt}:\mathrm{\Phi }_{}\mathrm{\Phi }_,\text{ }if\text{ }t>0,\text{ }e^{iLt}:\mathrm{\Phi }_+\mathrm{\Phi }_+,\text{ }if\text{ }t<0$$
(91)
so axiom 3 states that if $`\rho (t)\mathrm{\Phi }_{}`$ its evolution is only defined for $`t>0`$, and therefore the evolution operator $`e^{iLt\text{ }}`$cannot be physically inverted since its mathematical inverted operator $`e^{iLt}`$ corresponds to $`t<0`$ and therefore it is not well defined within space $`\mathrm{\Phi }_{}.`$ Namely the inverted evolution is forbidden by axiom 3. Therefore we have found that the new theory contains dynamical irreversible evolutions.
Of course $`t=0`$ is an arbitrary time so the condition $`t>0`$ physically simply means that operators $`e^{iHt\text{ }}`$and $`e^{iLt}`$ are not well defined for $`t\mathrm{}`$ for the state of space $`\mathrm{\Phi }_{}`$. Analogously, the condition $`t<0`$ means that the sane operators are not well defined for $`t+\mathrm{}`$ for the states of $`\mathrm{\Phi }_+.`$
## X Schrödinger and Heisenberg pictures and scattering experiments.
In the Schroedinger picture, if $`\rho (t)`$ is the time-variable state of the system and $``$ is a fixed dynamical variable, from axioms 2 and 3 we have:
$$<>_\rho (t)=Tr[\rho (t)R]=Tr[e^{iLt}\rho (0)R]=Tr[e^{iHt}\rho (0)e^{iHt}R]$$
(92)
According to axiom 3 $`\rho (t)\mathrm{\Phi }_{}`$ so, from eq. (91$`{}_{1}{}^{}),`$ we know that the last equation is only valid if $`t>0`$. Now from the cyclic property of the trace we also have that:
$$<>_\rho (t)=Tr[\rho (0)e^{iHt}R\text{ }e^{iHt}]$$
(93)
So we can define a time-variable Heisenberg operator:
$$R_H(t)=e^{iH(t)}R\text{ }e^{iH(t)}=e^{iL(t)}R$$
(94)
Then we have the Heisenberg picture equation:
$$<>_\rho (t)=Tr[\rho (0)R_H(t)]$$
(95)
But from eq. (91<sub>2</sub>) and since $`t<0`$ we know that the last time equation is only valid if $`R\mathrm{\Phi }_+.`$ This fact justifies both the choice of the operators space done in axiom 1, and what we have said in section 4.
In other words: in a scattering experiment (, ) the states are prepared at a time $`t_1`$ and propagate towards the future and therefore to times $`t>t_1`$ so according to eq. (91<sub>1</sub>) $`\rho \mathrm{\Phi }_{}.`$ At time $`t_2>t_1`$ dynamical variables $``$ are measured i. e. the S-matrix and the corresponding probabilities are obtained. But we can invert the procedure and propagate the dynamical variables $``$ and the corresponding operators $`R`$ towards the past down to time $`t_1<t_2.`$ Then we must propagate $`R`$ towards the past, therefore according to eq. (91$`{}_{2}{}^{})`$ $`R\mathrm{\Phi }_+.`$ So now we see the motivation of the choice of the spaces for $`\rho `$ and $`R`$ made in axioms 1 and 2, namely $`\rho \mathrm{\Phi }_{}`$ and $`R\mathrm{\Phi }_+.`$
## XI Equilibrium and decoherence.
We will face the problem of equilibrium in four steps: in the first one we will obtain a strong limit, in the second one a weak limit, in the third one the dominant time evolution components, and in the fourth one we will combine the two last ones to obtain some physical conclusions.
a.-From eqs. (57) and (59) we can deduce that if $`\rho (t)\mathrm{\Phi }_{}:`$
$`\rho (t)=\rho _0\rho (0)+{\displaystyle \underset{n}{}}[\rho _{n0}e^{i(z_n\omega _0)t}\rho (\overline{z_{n,}0})+\rho _{0n}e^{i(\omega _0z_n^{})t}\rho (\overline{0,z_n})+`$
$`+{\displaystyle _\mathrm{\Gamma }}[\rho _{z0}e^{i(z\omega _0)t}\rho (\overline{z,0})+\rho _{0z}e^{i(\omega _0z)t}\rho (\overline{0,z})]𝑑z+{\displaystyle \underset{jl}{}}\rho _{jl}e^{i\zeta _l^jt}\overline{\beta (\sigma _j,\zeta _l^j)}`$
$$_0^{\mathrm{}}𝑑\sigma [\rho _\sigma \beta (\sigma )+\underset{l}{}\rho _{\sigma l}e^{i\zeta _lt}\beta (\overline{\sigma ,\zeta _l})+_C\rho _{\sigma z}e^{izt}\beta (\overline{\sigma z})𝑑z]$$
(96)
where $`\zeta _l`$ and $`z_n`$ symbolize the complex poles. If we call, as it is traditional:
$`z_n=\overline{\omega }_n{\displaystyle \frac{i}{2}}\gamma _n,\text{ }\gamma _n>0`$
$$\zeta _l=\overline{\nu }_li\mathrm{\Gamma }_l,\text{ }\mathrm{\Gamma }_l>0$$
(97)
we have that:
$`\rho (t)=\rho _0\rho (0)+{\displaystyle \underset{n}{}}[\rho _{n0}e^{i(\overline{\omega }_n\omega _0)t}e^{\frac{1}{2}\gamma _nt}\rho (\overline{z_{n,}0})+\rho _{0n}e^{i(\omega _0\overline{\omega }_n)t}e^{\frac{1}{2}\gamma _nt}\rho (\overline{0,z_n})+`$
$`+{\displaystyle _\mathrm{\Gamma }}[\rho _{z0}e^{i(z\omega _0)t}\rho (\overline{z,0})+\rho _{0z}e^{i(\omega _0z)t}\rho (\overline{0,z})]𝑑z+{\displaystyle \underset{jl}{}}\rho _{jl}e^{i\overline{\nu }_l^jt}e^{\mathrm{\Gamma }_l^jt}\overline{\beta (\sigma _j,\zeta _l^j)}+`$
$$_0^{\mathrm{}}𝑑\sigma [\rho _\sigma \beta (\sigma )+\underset{l}{}\rho _{\sigma l}e^{i\overline{\nu }_lt}e^{\mathrm{\Gamma }_lt}\beta (\overline{\sigma ,\zeta _l})+_C\rho _{\sigma z}e^{izt}\beta (\overline{\sigma z})𝑑z]$$
(98)
For complex poles it is $`\gamma _n,\mathrm{\Gamma }_l>0,`$ so terms containing these positive gammas vanish when $`t+\mathrm{}.`$ The terms corresponding the continuous spectra do not vanish since the curves $`\mathrm{\Gamma }`$ and $`C`$ can be taken to be contained in the real axis almost anywhere. Then we obtain the strong limit:
$`\rho (t)\rho _{}(t)=\rho _0\rho (0)+{\displaystyle _\mathrm{\Gamma }}[\rho _{z0}e^{i(z\omega _0)t}\rho (\overline{z,0})+\rho _{0z}e^{i(\omega _0z)t}\rho (\overline{0,z})]𝑑z+`$
$$_0^{\mathrm{}}𝑑\sigma [\rho _\sigma \beta (\sigma )+_C\rho _{\sigma z}e^{izt}\beta (\overline{\sigma z})𝑑z]$$
(99)
so any state goes to a state of ”dynamical equilibrium”. We use the adjective ”dynamical” since it is a final state that it is a function of time. If we take into account normalization (76) we have:
$$Tr\rho (t)=Tr\rho _{}(t)=1$$
(100)
so:
$$\rho _0+_0^{\mathrm{}}\rho _\sigma 𝑑\sigma =1,$$
(101)
This is certainly a equation that the $`\rho _0`$ and $`\rho _\sigma `$ must satisfy, but in principle this is the only condition.
So in general the dynamical equilibrium state is not unique and depends of the initial conditions. This would be the general case.
b.-Moreover, if we consider that really the $`\rho `$ are functionals over the observables $`A`$ , since only the mean values $`<A>_\rho =A[\rho ]`$ are physically observed, and we use the Riemann-Lebesgue theorem, as in ref. , we obtain in a weak sense that:
$$\rho (t)\rho _{}=\rho _0\rho (0)+_0^{\mathrm{}}\rho _\sigma \rho (\sigma )𝑑\sigma $$
(102)
Therefore only the terms on the diagonal remain and we obtain a stationary final equilibrium as a weak limit. Only the state $`\rho (0),`$ corresponding to the ground state of the discrete spectrum (e. g. an oscillator), and the states $`\rho (\sigma )`$, corresponding to the diagonal states of the continuous one (e. g. the bath), remain in equilibrium. Thus we have proved that quantum decoherence takes place in our theory.
c.- But, using this method, we just obtain the limit, but we cannot see how this limit is attained. So we will use another approach. We know that the dominant component of the evolution towards equilibrium is given by the pole terms . In fact, this component is an excellent approximation for intermediate times: not too short times, so the Zeno effect would not be important, not too long times, so the Khalfin effect would not be important. Furthermore, experimentally we know that this intermediate period turns out to be very long, since, up to now, the Khalfin effect is not detected. So if we want to have a very good approximation of the evolution towards equilibrium we can neglect the regular background fields terms of curves $`\mathrm{\Gamma }`$ and $`C`$ and only consider the poles terms and the singular diagonal terms, precisely:
$`\rho (t)=\rho _0\rho (0)+{\displaystyle \underset{n}{}}[\rho _{n0}e^{i(\overline{\omega }_n\omega _0)t}e^{\frac{1}{2}\gamma _nt}\rho (\overline{z_{n,}0})+\rho _{0n}e^{i(\omega _0\overline{\omega }_n)t}e^{\frac{1}{2}\gamma _nt}\rho (\overline{0,z_n})`$
$$\underset{jl}{}\rho _{jl}e^{i\overline{\nu }_l^jt}e^{\mathrm{\Gamma }_l^jt}\overline{\beta (\sigma _j,\zeta _l^j)}+_0^{\mathrm{}}𝑑\sigma [\rho _\sigma \beta (\sigma )+\underset{l}{}\rho _{\sigma l}e^{i\overline{\nu }_lt}e^{\mathrm{\Gamma }_lt}\beta (\overline{\sigma ,\zeta _l})]$$
(103)
We will call:
$$\rho _{}=\rho _0\rho (0)+_0^{\mathrm{}}\rho _\sigma \rho (\sigma )𝑑\sigma $$
(104)
$`e^{\gamma _{\mathrm{min}}t}\rho _1(t)={\displaystyle \underset{n}{}}[\rho _{n0}e^{i(\overline{\omega }_n\omega _0)t}e^{\frac{i}{2}\gamma _nt}\rho (\overline{z_{n,}0})+\rho _{0n}e^{i(\omega _0\overline{\omega }_n)t}e^{\frac{i}{2}\gamma _nt}\rho (\overline{0,z_n})+`$
$$\underset{jl}{}\rho _{jl}e^{i\overline{\nu }_l^jt}e^{\mathrm{\Gamma }_l^jt}\overline{\beta (\sigma _j,\zeta _l^j)}+_0^{\mathrm{}}d\sigma \underset{l}{}\rho _{\sigma l}e^{i\overline{\nu }_lt}e^{i\mathrm{\Gamma }_l}\beta (\overline{\sigma ,\zeta _l})]$$
(105)
where we have made explicit the minimum of the $`\gamma `$ and the $`\mathrm{\Gamma }`$ in the l.h.s. of the second equation. Since $`\gamma ,\mathrm{\Gamma }>0`$ , when $`t\mathrm{}`$ we have:
$$\rho (t)\rho _{}$$
(106)
So, in this case, we obtain the usual stationary equilibrium which is not time dependent (but it still depends on the initial condition, through the $`\rho _0`$ and the $`\rho _\sigma ,`$ we will find an equilibrium independent of these conditions in section 14). The normalization condition is still eq. (101) and we have:
$$Tr\rho _1(t)=0$$
(107)
which is also a consequence of the fact that the ghost has vanishing trace.
So again we have obtained the usual equilibrium state and, as the off-diagonal terms vanish when $`t\mathrm{},`$ the phenomenon of decoherence is also proved.
The present method has been used to study decoherence in the cosmological case in papers .
## XII Conservation of the norm, the trace, and the energy. Lyapunov variables. Survival probability.
a.- In eq. (98) we have shown that some states of the theory vanish for $`t+\mathrm{}`$, precisely the ”ghost” states such that $`\gamma _n,\mathrm{\Gamma }_l>0.`$ Then we could wonder if the trace, the norm, and the energy are conserved in our theory. In fact, it is so, since we know that the trace of the off diagonal terms vanishes, so we have:
$$Tr\rho (t)=\rho _0+_0^{\mathrm{}}\rho _\sigma 𝑑\sigma =1$$
(108)
Also:
$$Tr[\rho (t)H]=\omega _0\rho _0+_0^{\mathrm{}}\sigma \rho _\sigma 𝑑\sigma =const$$
(109)
so the trace and the mean value of the energy are constant
In the pure state case these equations read:
$$<\psi |\psi >=1$$
(110)
$$<\psi |H|\psi >=const.$$
(111)
so the norm of pure states $`|\psi >`$ are also a constant.
b.- To continue let us repeat the reasonings of the beginning of this section in a different case : we will use another basis (precisely the one that can be obtained by the products of the basis of section 3) and the shorthand notation of section 3 and let it be $`N0`$. Let us consider the space $`𝒮(\varphi _{}\varphi _\text{ })`$ and a state $`\rho `$ $`𝒮(\varphi _{}\varphi _\text{ })`$ , that we can develop as:
$$\rho =\underset{i}{}\rho _i|i><i|+\underset{ij}{}\rho _{ij}|\overline{i}><\overline{j}|$$
(112)
where $`\rho _i𝒮,`$ $`\rho _{ij}𝒮_{}𝒮_+`$. Now as:
$$\rho _{ij}=\rho _{\sigma +\frac{1}{2}\nu ,\sigma \frac{1}{2}\nu }=\rho _{\rho \nu }$$
(113)
it turns out that $`\rho \mathrm{\Phi }_{},`$ since $`\rho _{\sigma \nu }`$ as a function of $`\nu `$ belongs to $`𝒮_{}^{(\sigma )},`$ then we can conclude that:
$$𝒮(\varphi _{}\varphi _\text{ })\mathrm{\Phi }_{}_{}\mathrm{\Phi }_{}^\times 𝒮^\times (\varphi _{}^\times \varphi _{}^\times )$$
(114)
So any function $`\beta \mathrm{\Phi }_{}^\times `$ (let say $`\rho \mathrm{\Phi }_{}`$ or $`\beta \overline{(\sigma ,\zeta _l)}\mathrm{\Phi }_{}^\times )`$ can also be expanded as in eq. (112) and then:
$$\beta (t)=e^{iHt}(\underset{i}{}\beta _i|i><i|+\underset{ij}{}\beta _{ij}|\overline{i}><\overline{j}|)e^{iHt}=\underset{i}{}\beta _i|i><i|+\underset{ij}{}\beta _{ij}e^{i(z_iz_j^{})t}|\overline{i}><\overline{j}|$$
(115)
Then as usual, if we call:
$$z_i=\omega _i\frac{i}{2}\gamma _i,\text{ }\gamma _i>0$$
(116)
we have:
$`\beta (t)={\displaystyle \underset{i}{}}\beta _i|i><i|+{\displaystyle \underset{ij}{}}\beta _{ij}e^{i(\omega _j\omega _i)t}e^{\frac{1}{2}(\gamma _i+\gamma _j)}|\overline{i}><\overline{j}|=`$
$$\underset{i}{}\beta _i|i><i|+\underset{IJ}{}\beta _{IJ}e^{i(\omega _J\omega _I)t}|\omega _I><\omega _J|+\underset{ijIJ}{}\beta _{ij}e^{i(\omega _j\omega _i)t}e^{\frac{1}{2}(\gamma _i+\gamma _j)t}|\overline{i}><\overline{j}|$$
(117)
where again the indices $`I,J,\mathrm{}`$ correspond to the real poles of the real continuous spectrum ($`ijIJ`$ means that either $`i`$ or $`j`$ or both correspond to ”ghost” states) and as $`\gamma _i+\gamma _j>0,`$ so if we make $`t\mathrm{}`$ we can see that we can expand $`\beta _{}(t)`$ and $`e^{\gamma _{\mathrm{min}}t}\beta _1(t)`$ as:
$$\beta _{}(t)=\underset{i}{}\beta _i|i><i|+\underset{IJ}{}\beta _{IJ}e^{i(\omega _J\omega _I)t}|\omega _I><\omega _J|$$
(118)
$$e^{\gamma _{\mathrm{min}}t}\beta _1(t)=\underset{ijIJ}{}\beta _{ij}e^{i(\omega _j\omega _i)t}e^{\frac{1}{2}(\gamma _i+\gamma _j)t}|\overline{i}><\overline{j}|$$
(119)
Finally on this basis the conservation of the norm and the energy read:
$$Tr\beta (t)=\underset{i}{}\beta _i=1$$
(120)
$$Tr(\beta (t)H)=\underset{i}{}\omega _i\beta _i=const.$$
(121)
so we have proved, in another basis, that the trace and the mean value of the energy are constant.
So up to now, every scalar we have introduced is time-constant and it seems impossible to define Lyapunov variables. But again let us try to find these results, using now another approach: considering only the regular component $`\rho _{reg}`$ of $`\rho `$ and using the restriction of the trace on the regular component, namely:
$$tr\rho _{reg}=\underset{i}{}<i|\rho _{reg}|i>=\underset{i}{}<\stackrel{~}{i}|\rho _{reg}|\overline{i}>$$
(122)
where the last equation can be obtained by analytic continuation. Then we have:
$$tr\rho _{reg}(t)=\underset{i}{}\rho _{ii}e^{\gamma _it}<\overline{i}|\overline{i}>=\underset{i}{}\rho _{ii}e^{i\gamma _it}\epsilon _i=\underset{I}{}\rho _{II}=const.$$
(123)
$$tr[\rho _{reg}(t)H]=\underset{i}{}\rho _{ii}(\overline{\omega }_i\frac{i}{2}\gamma _i)e^{\gamma _it}<\overline{i}|\overline{i}>=\underset{i}{}\rho _{ii}(\overline{\omega }_i\frac{i}{2}\gamma _i)e^{i\gamma _it}\epsilon _i=\underset{I}{}\omega _I\rho _{II}=const.$$
(124)
so again everything is constant and we cannot find Lyapunov variables.
From eq. (119) we can now prove eq. (62). In fact, in $`\rho _1(t)`$ or $`\beta _1(t)`$ there are only terms like $`|noghost><\overline{ghost}|,`$ $`|\overline{ghost}><noghost|,`$ and $`|\overline{ghost}><\overline{ghost}|.`$ If we compute $`\rho _1^2(t)`$ or $`\beta _1^2(t)`$ only the combination $`|\overline{ghost}><noghost|`$ $`noghost><\overline{ghost}|`$ survives, so these squares have only $`|\overline{ghost}><\overline{ghost}|`$ terms and their traces vanish. Similarly we can prove that the powers higher than two of $`\rho _1(t)`$ or $`\beta _1(t)`$ vanish.
c.- We have just learnt that the usual constants of quantum mechanics remain constant in the new theory so they are not Lyapunov variables, namely variable never-decreasing quantities. Nevertheless we can define Lyapunov variables in the new theory if we introduce the new ”metric” operators:
$$\stackrel{~}{M}=\underset{i}{}|\stackrel{~}{i}><\stackrel{~}{i}|,\text{ }\overline{M}=\underset{i}{}|\overline{i}><\overline{i}|$$
(125)
The role of these operators is to exchange bases $`\{|\stackrel{~}{i}>\}`$ and $`\{|\overline{i}>\}`$ since from eq. (35) we have:
$$\stackrel{~}{M}|\overline{i}>=\underset{j}{}|\stackrel{~}{j}><\stackrel{~}{j}|\overline{i}>=|\stackrel{~}{i}>$$
(126)
$$\overline{M}|\stackrel{~}{i}>=\underset{j}{}|\overline{j}><\overline{j}|\stackrel{~}{i}>=|\overline{i}>$$
(127)
this would happen, e. g. in the $`\beta (t)`$ of eq. (115). Therefore:
$$\stackrel{~}{M}:\varphi _{}\varphi _+,\text{ }\overline{M}:\varphi _+\varphi _{}$$
(128)
Then from eq. (35) we have:
$$\overline{M}\stackrel{~}{M}=\stackrel{~}{M}\overline{M}=1$$
(129)
$$<\stackrel{~}{i}|\overline{M}|\stackrel{~}{j}>=\delta _{ij}$$
(130)
$$<\overline{i}|\stackrel{~}{M}|\overline{j}>=\delta _{ij}$$
(131)
so the role of operators $`M`$ and $`\stackrel{~}{M}`$ is to eliminate the ”ghost” factor $`\epsilon _i`$ from eqs. (37), (38) and, as we can see, to make a variable what was previously a constant.
Then from eq. (98) we now have:
$$tr[\rho _{reg}(t)\stackrel{~}{M}]=\underset{i}{}\rho _{ii}e^{\gamma _it}<\overline{i}|\stackrel{~}{M}|\overline{i}>=\underset{I}{}\rho _{II}+\underset{iI}{}\rho _{ii}e^{\gamma _it}=$$
(132)
$`const.+{\displaystyle \underset{iI}{}}\rho _{ii}e^{\gamma _it}=Y(t)`$
where now $`Y(t)`$ is a Lyapunov variable, if e. g.: $`\rho _{ii}0`$ , and we have:
$$\stackrel{}{Y}(t)=\underset{iI}{}\rho _{ii}\gamma _ie^{\gamma _it}>0$$
(133)
because from eq.(97) $`\gamma _i>0`$ if $`iI.`$
If we want to find a Lyapunov variable without the condition $`\rho _{ii}0`$ we can introduce the ”linear entropy” :
$$tr[(\rho _{reg}(t)\stackrel{~}{M})^{}\rho _{reg}(t)\stackrel{~}{M}]=\underset{I}{}|\rho _{II}|^2+\underset{iI}{}|\rho _{ii}|^2e^{\gamma _it}=Y_G(t)$$
(134)
and in this case we surely have $`Y_G^{}(t)>0`$ since $`|\rho _{ii}|^2>0.`$
d.- We can also compute survival probabilities and found that in the theory we have non regular states with pure exponential survival probability. In fact let:
$$\rho (t)=|\overline{i(t)}><\overline{i(t)}|$$
(135)
$`|\overline{i(t)}>\varphi _+^\times .`$ At $`t=0`$ we can consider the observable $`R=|\psi ><\psi |\mathrm{\Phi }_+`$ and see how the probability $`p_i(t),`$ to measure the state $`\rho (t)`$ in the eigenstate $`|\psi ),`$ evolves: Using eq. (86) we found:
$$p_i(t)=<\psi |\overline{i(t)}><\overline{i(t)}|\psi >=e^{\gamma _it}<\psi |\overline{i(0)}><\overline{i(0)}|\psi >$$
(136)
which is an exponentially decaying survival probability with mean life $`\gamma _i^1.`$ This life-time turns out to be infinite for the non-ghost states $`|\omega _I>,`$ since in this case $`\gamma _I=0.`$ Then the non-ghost states are stable states and the ghost are unstable decaying states. The physical meaning of the exponents $`\gamma _i`$ is now clear, they are the inverse of the mean-life of the decaying processes within the theory.
In the case of the generic state of eq. (117) and a generic operator $`R=_\alpha \psi _\alpha |\psi _\alpha ><\psi _\alpha |`$ we have:
$`p_\rho (t)={\displaystyle \underset{ij\alpha }{}}\rho _{ij}\psi _\alpha e^{i(\omega _j\omega _i)}e^{\frac{1}{2}(\gamma _i+\gamma _j)t}<\psi _\alpha |\overline{i}><\overline{j}|\psi _\alpha >=`$
$`{\displaystyle \underset{IJ\alpha }{}}\rho _{IJ}\psi _\alpha <\psi _\alpha |I><J|\psi _\alpha >e^{i(\omega _J\omega _I)t}+`$
.
$$\underset{ijI}{}\rho _{ij}\psi _\alpha e^{i(\omega _j\omega _i)t}e^{\frac{1}{2}(\gamma _i+\gamma _j)t}<\psi _\alpha |\overline{i}><\overline{j}|\psi _\alpha >$$
(137)
The second term of the r.h.s. is clearly the dominant exponential component of the survival probability while the first one gives rise to the Zeno and Khalfin effects, as can be see in the examples .
## XIII Entropy.
As we have seen the irreversibility of the time evolution of the new theory allows us to introduce Lyapunov variables. In particular the equilibrium state $`\rho _{}(t)`$ defined in section 11 can be used to define a very important Lyapunov variable: a quantum conditional entropy . This entropy coincides with the usual conditional entropy in the classical limit, which has a remarkable property: it never decreases under a generic evolution driven by a Markov operator. It is, therefore, an excellent candidate for the phenomenological internal entropy. In this paper quantum conditional entropy is just the quantum analogue of classical conditional entropy, a extensive ever-growing Lyapunov variable that vanishes at equilibrium. A rigorous and complete study that relates this entropy with other kinds of entropy is missing nowadays (anyhow see , , ).
The naive definition of quantum conditional entropy of the state $`\rho `$ would be:
$$S=Tr[\rho \mathrm{log}(\rho \rho _{}^1)]$$
(138)
where $`\rho _{}`$ is a equilibrium state like the one of eq. (104). Of course, at equilibrium we have $`S=0`$.
In this section we will use the time evolution based in the dominant pole component of eq. (103) (so $`N=0)`$:
$$\rho (t)=\rho _0\rho (0)+\underset{n}{}[\rho _{n0}e^{i(\overline{\omega }_n\omega _0)t}e^{\frac{1}{2}\gamma _nt}\rho (\overline{z_{n,}0})+\rho _{0n}e^{i(\omega _0\overline{\omega }_n)t}e^{\frac{1}{2}\gamma _nt}\rho (\overline{0,z_n})+$$
(139)
$`{\displaystyle \underset{jl}{}}\rho _{jl}e^{i\overline{\nu }_l^jt}e^{\mathrm{\Gamma }_l^jt}\overline{\beta (\sigma _j,\zeta _l^j)}+{\displaystyle _0^{\mathrm{}}}𝑑\sigma [\rho _\sigma \beta (\sigma )+{\displaystyle \underset{l}{}}\rho _{\sigma l}e^{i\overline{\nu }_lt}e^{\mathrm{\Gamma }_lt}\beta (\overline{\sigma ,\zeta _l})]=\rho _{}+e^{\gamma _{\mathrm{min}}t}\rho _1(t)`$
where, in the second term of the r. h. s. we have made explicit the slowest dumping factor so $`\gamma _{\mathrm{min}}=\mathrm{min}(\gamma _n,\mathrm{\Gamma }_l)>0`$. All other factors contained in eq. (139) are oscillatory and they have a faster decrease than the slowest dumping factor. Since the trace of the off diagonal vanishes and using eq. (62), we know that:
$$Tr\rho _1(t)=0,\text{ }Tr\rho _1^2(t)=0,\text{ }\rho _1^n(t)=0,\text{ }if\text{ }n>2$$
(140)
We can also introduce a diagonal factor $`\rho _{}^1`$ among the factors $`\rho _1(t)`$ and the result will be the same, since being these matrices $`\rho _{}^1`$ diagonal they only modify the coefficient in the expansion of $`\rho _1(t)`$ while eq. (62) remains valid. Introducing eq. (139) in eq. (138) we have:
$$S(t)=Tr[(\rho _{}(t)+e^{\gamma _{\mathrm{min}}t}\rho _1(t))\mathrm{log}(1+e^{\gamma _{\mathrm{min}}t}\rho _1(t)\rho _{}^1(t))]$$
(141)
and expanding the logarithm we obtain:
$$S(t)=\frac{1}{2}e^{2\gamma _{\mathrm{min}}t}Tr[\rho _1^2(t)\rho _{}^1(t)]+\mathrm{}$$
(142)
where the dots symbolize terms with higher powers of $`\rho _1(t).`$ Then from eq. (140) we can conclude that $`S(t)=0`$, namely that this naive entropy is constant and it always coincides with equilibrium entropy. This result is analogous to those obtained in eq. (108) or eq. (109). So we can conclude, as in the pervious section, that if we do not use the metric operators $`\overline{M}`$ and $`\stackrel{~}{M}`$, or eventually some projector $`P,`$ we will always obtain constant naive quantities. Thus we must introduce, somehow, these operators in the naive definition (138). This is the new element that must be added to obtain a growing entropy as announced in section 9.
The most general and satisfactory solution is to define a projector:
$$P:\mathrm{\Phi }_{}\mathrm{\Phi }_P$$
(143)
and consider new projected density matrices $`\stackrel{~}{\rho }=P\rho `$ in such a way that a set of necessary properties should be fulfilled. These properties are:
1.-
$$Tr\stackrel{~}{\rho }(t)=Tr\stackrel{~}{\rho }_{}=1Tr\stackrel{~}{\rho }_1(t)=0$$
(144)
in order that the new projected matrices $`\stackrel{~}{\rho }`$ would have the same normalization properties as the old $`\rho .`$ In short: $`Tr\rho =Tr\stackrel{~}{\rho }.`$ In this sense $`P`$ is completely different than $`\overline{M}`$ or $`\stackrel{~}{M}`$ that transform a constant trace into a variable one.
2.-
$$Tr\stackrel{~}{\rho }_1^20,\text{ }\stackrel{~}{\rho }_1^n(t)0,n>2$$
(145)
and this must also be the case if some factors $`\stackrel{~}{\rho }_{}^1`$ would be included in the product. From these properties we would obtain the growing of the entropy.
3.-
$$\stackrel{~}{\rho }_{}=P\rho _{}=\rho _{}$$
(146)
in such a way that in the limit $`t\mathrm{}`$ the projection $`P`$ turns out to be irrelevant and we reobtain the usual Gibbs entropy for the equilibrium.
The most general $`P`$ acting on $`\mathrm{\Phi }_{}`$ would be:
$$P=\underset{ij}{}p_{ij}|\overline{i})(\stackrel{~}{j}|+\underset{ij}{}q_{ij}|\stackrel{~}{i})(\stackrel{~}{j|}$$
(147)
where, for the sake of simplicity we have, once more, trivialized the notation and we have written the basis $`\{\beta (\overline{..})\}`$ and $`\{\beta (\stackrel{~}{..})\}`$ with just one index as $`\{|\overline{i})\}`$ and $`\{|\stackrel{~}{i})\}`$. The most general matrix of $`\mathrm{\Phi }_{}`$ reads:
$$\rho =\underset{i}{}\rho _i|\overline{i})$$
(148)
namely eq. (45) in the new notation. If $`p_{ij},q_{ij}𝒮𝒮𝒮𝒮`$ then the behavior at infinities is good enough and then $`P\rho \mathrm{\Phi }_P.`$ Now if we symbolize the diagonal states by $`|I)`$ and the off diagonal ghost by $`|\overline{\mu })`$ we have that:
$$\rho =\underset{I}{}\rho _I|I)+\underset{\mu }{}\rho _\mu |\overline{\mu })$$
(149)
and
$`P={\displaystyle \underset{IJ}{}}p_{IJ}|I)(J|+{\displaystyle \underset{I\mu }{}}p_{I\mu }|I)(\stackrel{~}{\mu }|+{\displaystyle \underset{\mu I}{}}p_{\mu I}|\overline{\mu })(I|+`$
$$\underset{\mu \nu }{}p_{\mu \nu }|\overline{\mu })(\stackrel{~}{\nu }|+\underset{\mu I}{}q_{\mu I}|\stackrel{~}{\mu })(I|+\underset{\mu \nu }{}q_{\mu \nu }|\stackrel{~}{\mu })(\stackrel{~}{\nu }|$$
(150)
But in order to satisfy property 3 it must be:
$$p_{IJ}=\delta _{IJ},\text{ }p_{\mu I}=q_{\mu I}=0$$
(151)
Let us now compute:
$$P\rho =\underset{I}{}\rho _I|I)+\underset{I\mu }{}p_{I\mu }\rho _\mu |I)+\underset{\mu \nu }{}p_{\mu \nu }\rho _\nu |\overline{\mu })+\underset{\mu \nu }{}q_{\mu \nu }\rho _\nu |\stackrel{~}{\mu })$$
(152)
and let us compute the traces of $`\rho `$ and $`P\rho :`$
$$Tr\rho =\underset{I}{}\rho _I,\text{ }TrP\rho =\underset{I}{}\rho _I+\underset{I\mu }{}p_{I\mu }\rho _\mu $$
(153)
Accordingly to condition 1 these norms must be equal for every $`\rho _\mu `$ thus: $`p_{I\mu }=0.`$
Also, as $`P^2=P`$, it must be:
$$\text{ }\underset{\lambda }{}p_{\mu \lambda }p_{\lambda \nu }=p_{\mu \nu }\text{ }\underset{\lambda }{}q_{\mu \lambda }p_{\lambda \nu }=q_{\mu \nu }$$
(154)
Then these are the conditions that the coefficients of the projector must satisfy. These equations have solutions, since the first one is satisfied for any projector in the ”ghost” space and the second would be satisfied if $`p_{\mu \nu }=q_{\mu \nu }.`$ But, of course, more general solutions can be found. So the solution is not unique and we would have many possible projectors that will originate many possible entropies (out of equilibrium), as we will see. We also must remark that the condition $`P^2=P`$ is not strictly necessary, so we also can develop a theory where $`P`$ is not a projector. Finally:
$$P\rho _1=\underset{\mu \nu }{}p_{\mu \nu }\rho _\nu |\overline{\mu })+\underset{\mu \nu }{}q_{\mu \nu }\rho _\nu |\stackrel{~}{\mu })$$
(155)
Thus $`Tr\stackrel{~}{\rho }_1^2=Tr(P\rho _1)^2`$ $`0,`$ $`\stackrel{~}{\rho }_1^n(t)0,n>2`$ and these results will be valid if we intercalate $`\rho _{}^1`$ factors. In fact, the ”q” term in $`P`$ introduces the $`|\stackrel{~}{ghost}>`$ that does not vanish when multiplied by the $`|\overline{ghost}>.`$ Then condition 2 is also satisfied.
Now we can redefine our quantum conditional entropy as:
$$S=Tr[\stackrel{~}{\rho }\mathrm{log}(\stackrel{~}{\rho }\rho _{}^1)]$$
(156)
Repeating the calculations we obtain:
$$S(t)=\frac{1}{2}e^{2\gamma _{\mathrm{min}}t}Tr[(\stackrel{~}{\rho }_1(t))^2\rho _{}^1(t)]+\mathrm{}$$
(157)
where now the dots symbolize terms that vanish faster than the first one. But with this new definition $`S(t)0`$ and:
$$\underset{t\mathrm{}}{lim}S(t)=0$$
(158)
namely $`S(t)`$ goes to the equilibrium entropy when $`t\mathrm{}.`$ So, from eq. (157) we can say that near to equilibrium our quantum conditional entropy is negative and it always grows, but of course, far from the equilibrium, the evolution may have not these properties.
So we realize that the ever-growing property of the entropy is not a quantum property but just a classical one. Thus, for the sake of completeness we will sketch in the appendix the relation between the classical case and the quantum case. Then, using our equations we will see that the classical analogue of our quantum conditional entropy never decreases.
But, for each $`P`$ there is a different entropy, both in the quantum and the classical cases. Really this is not a problem since all these entropies coincide in the equilibrium limit, due to condition 3 (also near equilibrium all of them have the same dominant dumping factor). So we have one and only one equilibrium thermodynamic entropy, the Gibbs one, and many non equilibrium entropies, that depend on the choice of $`P,`$ namely the choice of the apparatus that measures these entropies. So our position is that, even if the arrow of time is intrinsic and defined by space $`\mathrm{\Phi }_{}`$ the definition of the out-of-equilibrium entropy is observer (or measurement apparatus) dependent. But all these entropies grow in the same time direction and therefore share the same arrow of time.
## XIV Thermalization
From eq. (24) we know that our method allows to split the hamiltonian in two non interacting parts a discrete one related with the oscillating and damped modes of the eventual of $`n`$ oscillators of our model and a continuous one , related with the field or the bath $`\omega .`$ In eq. (24) there is no interaction between the discrete and continuous modes so we may ask ourselves how a bath in thermal equilibrium can thermalize the oscillators in our model.
Obviously the thermalization can only be seen if we go to the basis where there is some interaction between the oscillators and the bath, namely the basis of the unperturbed hamiltonian $`H_0`$ that we can diagonalize as:
$$H_0=\underset{n}{}E_n^{(0)}|E_n^{(0)}><E_n^{(0)}|+_0^{\mathrm{}}E^{(0)}|E^{(0)}><E^{(0)}|dE^{(0)}$$
(159)
while the perturbed hamiltonian reads:
$$H=\underset{n}{}z_n|\overline{f_n}><\stackrel{~}{f_n}|+_0^{\mathrm{}}\omega |\overline{f_\omega }><\stackrel{~}{f_\omega }|d\omega $$
(160)
But, $`H=H_0+V`$ so, in the basis $`\{|E_n^{(0)},E^{(0)}>\},`$ we can see the interaction and the thermalization phenomena, but not in the basis of eq. (160)
In fact, any initial $`\rho ,`$ in a thermal bath, can be written as:
$$\rho =\underset{nm}{}\rho _{nm}|E_n^{(0)}><E_m^{(0)}|+Z_0^{\mathrm{}}e^{\beta E^{(0)}}|E^{(0)}><E^{(0)}|dE^{(0)}$$
(161)
(in this section we do not consider the off-diagonal components of the bath, since we consider that the bath is always in equilibrium, so these components take no part in the thermalization procedure). The unperturbed bath is a thermal state at temperature $`\beta ^1`$ and $`Z`$ is a normalization coefficient. Then, introducing:
$$I=\underset{n^{}}{}|\overline{f_n^{}}><\stackrel{~}{f_n^{}}|+_0^{\mathrm{}}|\overline{f_\omega }><\stackrel{~}{f_\omega }|d\omega $$
(162)
in each term and tracing away the bath, because we are only interested in the oscillator’s modes term $`\rho _{O\text{ ,}}`$we have.
$`\rho _O={\displaystyle \underset{nm}{}}\rho _{nm}{\displaystyle \underset{n^{}m^{}}{}}|\overline{f_n^{}}><\stackrel{~}{f_n^{}}|E_n^{(0)}><E_m^{(0)}|\overline{f_m^{}}><\stackrel{~}{f_m^{}|}+`$
$`Z{\displaystyle _0^{\mathrm{}}}e^{\beta E^{(0)}}{\displaystyle \underset{n^{}m^{}}{}}|\overline{f_n^{}}><\stackrel{~}{f_n^{}}|E^{(0)}><E^{(0)}|\overline{f_m^{}}><\stackrel{~}{f_m^{}|}dE^{(0)}=`$
$`{\displaystyle \underset{n^{}m^{}}{}}|\overline{f_n^{}}><\stackrel{~}{f_m^{}}|[{\displaystyle \underset{nm}{}}\rho _{nm}<\overline{f_n^{}}|E_n^{(0)}><E_m^{(0)}|\stackrel{~}{f_m^{}}>+`$
$$Z_0^{\mathrm{}}e^{\beta E^{(0)}}<\overline{f_n^{}}|E^{(0)}><E^{(0)}|\stackrel{~}{f_m^{}}>dE^{(0)}]$$
(163)
To see what really is going on let us go to the Friedrichs model of ref. and consider the case where $`V0`$ and $`t\mathrm{},`$ using the equations (5.20) of this paper, which translated to our notation reads:
$`|<\overline{f_n^{}}|E_n^{(0)}>|^20,\text{ }|<E_m^{(0)}|\stackrel{~}{f_m^{}}>|^20`$
$$|<\overline{f_n^{}}|E^{(0)}>|^2\delta (E^{(0)}E_n^{}^{(0)}),\text{ }|<E^{(0)}|\stackrel{~}{f_m^{}>}|^2\delta (E^{(0)}E_m^{}^{(0)})$$
(164)
where we have neglected the terms like $`|<\overline{f_n^{}}|E^{(0)}>|^2`$ since they vanish in the limit when $`V0`$ because they relate the discrete eigenkets with the continuous ones. Using these equations we obtain the result:
$$\rho _0\underset{n}{}|\overline{f_n}><\stackrel{~}{f_n}|e^{\beta E_n^{(0)}}\underset{n}{}|E_n^{(0)}><E_n^{(0)}|e^{\beta E_n^{(0)}}$$
(165)
so the oscillators will be thermalized when $`t\mathrm{}`$, for small interactions.
## XV Local and global time-asymmetry.
Let us now precise the meaning of two important words: ”conventional” and ”substantial”:
-In mathematics we use to work with identical objects, like points, the two directions of an axis, the two semicones of a light cone, the two time orientations of a time-oriented manifold, etc.
-In physics there are also identical objects, like identical particles, the two spin directions, the two minima of a typical ”two minima” potential, etc.
-When (, ) we are forced to call two identical objects by different names we will say that we are establishing a conventional difference, e.g.: when we call $`e_1`$ and $`e_2`$ two electrons, or ”up” and ”down” two spin directions, or ”right” and ”left” two minima of a symmetric potential curve, while
\- if we call two different objects by different names we will say that we are establishing a substantial difference.
The problem of time asymmetry is that, in all time-symmetric normal physical theories, usually the difference between past and future is just conventional. In fact, if we can change the word ”past” by the word ”future”, in these theories, the theory remains valid and nothing actually changes. But we have the clear psychological filling that the past is substantially different from the future. Thus the problem of the arrow of time is to find theories where the past is substantially different from the future, such that the usual well established physics would remain valid. As we will see our minimal irreversible quantum mechanics would be one of these theories.
But, up to now, we have just postulated that the states $`\rho \mathrm{\Phi }_{}`$ and the observables $`R\mathrm{\Phi }_+,`$ namely we have established a local time asymmetry within a particular quantum system . To establish a global time asymmetry for the whole universe is quite a different task.
It is clear that if we content ourselves with the local time asymmetry we must face several problems. E. g. as $`\mathrm{\Phi }_\text{ }`$and $`\mathrm{\Phi }_+`$ are only conventionally different there is no reason to justify that all the states would belong to a space $`\mathrm{\Phi }_{}`$ and all the operators to a space $`\mathrm{\Phi }_+`$ in all the systems of the universe. Most likely the states would belong to $`\mathrm{\Phi }_\text{ }`$ in the 50% of the system of the universe and to $`\mathrm{\Phi }_+`$ in the other 50%. For the operators we would have the inverted situation. Then: what would happen when a system of the first class interacts with a system of the second class? Perhaps somehow would an average arrow of time be established?
To scape from this dilemma it is clear that we need a cosmological model that correlates all the local arrows of time, and we know that, nowadays, cosmological models are never completely satisfactory. Nevertheless the more satisfactory, realistic, and simplest model to solve the problem is Reichenbach global system that we will further explain.
## XVI Reichenbach global system.
The set of irreversible processes within the universe, each one beginning in an unstable non-equilibrium state, can be considered a global system , . Namely, every one of these branching processes began in a non-equilibrium state, such that, this state was produced by a previous process of the set. E. g.: Gibbs ink drop (initial unstable state) spreading in a glass of water ( irreversible process) it only may exist if it was produced by an ink factory (since the probability to concentrate an ink drop by a fluctuation in a glass, where the ink is nixed with the water, is extremely small). This factory extract the necessary energy from an oven, where coal (initial unstable state) is burning (branch irreversible process); in turn coal was created with energy coming from the sun, where $`H`$ (initial unstable state) is burning (branch irreversible process); finally $`H`$ was created using energy obtained from the unstable initial state of the universe (the absolute initial state of the global system). It is now clear that within the global system all the arrows of time point to the same direction since all of them are originated in the same global initial unstable state (why the initial state is unstable is explained in papers and and it is not of our concern in this paper).
Of course Reichenbach global system was originally imagined as a model of the classical universe. But we can also imagine a quantum global system.
Let us draw the ordinary diagram of a scattering experiment (fig. 2) to have a graphic idea of the nature of the unstable states. In the center of the diagram there is a black box that symbolizes any scattering process. A set of stables ”in” states $`a_1,a_2,\mathrm{}`$ is transformed by the scattering process in another set of stable ”out” states $`b_1,b_2,\mathrm{}`$. It is a reversible process because the evolution equations are time-reversible, so we can interchange the ”in” and ”out” states and all the results remain valid. In fact, fig. 2 is essentially symmetric. The variation of the quantum entropy vanishes in this process .
Now, let us cut the black box in two parts by the dotted line draw at $`t=0`$. Then, we can consider the right side of the figure, namely fig. 3. This figure was introduced by A. Bohm , so we will call this kind of figures ”Bohm diagrams”. In fig. 3 the set of stable ”incoming” states creates a set of unstable states, $`u_1,u_2,\mathrm{}`$ which are growing states and they belong to space $`\varphi _+`$ . (e. g. radiation exciting an electron of the ground state). As the states of $`\varphi _+`$ are linear combinations of the states of $`\varphi _{}^\times `$, in some sense they can also be considered also as growing states and as such they can be symbolized as horizontal lines inside the half box. Fig. 3 is asymmetric and it symbolizes an irreversible creation process. The evolution equation is still time-symmetric but irreversibility is introduced by the growing nature of the states of space $`\varphi _{}^\times `$ or by the non-invertible time evolution operator acting on $`\varphi _+`$ for $`t<0`$ . The variation of the entropy is negative in this process. This is not contradictory since in every creation process there is an incoming energy and then we can consider the system as an open one.
We can also consider the second half of figure 1, namely fig. 4. It is the Bohm diagram of a decaying process where a set of unstable decaying states $`u_1,u_2,\mathrm{},`$ that are linear combinations of the basis of space $`\varphi _+^\times ,`$ is transformed in a set of stable ”outgoing” stable states of $`\varphi _{}`$ (e. g.: an excited electron decaying into the ground state and emitting radiation). Fig 4 is asymmetric and symbolizes a decaying irreversible process. Again, evolution equations are still time-symmetric but the decaying nature of the states of space $`\varphi _+^\times `$ introduces the irreversibility, etc., etc. The variation of the entropy is positive in this process. These would be the diagrams corresponding to local processes.
But Bohm diagrams allow us to also see also quantum structure of a global system ( fig. 5). The universe is represented by a set of branching scattering processes with one initial unstable state symbolized by the cut box (at ”big-bang” time $`t=0)`$ in the far right. Each subsystem going from an unstable state to equilibrium (the ink drop spreading in the water, the sugar lump solving in the coffee,…) is symbolized by a decaying process as the one of fig 4, namely the diagram in the shaded box of figure 5. The creation of an unstable state is symbolized by a creation process (like the one of fig. 3) where energy comes from a previous decaying process (the ink factory with its oven). One of these larger subsystem with its source of energy is represented in the dotted box in fig. 5 . The overall process is irreversible, because fig. 5 is asymmetric, and if we would have a model of this universe ( see some simple models in ref., ) the state of the universe must belong to some global space $`\varphi _{}^G`$ or $`\mathrm{\Phi }_{}^G.`$ Therefore in this diagram there is a clear arrow of time. But in the previous diagrams (fig. 3 or fig. 4) the arrow of time was a ”local” one, while this diagram has one of the most important characteristic of the observed time asymmetry: it is global. This is the way to introduce the arrow of time in a time-symmetric dynamical formalism: by a global and generalized symmetry breaking process
But we must remember that the difference between the global $`\mathrm{\Phi }_{}^G`$ and the global $`\mathrm{\Phi }_+^G`$ of the whole universe global system is just conventional since these two spaces are identical. Thus physics is the same in $`\mathrm{\Phi }_{}^G`$ than in $`\mathrm{\Phi }_+^G.`$ Think in a cosmological model, (fig. 5) life will be the same in this universe (with a quantum state in space $`\mathrm{\Phi }_{}^G)`$ than in the universe of fig. 6, the time inverted image of fig. 5, (with a t-inverted quantum state in space $`\mathrm{\Phi }_+^G).`$ In fact, since in both models of the universe (if completely computed) all the arrows of time must point to the same direction, there is no physical way to decide if we are in one model or in the other. So both models are identical. Thus the choice between $`\mathrm{\Phi }_{}^G`$ and $`\mathrm{\Phi }_+^G`$ is just conventional and physically irrelevant (as to choose one of the two minima of the potential in spontaneous symmetry breaking problem).
But once this choice is made a substantial difference is established in the model e. g. the only time evolution operator is $`U_{}(t)=e^{iHt},t>0,`$ and it cannot be inverted, we only have equilibrium towards the future, entropy growing, etc. .
Once the $`\mathrm{\Phi }_{}^G`$ or the $`\mathrm{\Phi }_+^G`$ is chosen in the global system we are forced to choose the corresponding spaces in the local subsystems, if we want to study these subsystems as isolated systems, and a global arrow of time is established. This fact solves the problem stated in the preceding section.
Thus the choice between $`\mathrm{\Phi }_{}^G`$ and $`\mathrm{\Phi }_+^G`$ is trivial and unimportant (but this choice must be a global one), that is why the arrow of time is not introduced by hand in our theory. The important choice is between $`^G`$ and $`\mathrm{\Phi }_{}^G`$ (or $`\mathrm{\Phi }_+^G)`$ as the space of our physical admissible states. And we are free to make this choice, since a good physical theory begins by the choice of the best mathematical structure to mimic real nature. Thus, our thesis is essentially that time asymmetric mathematical structures like ours mimic in the more economical way the time-asymmetric of the Nature where we live, than time-symmetric mathematical structures.
## XVII Other results.
The main results related with quantum mechanics are stated in the above sections. But we must comment that using the present formalism all the relevant results of irreversible statistical mechanics can also be obtained, e. g. all the results of the book , as it is proved in ref. , because the main $`\mathrm{\Pi }`$ projector of the quoted book can be defined using Gel’fand triplets. Also, in some simple cases, we can go from the quantum models to the classical ones , where we find the same philosophy, in classical cases. Chaotic models like Baker’s transformation and Renyi’s maps, are also treated with the same method, with good results , . Other interesting results are contained in papers , , , , and . So what we have explained is just the quantum axiomatic formalism of a general method to deal with irreversible processes.
## XVIII Conclusion.
We claim that our axiomatic formalism condenses the main ideas of how to solve the problem of time asymmetry pioneered by many authors that appears in the bibliographical references. This method is based in the reasonings we have made in the introduction and it yields correct physical results that coincide with those obtained by other methods (coarse-graining, traces, stochastic noises, lost of information, etc.).
We do not foresee any cross-experiment to settle which of the quoted methods (coarse-graining, traces, stochastic noises, lost of information, etc.) or ours is the correct one, because we think that they are somehow complementary and that they only explain the real physical world in different ways. Therefore we believe that we must solve at least three problems in order to complete our theory:
a.- The methods of coarse-graining: traces, projections, stochastic noises, lost of information, etc. have clear physical motivations. On the contrary our method is just based on the search of the minimal mathematical structure to explain time-irreversibility. Even if this reason can be a sufficient one for mathematical minded readers, it turns out to be not so eloquent for physical minded ones. So we must find the relation among all the method because most probably they are all based in the same or very similar (eloquent) physical bases. This work was already began in paper where it is shown that causality is the real reason for the choice of the proper space of physical states.
b.-We have admitted, under eq. (14), that there is not a unique way to define space $`\varphi _{}.`$ We must find the necessary and sufficient condition to define a unique space $`\varphi _{}`$ or at least a class of spaces all of them endowed with enough properties to explain all phenomena related with time-asymmetry. Lax-Phillips scattering theory, recently redeveloped by Boris Pavlov , and the quoted paper seem the more promising ways to solve the problem.
c.- Essentially space $`\mathrm{\Phi }_{}`$ is the one where all out-of-equilibrium entropies turn out to be growing. But we have said, at the beginning of section 13, that a complete study that single out, in an indisputable way, an entropy, that generalizes the phenomenological thermodynamical entropy at equilibrium, is missing. Most likely only when this task will be accomplished we will be able to motivate the choice of space $`\mathrm{\Phi }_{}`$ convincingly.
Meanwhile we just claim that our formalism is a contribution to the understanding of time-asymmetry.
## XIX Acknowledgment.
This work was partially supported by grants: CI1-CT94-0004 and PSS-0992 of the European Community, PID-0150 of CONICET (National Research Council of Argentina), EX-198 of the Buenos Aires University, and 12217/1 of Fundación Antorchas and also grants from the Fondation pour la Receherche Fondamentale OLAM and the British Council.
## XX Appendix . Wigner function integral and classical entropy.
Let $`\rho `$ be a density matrix of Liouville space $`=()`$ and let $`\{|q>\}`$ be the configuration or position basis of the Hilbert space $`.`$ The Wigner function corresponding to state $`\rho `$ is a real but not positive definite quantity that reads ( eq. (2.1)):
$$\rho _W(q,p)=\pi ^1<q\lambda |\rho |q+\lambda >e^{2i\lambda p}𝑑\lambda $$
(166)
This equation is valid if the wave function has only one variable. If it would have $`n`$ variables the $`\pi ^1`$ must be changed by $`\pi ^n.`$ It can be proved that:
$$L\rho _W(q,p)=\pi ^1<q\lambda |L\rho |q+\lambda >e^{2i\lambda p}𝑑\lambda +O(\mathrm{})$$
(167)
where $`L`$ is respectively the classical and quantum Liouville operator. In the classical limit $`\mathrm{}0`$ therefore $`\rho _W`$ can be considered as the classical distribution function corresponding to $`\rho `$. As in the classical regime we practically work in this limit we will consider that eq.(166) is the relation between the quantum density matrix and the classical distribution function. In fact, even if $`\rho _W`$ is not generally positive definite, using the Wigner integral from classical equations we can obtain quantum equations and vice versa, as we will see in a few examples. E.g., let us observe that ( eq. (2.6):
$$\rho _W=\rho _W(q,p)𝑑q𝑑p=$$
(168)
$`={\displaystyle 𝑑q<q\lambda |\rho |q+\lambda >\delta (\lambda )𝑑\lambda }=Tr\rho `$
so to the classical norm corresponds the quantum trace. Also, if we define the classical analogue of operator $`O`$ as ( eq. (2.12)):
$$O_W(q,p)=𝑑ze^{ipz}<q\frac{1}{2}z|O|q+\frac{1}{2}z>$$
(169)
we obtain that :
$$(\rho _W|O_W)=\rho _W(q,p)O_W(q,p)𝑑p𝑑q=Tr(\rho O)$$
(170)
Therefore to the inner product in classical Liouville space corresponds the inner product in the quantum Liouville space. Finally if $`O_1`$ and $`O_2`$ are two operators and $`O_P=O_1O_2`$ it is ( eq. (2.59)):
$$O_{PW}(q,p)=O_{1W}(q,p)e^{\frac{\mathrm{\Lambda }}{2i}}O_{2W}(q,p)$$
(171)
where:
$$\mathrm{\Lambda }=\stackrel{}{\frac{}{p}}\stackrel{}{\frac{}{q}}\stackrel{}{\frac{}{q}}\stackrel{}{\frac{}{p}}$$
(172)
Therefore if $`O_1=O_2`$ we have that:
$$O_{PW}(q,p)=[O_{1W}(q,p)]^2$$
(173)
This fact completes the analogy between classical and quantum spaces implemented by the Wigner integral.
As Wigner integral (166) is a linear mapping the quantum evolution equation (139) can be reproduced in the ”classical case”as
$$\rho _W(t)=\rho _W+e^{\gamma _{\mathrm{min}}t}\rho _{1W}(t)$$
(174)
where, somehow, $`\rho _W`$ could be considered as a classical distribution function. It is not so because it is not positive definite. But we will see that $`\rho _W`$ has this property near equilibrium, where we know that, due to the appearance of decoherence and correlations, the state is classic.
In fact, the classical states are defined as those having decoherence and correlations (this fact can be seem in ref. , using coarse-graining method or using our method, in section 11 of this paper for decoherence and in ref. for correlations), namely those that can be expanded as:
$$\rho =\underset{I}{}\rho _I|I>_{cor}<I|_{cor},\text{ }\rho _I0$$
(175)
where $`|I>_{cor}`$ is a no-ghost state such that its position $`<q|I>_{cor}`$ and its momentum $`<p|I>_{cor}`$ are correlated, i. e. they are as defined as possible (e. g. as the ground state of a H atom). Precisely, they are as defined as it is allowed by the uncertainty principle, around a point $`(q_I,p_I)`$. The quantum equilibrium state $`\rho _{}`$ is one of these states. Let $`\rho _{W\text{ }}`$ be the classical equilibrium states . The Wigner integral of this state gives:
$$\rho _W(q,p)=\pi ^1_{\mathrm{}}^+\mathrm{}<q\lambda |\rho _{}|q+\lambda >e^{2i\lambda p}𝑑\lambda =$$
(176)
$`=\pi ^1{\displaystyle \underset{I}{}}\rho _I{\displaystyle _{\mathrm{}}^+\mathrm{}}<q\lambda |I>_{cor}<I|_{cor}|q+\lambda >e^{2i\lambda p}𝑑\lambda =`$
$`\pi ^1{\displaystyle \underset{I}{}}\rho _I{\displaystyle _{\mathrm{}}^+\mathrm{}}<q_I\lambda |I>_{cor}<I|_{cor}|q_I+\lambda >\delta (\lambda )e^{2i\lambda p}𝑑\lambda =`$
$`=\pi ^1{\displaystyle \underset{I}{}}\rho _I{\displaystyle _{\mathrm{}}^+\mathrm{}}|<q_I|I>_{cor}|^2e^{2i\lambda p}\delta (\lambda )𝑑\lambda =\pi ^1{\displaystyle \underset{I}{}}\rho _I|<q_I|I>_{cor}|^20`$
since the functions $`<q_I|I>_{cor},`$ $`<I|_{cor}|q_I>`$ vanish except around a definite value of $`q_I,`$ so the states $`<q\pm \lambda |`$ are only different from zero when $`qq_I`$ and $`\lambda 0,(`$this is the reason of the factor $`\delta (\lambda ))`$ and because of eq. (79) which is valid for $`\rho _{}.`$ Then we conclude that the equilibrium states and the quantum states near $`(q_I,p_I)`$, have the property $`\rho _W(q,p)0`$ and therefore these Wigner functions can be considered classical states. These states disappear, if we wait long enough, due to the dumping term in eq. (174)
So, in this classical regime, a usual conditional entropy can be defined, and any quantum formula can be reobtained in the classical case. Of course we can directly work in this regime if we consider the classical analogue of our formalism and the poles of the classical Liouville operator . But now we know, from papers and , that this classical conditional entropy never decreases. So we can consider $`\rho _{W\text{ }}`$ and make a projection obtaining $`\stackrel{~}{\rho }_W`$ and define:
$$S(\stackrel{~}{\rho }_W)=_\mathrm{\Gamma }\stackrel{~}{\rho }_W\mathrm{log}\frac{\stackrel{~}{\rho }_W}{\rho _W}dqdp$$
(177)
As this entropy is never decreasing we know that $`\stackrel{}{S(\stackrel{~}{\rho }_W)}0.`$ Also as $`\rho _W`$ is given by the sum of local terms (175) this definition contains the notion of local equilibrium entropy. Now we can repeat everything we have done in the quantum case, since condition 1, 2, and 3 can be translated to the quantum case. Precisely: Condition 1 using eq. (168). Condition 2 using eq. (171). Condition 3 using eq. (166). So we obtain:
$$S(\stackrel{~}{\rho }_W)=\frac{1}{2}e^{2\gamma _{\mathrm{min}}t}_\mathrm{\Gamma }\frac{\stackrel{~}{\rho }_{1W}^2}{\rho _W}𝑑q𝑑p+\mathrm{}$$
(178)
where the dots symbolize higher order terms. Thus:
$$\stackrel{}{S}(\stackrel{~}{\rho }_W)=\gamma _{\mathrm{min}}e^{2\gamma _{\mathrm{min}}t}_\mathrm{\Gamma }\frac{\stackrel{~}{\rho }_{1W}^2}{\rho _W}𝑑q𝑑p+\mathrm{}$$
(179)
So now, it is clear that $`\stackrel{}{S}(\stackrel{~}{\rho }_W)>0`$ and that we have obtained the second law of the thermodynamics. This is not at all surprising since it is well known that, by a projection, the constant $`S(\rho _W)`$ can be transformed in the variable $`S(\stackrel{~}{\rho }_W).`$ We have just showed how our method works and also how we can obtain the classical eqs. (174) and (178).
So our $`S(\stackrel{~}{\rho }_W)`$ has all the properties that we have announced and we can claim that it is a good candidate to play the role of internal entropy.
## XXI Figures.
Fig 1. Plane $`\nu `$ and curve $`C.`$
Fig 2. A scattering experiment.
Fig. 3. The left hand side of fig 1: a growing process.
Fig. 4. The right hand side of fig. 1: de decaying process.
Fig. 5. The quantum branch system of the universe.
Fig. 6. Fig. 4 time-inverted.
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# 1 Introduction
## 1 Introduction
Recently, hybrid algebraic-Schrödinger approaches for the study of transition intensities in molecules have been introduced \[ReferencesReferences\]. In these approaches one relies on Schrödinger/algebraic correspondences to introduce operators in the algebra which are functions of configuration space parameters. For the $`U(2)`$ model considered in \[References,References\] one does so by relying on correspondences between a particular dynamical symmetry and a Schrödinger Hamiltonian \[References\]. For more complicated algebras such correspondences are not available. Further, the inner-product of two algebraic states is not equal to the Schrödinger inner-product of their images. Since, in some circumstances, the inner-product of the algebra has a natural interpretation as a Schrödinger overlap \[References\] such shortcomings are an impediment to understanding molecular structure.
It would therefore be useful to have a generic correspondence —independent of a particular dynamical symmetry— for the interpretation of algebraic parameters as geometric quantities which reside in configuration space (which is noticeably absent in the algebraic models). We provide such an interpretation by exploiting the many approximate correspondences between algebraic and Schrödinger pictures. We demonstrate how one may turn an algebraic hamiltonian into a traditional Schrödinger (traditional meaning kinetic plus potential—with no coordinate dependent mass terms) single particle hamiltonian which will give the same results to leading order in $`N`$, the label for the symmetric representation of the algebra. Our procedure is most similar to the re-quantization used by \[References\] to study spectra. We proceed further to study the wavefunctions and inner product. Our approach differs from other studies which have examined the correspondence between algebraic and Schrödinger inner products \[References\] since they rely on particular dynamical symmetries. We work out our procedure explicitly for $`U(3)`$ and conclude that different chains correspond not only to different geometries but different scales.
## 2 Algebraic Models: Constructing a Connection with Configuration Space
Numerous correspondences between algebraic and Schrödinger methods and their classical limits can be succinctly summarized in the following simple diagram:
(2.1)
Hamiltonians on the right hand side are written in terms of geometric parameters whereas those on the left are written in terms of algebraic ones. Thus, a careful exploitation of the horizontal correspondences should relate the different parameters. We propose to do so by travelling counter-clockwise around the diagram—re-quantizing the algebraic model.
The vertical correspondences are extremely well defined and have been thoroughly discussed in the literature \[References,ReferencesReferences\]. The relationship between the $`U(n)U(n1)`$ and the simple harmonic oscillator (SHO) has also been discussed \[References\] and can be formalized by considering the contraction limit of the algebra \[References,References\]. Similarly, if one finds appropriate coordinates in a patch of the coherent state phase space (CSPS) one may embed this region within a standard phase space and have a perfect copy of the dynamics for local trajectories \[References\].
The horizontal correspondences are approximate on the whole given that the phase spaces are topologically different (compact versus non-compact) in the classical regimes and one has a finite versus infinite dimensional Hilbert space in the corresponding quantum cases. However, with careful selection of coordinates it will be seen that these incompatabilities will not effect the lower bound states in a re-quantization scheme.
It is to the selection of these coordinates we now turn. We must find coordinates in a patch of CSPS which behave like positions and momenta and then identify these as coordinates covering the entire canonical phase space. In the language of geometry, we wish to find coordinates in which the symplectic form determining the dynamics is in Darboux form: $`\omega =dpdq`$ \[References\]. Historically, it has been useful to think of the imaginary part of the coherent state parameter $`\alpha `$ as a momenta \[References,References\]. We adopt that convention here to eliminate any remaining coordinate choice ambiguity.
Examining the action in the propagator path integral \[References,References\]: one finds that the classical hamiltonian is given by the coherent state expectation, $`H=\alpha |\widehat{H}|\alpha `$; $`|\alpha =\frac{1}{\sqrt{N!}}(\sqrt{1\alpha \alpha ^{}}\sigma ^{}+\alpha \tau ^{})^N|0`$ denotes the coherent state in group coordinates \[References\]; $`\sigma `$ and $`\tau `$ are typically scalar and tensor operators of $`O(3)`$ respectively; and the hamiltonian evolution is determined by the symplectic form $`\omega =i\mathrm{}d\alpha |d|\alpha `$ which in group coordinates is explicitly:
$`\omega =Ni\mathrm{}d\alpha _\mu ^{}d\alpha _\mu .`$ (2.2)
Deviating slightly from the typical development we let
$`\alpha _\mu =\sqrt{{\displaystyle \frac{\overline{m\omega }}{2}}}({\displaystyle \frac{q_\mu }{\sqrt{N}}}+{\displaystyle \frac{i}{\overline{m\omega }}}{\displaystyle \frac{p_\mu }{\sqrt{N}}}),`$ (2.3)
where $`q`$ and $`p`$ are now dimensionful quantities and $`\overline{m\omega }`$ has units of length over momentum and has the physical interpretation of the ratio of the natural distance and momenta scales of the problem (see appendix A). In these coordinates the symplectic two-form is in standard Darboux form $`\omega =dp_\mu dq_\mu `$, or equivalently $`\{q,p\}_{\mathrm{PB}}=1`$ .
We have been extremely explicit in order to contrast our choice with those of the literature. For instance in the coordinates of \[References\] $`\{\stackrel{~}{q},\stackrel{~}{p}\}_{\mathrm{PB}}=N`$. This point is obfuscated because the correct equations of motion are maintained by dividing the classical action (and hence $`\omega `$) by $`N`$. Such a procedure, however, leads to improper quantization. Similarly in terms of the projective coordinates \[References\] used to establish algebraic-geometric correspondences in such references as \[References,References\] one has $`\{\stackrel{~}{q},\stackrel{~}{p}\}_{\mathrm{PB}}=N(1+\stackrel{~}{q}^2+\stackrel{~}{p}^2)^2`$. Thus only near the domain $`\stackrel{~}{p}^2+\stackrel{~}{q}^20`$ is it suitable to interpret these coordinates as the natural coordinates of a cotangent bundle over configuration space.
Classical Hamiltonians obtained from our coordinates usually have contributions such as $`p_r^4`$ and $`p_r^2r^2`$. This leads to canonical quantization ordering ambiguities as well as higher degree differential equations. However, if we focus on the lowest bound states we may use an approximate hamiltonian which is valid in the classical regions whose paths most greatly contribute to these states. In the spirit of the stationary phase approximation, the appropriate region would be near the fixed point of the hamiltonian flow, i.e. where $`dH=0`$.
Fortunately, since the CSPS limit of algebraic hamiltonians typically have a convenient momenta dependence, i.e. $`H=p_r^2f(p_r^2,r^2)+V_{\mathrm{eff}}(r)`$, this condition implies we are looking in the region of reduced phase space near $`p_r=0`$ and $`r=r^{}`$ where $`V_{\mathrm{eff}}^{}(r^{})=0`$. These conditions are eminently reasonable.
Near this region the dynamics can be given by approximating the hamiltonian:
$`H(p_r,r)H|_{p_r=0,r=r^{}}+{\displaystyle \frac{1}{2}}\left(V_{\mathrm{eff}}^{\prime \prime }(r^{})(rr^{})^2+H_{1^2}|_{p=0,r=r^{}}p_r^2\right),`$ (2.4)
where all other lower order terms vanish. Hence, at least locally around the fixed point of the system, our ordering ambiguities are resolved. If $`r^{}`$ is sufficiently far from the $`r=0`$ phase space boundary one can sensically re-quantize.
Given our coordinate choice each higher derivative in the taylor expansion of $`H`$ in equation 2.4 will be down by a factor of $`\frac{1}{\sqrt{N}}`$. Thus this approximation may be considered as an expansion in large $`N`$. Note, however that this procedure is distinct from other ‘large $`N`$’ techniques such as the contraction limit \[References, References\]. Indeed, we will see in section 3.2 that this procedure, in the specific instance for Hamiltonians with the dynamical symmetry $`U(n)U(n1)`$, essentially reproduces the contraction limit results.
We have carried out this discussion in a reduced phase space with a radial coordinate. That is, due to $`O(3)`$ invariance the fixed point is not a true minima in the global phase space. In instances where one has a true fixed point in the global phase space (e.g. when the minima is at $`r=0`$) a similar procedure applies.
The introduction of the scales in equation 2.3 appears arbitray. However, since we desire a generic correspondence which preserves the inner product structure of the algebra in the Schrödinger picture, once a scale is picked for one Hamiltonian it must be fixed for others, i.e. any scale change would influence the overlap of eigenstates which should be equated with the inner product of the representation. Thus, the scale for each hamiltonian must be functionally related to the scale of another. If the relationship is not trivial $`\overline{m\omega }`$ may depend on the parameters within the hamiltonian. This dependence of the scale ratio, $`\overline{m\omega }`$, will be suppressed until such subtleties must be considered, at which time we will supplement our notation with subscripts.
## 3 An Explicit Example: $`U(3)`$
### 3.1 Coherent State Limit
The algebraic approach for 2D problems was presented by \[References\]. One considers symmetric representations of $`u(3)`$ realized with the two chains of interest:
$`\begin{array}{cc}U(3)U(2)O(2)\hfill & \mathrm{I}\hfill \\ U(3)O(3)O(2)\hfill & \mathrm{II}\hfill \end{array}.`$ (3.3)
We use the same notation as \[References\] for the generators. However, we select a different $`O(3)`$ subgroup which is generated by $`\widehat{R}_+`$, $`\widehat{R}_{}`$, and $`\widehat{l}`$. Our choice introduces more ‘coordinate like’ terms of the classical Hamiltonian. The choice makes no difference for spectra and only an overall phase for FC overlaps \[References\]. Please note that the $`O(3)`$ group is a dynamical symmetry subgroup and does not have the interpretation of a rotation in configuration space.
The general Hamiltonian of a $`U(3)`$ model is
$$H=ϵ\widehat{n}+\delta \widehat{n}(\widehat{n}+1)+\beta \widehat{l}^2A\widehat{W}^2,$$
(3.4)
where $`ϵ`$, $`\delta `$, and $`A`$ are taken as positive or $`0`$. Setting $`A=0`$ ($`ϵ=\delta =0`$) gives a Hamiltonian with dynamical symmetry I (II). The basis corresponding to chain I is labeled by the eigenvalues of $`\widehat{n}`$ and $`\widehat{l}`$ respectively. The basis corresponding to chain II is labeled by the eigenvalues of $`\widehat{W}^2=\frac{1}{2}(\widehat{R}_+\widehat{R}_{}+\widehat{R}_{}\widehat{R}_+)+\widehat{l}^2`$ and again $`\widehat{l}`$.
Taking the appropriate expectations one may calculate the classical limits of the various operators. The results are enumerated in Table 1.
We reduce the phase space in the natural way:
$$\begin{array}{cc}q_x=r\mathrm{cos}\theta & q_y=r\mathrm{sin}\theta \\ p_x=p_r\mathrm{cos}\theta p_\theta r\mathrm{sin}\theta & p_y=p_r\mathrm{sin}\theta +p_\theta r\mathrm{cos}\theta ,\end{array}$$
(3.5)
implying
$$qq=r^2,pp=p_r^2+\frac{l_{\mathrm{cl}}^2}{r^2},$$
(3.6)
where $`l_{\mathrm{cl}}^2=p_\theta ^2r^4`$ is a constant of the motion.
### 3.2 $`U(2)`$ Chain
We work out the re-quantization procedure in the case where $`A=\beta =\delta =0`$. In this instance we have the Hamiltonian on the reduced phase space:
$`H_{\mathrm{cl}}=ϵ{\displaystyle \frac{\overline{m\omega }}{2\mathrm{}}}\left({\displaystyle \frac{1}{(\overline{m\omega })^2}}(p_r^2+{\displaystyle \frac{l_{\mathrm{cl}}^2}{r^2}})+r^2\right).`$ (3.7)
The canonical quantization scheme yields the following prescription:
$`p_r^2{\displaystyle \frac{\mathrm{}^2}{r}}{\displaystyle \frac{d}{dr}}\left(r{\displaystyle \frac{d}{dr}}()\right),l\mathrm{}l,`$ (3.8)
where $`l`$ now labels irreps. of $`SO(2)`$. This yields the quantum Hamiltonian:
$`H=ϵ{\displaystyle \frac{\mathrm{}}{2\overline{m\omega }}}\left({\displaystyle \frac{1}{r}}{\displaystyle \frac{d}{dr}}\left(r{\displaystyle \frac{d}{dr}}()\right)+{\displaystyle \frac{l^2}{r^2}}+{\displaystyle \frac{(\overline{m\omega })^2}{\mathrm{}^2}}r^2\right),`$ (3.9)
which we immediately recognize as a multiple of a $`2D`$ circular oscillator with energy levels spaced by $`\mathrm{\Delta }E=ϵ`$ \[References\]. This result perfectly agrees with direct algebraic evaluation.
We proceed similarly for $`\delta `$ different from $`0`$. The classical Hamiltonian is then of the form:
$$H=N(ϵ+2\delta )n_{\mathrm{cl}}+\delta N(N1)(n_{\mathrm{cl}})^2.$$
(3.10)
In Cartesian coordinates the conditions for a fixed point are:
$$\begin{array}{c}dH=0=N[ϵ+2\delta +2\delta (N1)n_{\mathrm{cl}}]\hfill \\ \times \left(\frac{n_{\mathrm{cl}}}{x}dx+\frac{n_{\mathrm{cl}}}{y}dy+\frac{n_{\mathrm{cl}}}{p_x}dp_x+\frac{n_{\mathrm{cl}}}{p_y}dp_y\right),\hfill \end{array}$$
(3.11)
which is true at $`x=y=p_x=p_y=0`$ (hence our use of Cartesian as opposed to radial coordinates). Near this point to second order our Hamiltonian behaves like:
$$HN(ϵ+2\delta )n_{\mathrm{cl}}.$$
(3.12)
Re-quantization proceeds exactly as in the previous case leading to an energy spacing of $`\mathrm{\Delta }Eϵ+2\delta `$. The corresponding algebraic Hamiltonian ($`H=ϵ\widehat{n}+\delta \widehat{n}(\widehat{n}+1)`$) has spacing for the lowest levels of $`\mathrm{\Delta }E=ϵ+2\delta `$ .
Additionally the same low energy eigenfunctions diagonalize both $`U(2)`$ Hamiltonians in both the algebraic and approximate re-quantizing prescriptions (i.e. in both prescriptions the inner-product of the lowest eigenfunctions of the Hamiltonain with $`\delta =0`$ and $`\delta 0`$ is $`1`$).
### 3.3 Radially Displaced Oscillators—Near the $`O(3)`$ Chain
In the situation where our Hamiltonian is near the $`O(3)`$ dynamical symmetry the fixed point of the Hamiltonian flow is not at $`r=0`$.
The potential is given by:
$$\begin{array}{c}V_{\mathrm{eff}}(r)=\frac{\overline{m\omega }}{\mathrm{}}[(XZ)\frac{1}{2}r^2+(Y+Z)\frac{\overline{m\omega }}{4N\mathrm{}}r^4\\ +X\frac{1}{2}\frac{1}{(\overline{m\omega })^2}\frac{l^2}{r^2}+Y\frac{1}{4N\mathrm{}(\overline{m\omega })^3}\frac{l^4}{r^4}],\end{array}$$
(3.13)
where we have made the definitions $`X=ϵ+2\delta +\beta `$, $`Y=\delta (N1)`$, and $`Z=4A(N1)`$. The physically interesting regime \[References\] is ‘near’ the $`O(3)`$ chain, i.e. $`\frac{X}{Z}`$ and $`\frac{Y}{Z}`$ are both small. In this limit it is easy to find a perturbative solution for the condition $`V_{\mathrm{eff}}^{}(r^{})=0`$ :
$$(r^{})^2\frac{N\mathrm{}}{\overline{m\omega }}\left(1+\frac{X}{Z}(\frac{l^2}{(N\mathrm{})^2}1)\frac{Y}{Z}\right),$$
(3.14)
where now it is clear that the presumption that $`r^{}`$ is large depends on $`N`$ being large. Note that having $`X`$ or $`Y`$ different from $`0`$ decreases $`r^{}`$. That is, the $`O(3)`$ chain represents the maximum possible displacement. $`r^{}`$ may increase with increasing $`l`$ however, corresponding to a larger angular momentum barrier.
Since $`l`$ can be as large as $`N`$ we regard $`l^2/N^2`$ as an independent quantity—even though in practice it is of the same order as terms which go $`1/N^2`$. In anticipation of this we dropped the term going with $`\frac{l^2}{N^4}`$ for simplicity.
When one moves slightly off the $`O(3)`$ chain different $`l`$ subspaces which were degenerate split resulting in the lowest energy difference, $`\mathrm{\Delta }E`$, having a slight $`l`$ dependence. As we’ve seen this dependence has entered in lowest order in the form of terms $`\frac{l^2}{N^2}`$. Since the ratio $`\frac{l}{N}`$ can be regarded as a separate independent small quantity ignoring terms like $`O(\frac{1}{\sqrt{N}})`$ does not mean we loose all $`l`$ dependence in our calculation. Our results for the energy level spacing will only be correct to leading order in $`N`$. The $`l`$ space splitting behavior will be given correctly—but the actual values of the splitting will be insignificant compared to other (non-$`l`$ dependent) contributions we have ignored for smaller $`l`$.
Continuing we see:
$`V_{\mathrm{eff}}^{\prime \prime }(r^{})2Z{\displaystyle \frac{\overline{m\omega }}{\mathrm{}}}\left[1+{\displaystyle \frac{X}{Z}}\left(3{\displaystyle \frac{l^2}{(N\mathrm{})^2}}1\right)\right],`$ (3.15)
$`H_{p_r^2}|_{p_r=0,r=r^{}}{\displaystyle \frac{Z}{2\mathrm{}(\overline{m\omega })}}\left[1+{\displaystyle \frac{X}{Z}}({\displaystyle \frac{l^2}{(N\mathrm{})^2}}+1)+{\displaystyle \frac{Y}{Z}}(2{\displaystyle \frac{l^2}{(N\mathrm{})^2}}+1)\right].`$ (3.16)
Near the fixed point the Hamiltonian’s behavior is given by equation 2.4. However, the radial quantum mechanics depends crucially on the point $`r=0`$. Since equation 2.4 does not reproduce the original Hamiltonian’s behavior near this point the approximation becomes suspect. It would be more reasonable to quantize:
$$\frac{1}{2}H_{p_{r}^{}{}_{}{}^{2}}|_{p=0,r=r^{}}p_r^2+\stackrel{~}{V}(r),$$
(3.17)
where $`\stackrel{~}{V}(r)`$ has the same behavior near the fixed point as the approximate Hamiltonian and the behavior near $`r=0`$ of the original Hamiltonian. That is:
$`\underset{r0}{lim}r^4\stackrel{~}{V}(r)`$ $`=`$ $`{\displaystyle \frac{Y}{4}}{\displaystyle \frac{l^4}{N^2\mathrm{}^2(\overline{m\omega })^2}}`$ (3.18)
$`\stackrel{~}{V}(r)`$ $`=`$ $`H|_{p_r=0,r=r^{}}+{\displaystyle \frac{1}{2}}V_{\mathrm{eff}}^{\prime \prime }(r^{})(rr^{})^2+O((rr^{})^3).`$ (3.19)
Additionally, of course, $`\stackrel{~}{V}`$ should not introduce any other minima at other locations.
Following the prescription for the harmonic oscillator we use the radial quantization scheme (equation 3.8), substituting for the wavefunction, $`\mathrm{\Psi }=\frac{\psi }{\sqrt{r}}`$, we obtain:
$$H_{p_{r}^{}{}_{}{}^{2}}\frac{\mathrm{}^2}{2}\psi ^{\prime \prime }+\left[\stackrel{~}{V}(r)\frac{H_{p_{r}^{}{}_{}{}^{2}}\mathrm{}^2}{8r^2}\right]\psi =E\psi ,$$
(3.20)
where $`\psi (0)=0`$ and $`\psi (r\mathrm{})=0`$. The form of this equation is that of a 1-D Schrödinger equation except the left boundary condition is applied at $`r=0`$. At this point the $`r0`$ dependence of $`\stackrel{~}{V}`$ becomes crucial. Since $`\stackrel{~}{V}`$ blows up more quickly than the counter term blows down the sum of the two still has a minima, although it should be shifted inward from $`r^{}`$. Calling the new minima $`r^{}=r^{}(1ϵ)`$ and using our knowledge of the behavior of $`\stackrel{~}{V}`$ near $`r^{}`$ we find that the condition $`\frac{d}{dr}|_{r=r^{}}\left[\stackrel{~}{V}(r)\frac{H_{p_{r}^{}{}_{}{}^{2}}\mathrm{}^2}{8r^2}\right]=0`$ yields:
$$ϵ\frac{1}{16N^23}\left[1\frac{16N^2}{16N^23}\left(4\frac{X}{Z}(\frac{l^2}{N^2}1)\frac{Y}{Z}(2\frac{l^2}{N^2}+3)\right)\right]$$
(3.21)
This can be expanded in $`N`$ with leading contribution $`O(\frac{1}{N^2})`$. Recalling that $`N`$ essentially counts the number of bound states (hence effectively measures the depth of the potential) it is reasonable that $`\frac{1}{N^2}`$ should parameterize $`ϵ`$. Of course, since we are only working in leading order in $`N`$ we have $`ϵ0`$ and consequently $`r^{}r^{}`$. Our primary intent of displaying equation 3.21 is that the condition of $`N`$ being large is now clarified. Since the prefactor must be small we find $`2N1`$.
Returning to equation 3.20 we notice that in our approximation ($`r^{}`$ or $`r^{}`$ is large) a solution with the left boundary condition at $`r=\mathrm{}`$ will, to a good approximation, satisfy the boundary condition at $`r=0`$. In this instance the lower eigenvalues are given by $`E\mathrm{}\omega `$ where
$$\omega ^2=H_{p_{r}^{}{}_{}{}^{2}}\left[\stackrel{~}{V}^{\prime \prime }(r^{})\frac{3H_{p_{r}^{}{}_{}{}^{2}}\mathrm{}^2}{4(r^{})^4}\right].$$
(3.22)
Given the properties of $`\stackrel{~}{V}`$ we have $`\stackrel{~}{V}^{\prime \prime }(r^{})\stackrel{~}{V}^{\prime \prime }(r^{})=V_{\mathrm{eff}}^{\prime \prime }(r^{})`$. The second term is $`O(\frac{1}{N^2})`$ and may be ignored.
Thus we have
$$\mathrm{\Delta }EZ\left[1+2\frac{X}{Z}\frac{l^2}{N^2}+\frac{1}{2}\frac{Y}{Z}\left(2\frac{l^2}{N^2}+1\right)\right].$$
(3.23)
We see that on the $`O(3)`$ chain ($`X=Y=0`$) we have a spacing of $`\mathrm{\Delta }EZ4AN`$, as compared to the exact algebraic expression:
$$\mathrm{\Delta }E=A\left[N(N+1)(N2)(N1)\right]=A(4N2),$$
(3.24)
which agrees to leading order in $`N`$ as advertised.
Next we compute the induced harmonic dilatation constant of the geometry (see Appendix A):
$`\alpha ^2{\displaystyle \frac{\omega }{\mathrm{}H_{1^2}}}2{\displaystyle \frac{(\overline{m\omega })}{\mathrm{}}}\left[1+{\displaystyle \frac{X}{Z}}\left({\displaystyle \frac{l^2}{N^2}}1\right){\displaystyle \frac{Y}{Z}}\left({\displaystyle \frac{l^2}{N^2}}+{\displaystyle \frac{1}{2}}\right)\right].`$ (3.25)
Note, that the ‘concavity’ corrections depend explicitly on the algebraic parameters, indicative of the fact that our calculations have established relationships between algebraic Hamiltonians and the geometry of configuration space.
Equations 3.25 and 3.14 for the harmonic dilatation and radial displacement can be easily related to experimental data and provide useful assistance when analyzing transition intensities as to be detailed in the upcoming publication \[References\].
### 3.4 Near $`O(3)`$ to $`O(3)`$ Inner Product
We wish to calculate the inner-product of eigenstates of a hamiltonian corresponding to slightly off $`O(3)`$ symmetry to one on $`O(3)`$. In either case the approximate wavefunctions are:
$$\psi _0^{\mathrm{disp}}\frac{\sqrt{\alpha }}{\pi ^{\frac{1}{4}}}\frac{e^{\frac{1}{2}[\alpha (rr^{})]^2}}{\sqrt{r}},$$
(3.26)
where we have normalized it on the interval $`[\mathrm{},\mathrm{}]`$ (as opposed to $`[0,\mathrm{}]`$) as this only introduces errors of order $`\frac{1}{N}`$. The parameters $`\alpha `$ and $`r^{}`$ are determined by 3.25 and 3.14. We use $`\alpha ^{}`$ and $`r_{}^{}{}_{}{}^{}`$ as the parameters of the wavefunction exactly on the $`O(3)`$ chain (determined by the same equations except $`X=Y=0`$). Note that there is an $`l`$ dependence hidden in both $`r^{}`$ and $`\alpha `$.
Since both wavefunctions are essentially $`0`$ for $`r<0`$ we may again calculate the overlap on the interval $`[\mathrm{},\mathrm{}]`$. Thus we wish to evaluate the integral
$$I_{0,0}_{\mathrm{}}^{\mathrm{}}\varphi (\alpha ;rr^{})\varphi (\alpha ^{};rr_{}^{}{}_{}{}^{})𝑑r,$$
(3.27)
where $`\varphi `$ is the ground state wavefunction of a 1D SHO. The integral is easily evaluated to \[References\]:
$$I_{0,0}\left[\frac{2\alpha \alpha ^{}}{\alpha ^2+\alpha ^2}\right]^{\frac{1}{2}}\times \mathrm{exp}\left[\frac{(\alpha \alpha ^{}(r_{}^{}{}_{}{}^{}r^{}))^2}{2(\alpha ^2+\alpha ^2)}\right].$$
(3.28)
As stated previously the ratio of scales $`\overline{m\omega }`$ need not be the same in both scenarios and may indeed depend on the value of $`X`$,$`Y`$, and $`Z`$. However, this dependance must be smooth. Since the difference between these two scenarios is perturbative one concludes that $`\frac{\overline{m\omega }}{\overline{m\omega }^{}}=1+\delta `$ where $`\delta O(X/Z,Y/Z)`$. Given this, we have:
$$I_{0,0}1+O((\frac{X}{Z})^2,(\frac{Y}{Z})^2)$$
(3.29)
This is the exact overlap one would expect from perturbation theory of the algebraic model as spelled out in Appendix B.
### 3.5 $`O(3)`$ to $`U(2)`$ Inner-Products
Given equation 3.29, to our order the overlaps from a hamiltonian near $`O(3)`$ to the $`U(2)`$ chain should be exactly the same as the exact $`O(3)`$-$`U(2)`$ factors.
For the $`U(2)`$ chain we have the ground state wavefunction \[References\]:
$$\psi _0^{U(2)}=\sqrt{\frac{2}{|l|!}}(\alpha _{U(2)})^{|l|+1}r^{|l|}e^{\frac{1}{2}[\alpha _{U(2)}r]^2},$$
(3.30)
where $`(\alpha _{U(2)})^2=\frac{\overline{m\omega }_{U(2)}}{\mathrm{}}`$.
Thus our overlap becomes:
$`I_{0,0}({\displaystyle \frac{\alpha _{U(2)}}{\alpha }},\alpha r^{})=\sqrt{{\displaystyle \frac{2\alpha }{|l|!}}}{\displaystyle \frac{(\alpha _{U(2)})^{|l|+1}}{\pi ^{\frac{1}{4}}}}\mathrm{exp}\left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{\alpha ^2\alpha _{U(2)}^2}{\alpha _{+}^{}{}_{}{}^{2}}}r_{}^{}{}_{}{}^{2}\right]`$
$`\times {\displaystyle _0^{\mathrm{}}}drr^{|l|+\frac{1}{2}}\mathrm{exp}[{\displaystyle \frac{1}{2}}(\alpha _+r{\displaystyle \frac{\alpha ^2r^{}}{\alpha _+}})^2]`$
where $`\alpha _+^2=\alpha _{U(2)}^2+\alpha ^2`$, and $`\alpha `$ is again defined by 3.25 with $`X=Y=0`$ . The integral may be expressed in terms of confluent hypergeometric functions:
$`{\displaystyle _0^{\mathrm{}}}𝑑rr^{|l|+\frac{1}{2}}\mathrm{exp}\left[{\displaystyle \frac{1}{2}}(\alpha _+r{\displaystyle \frac{\alpha ^2r^{}}{\alpha _+}})^2\right]=2^{\frac{2l1}{4}}{\displaystyle \frac{1}{\alpha _{+}^{}{}_{}{}^{\frac{3+2l}{2}}}}\mathrm{exp}\left[{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\alpha ^2r^{}}{\alpha _+}}\right)^2\right]`$
$`\times \{\mathrm{\Gamma }\left({\displaystyle \frac{3+2l}{4}}\right)_1F_1({\displaystyle \frac{3+2l}{4}},{\displaystyle \frac{1}{2}};{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\alpha ^2r^{}}{\alpha _+}}\right)^2)`$
$`+\sqrt{2}\left({\displaystyle \frac{\alpha ^2r^{}}{\alpha _+}}\right)\mathrm{\Gamma }\left({\displaystyle \frac{5+2l}{4}}\right)_1F_1({\displaystyle \frac{5+2l}{4}},{\displaystyle \frac{1}{2}};{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\alpha ^2r^{}}{\alpha _+}}\right)^2)\}.`$
We note that $`\frac{\alpha ^2r^{}}{\alpha _+}N`$ and hence use an asymptotic expansion for the hypergeometric functions \[References\]. Keeping only the most dominant contribution of the integral we find:
$$I_{0,0}\frac{\pi ^{\frac{1}{4}}}{\sqrt{|l|!}}\left(\frac{\alpha _{U(2)}}{\alpha }\right)^{|l|+1}\left(\frac{\alpha }{\alpha _+}\right)^{4+2|l|}\mathrm{exp}\left[\frac{1}{2}\frac{\alpha ^2\alpha _{U(2)}^2}{\alpha _{+}^{}{}_{}{}^{2}}r_{}^{}{}_{}{}^{2}\right](\alpha r^{})^{\frac{5+2|l|}{2}}.$$
(3.33)
Substituting the expressions for the harmonic scales, and letting $`\zeta =\frac{\overline{m\omega }_{U(2)}}{\overline{m\omega }_{O(3)}}`$ we find
$$I_{0,0}\frac{\pi ^{\frac{1}{4}}}{\sqrt{|l|!}}\left(\frac{\zeta }{2}\right)^{|l|+1}\left(\frac{2}{\zeta +2}\right)^{2+|l|}\mathrm{exp}\left[N\frac{1}{1+\frac{2}{\zeta }}\right](2N)^{\frac{5+2|l|}{4}}.$$
(3.34)
### 3.6 Scale Changes
Comparing the analytic expression of the $`U(2)`$/$`O(3)`$ overlaps (equation 3.34 calculated by approximately requantizing) and the expression from appendix equation B.5 (calculated by direct algebraic means) fixes the scale dependance. Comparing the exponential dependance of the two expressions,
$$e^{\frac{\mathrm{ln}\mathrm{\hspace{0.17em}2}}{2}N}\mathrm{vs}.e^{\left[N\frac{1}{1+\frac{2}{\zeta }}\right]},$$
(3.35)
leads to the conclusion that:
$$\zeta =\frac{\overline{m\omega }_{U(2)}}{\overline{m\omega }_{O(3)}}\frac{2\mathrm{ln}2}{2\mathrm{ln}2}.$$
(3.36)
The rest of the dependance can be fixed by allowing $`\zeta `$ to have $`\frac{\mathrm{log}N}{N}`$ and $`\frac{1}{N}`$ corrections. To see this we let:
$$\zeta \frac{2\mathrm{log}2}{2\mathrm{log}2}\left(1\frac{4}{\mathrm{log}2(2\mathrm{log}2)}\frac{\mathrm{log}\gamma }{N}\right),$$
(3.37)
or equivalently
$$(1+\frac{2}{\zeta })^1\frac{\mathrm{log}2}{2}\frac{\mathrm{log}\gamma }{N}.$$
(3.38)
Substituting into the expression 3.34 we find the first order term of the large $`N`$ asymptotic:
$$I_{0,0}\frac{\pi ^{\frac{1}{4}}}{\sqrt{|l|!}}\left[\frac{\mathrm{log}2}{2}\right]^{|l|+1}\left[\frac{2\mathrm{log}2}{2}\right](2N)^{\frac{5+2|l|}{4}}e^{N\frac{\mathrm{log}2}{2}}\gamma .$$
(3.39)
Comparing with the algebraic expression we conclude that
$$\gamma =\frac{2^{\frac{3}{4}}}{2\mathrm{log}2}(\mathrm{log}2)^{(|l|+1)}N^{\frac{3}{2}}$$
(3.40)
That is, in order for the Schrödinger and algebraic prescriptions to be commensurate for the leading asymptotic in $`N`$, we must have the first orders of the asymptotics of $`\zeta `$:
$`\zeta {\displaystyle \frac{2\mathrm{log}2}{2\mathrm{log}2}}\left[1+{\displaystyle \frac{4}{\mathrm{log}2(2\mathrm{log}2)}}\left\{{\displaystyle \frac{3}{2}}{\displaystyle \frac{\mathrm{log}N}{N}}+(|l|\mathrm{log}\mathrm{log}2+c){\displaystyle \frac{1}{N}}\right\}\right],`$ (3.41)
$`c=\mathrm{log}\mathrm{log}2\mathrm{log}{\displaystyle \frac{2^{\frac{3}{4}}}{2\mathrm{log}2}}`$ (3.42)
Calculating the inner product of the eigenstates of an algebraic hamiltonian on the $`U(2)`$ chain with the eigenstates of another hamiltonian off of the $`U(2)`$ chain is not simply analogous to the overlap of radially displaced oscillators, but analagous to the matrix elements of an operator which radially displaces and dilatates (much like the operator matrix elements calculated in \[References\]) changing the natural scale of the problem. The degree of dilatation depends on the proximity of the second hamiltonian to either chain. For chains near $`U(2)`$ the dilatation paramater is essentially $`1`$. As one moves nearer to the $`O(3)`$ chain the dilatation parameter increases, approaching the value given by equation 3.41. This has considerable consequences for hybrid algebraic-Schrödinger analysis of molecules \[References\].
## 4 Conclusions
We have provided a procedure, via requantization, to convert a hamiltonian of an algebraic model into a Schrödinger hamiltonian which will give the same results for the lower states to leading order in $`N`$. The procedure is optimized for algebras which are interpreted as single particle hamiltonians. Thus, it seems most applicable to molecular models such as the vibron model $`U(4)`$ \[References\], the anharmonic oscillator $`U(2)`$ \[References\], and the two dimensional $`U(3)`$ model considered here. Although the prescription should generate Schrödinger hamiltonians in other models it remains to be seen whether or not such results would have appropriate many-body interpretations.
We have carried out the requantization process in detail for the limiting cases of the $`U(3)`$ model. By demanding that the wavefunction overlap of the requantized system agree with the inner product of the algebraic model we have demonstrated that different chains of the $`U(3)`$ model correspond to not only different geometries but different scales.
## 5 Acknowledgements
This work was performed in part under the U.S. Department of Energy, Contract No. DE-FG02-91ER40608. I extend my deepest gratitude to my advisor, Prof. Franco Iachello, for introducing me to the problem and his suggestions and feedback on this work. I thank Prof. Dimitri Kusnezov for also making suggestions on this manuscript. Finally, I would like to thank the Department of Energy’s Institute for Nuclear Theory at the University of Washington for its hospitality during the completion of this work.
Appendices
## Appendix A Scales in the Schrödinger Picture
Consider a traditional Schrödinger Hamiltonian (1-D for simplicity),
$$H=\frac{\mathrm{}^2}{2m}\frac{d^2}{dx^2}+V_0f(\overline{\alpha }x),$$
(A.1)
with the distance scale $`\overline{\alpha }`$ and momenta scale given by $`\frac{\overline{\alpha }}{\overline{m\omega }}=(mV_0)^{\frac{1}{2}}`$.
We on occasion refer to the induced harmonic dilatation constant—which we define as the distance scale (dilatation constant of the wavefunctions) one obtains by approximating the potential about a minimum to second order. For a SHO $`f(x)=\frac{1}{2}x^2`$ and $`V_0=\frac{m\omega ^2}{\overline{\alpha }^2}`$. In this instance $`\overline{m\omega }=m\omega `$ and we define the distance scale $`\overline{\alpha }^2=\frac{m\omega }{\mathrm{}}`$. If $`x^{}`$ denotes the minimum of a more complicated potential we see that the distance scale associated with its harmonic approximation is
$$\alpha ^2=\frac{\overline{m\omega }}{\mathrm{}}\sqrt{f^{\prime \prime }(\overline{\alpha }x^{})}.$$
(A.2)
The induced scale depends on two quantities: (1) the ratio of the true distance and momenta scales; (2) the concavity of the potential about the minima.
## Appendix B Algebraic Calculations of Overlaps near $`O(3)`$
We wish to find the overlaps of eigenstates of the hamiltonian
$$H=A\left(\frac{(ϵ+\delta )}{A}\widehat{n}+\frac{\delta }{A}\widehat{n}(\widehat{n}+1)+\frac{\beta }{A}\widehat{l}^2\widehat{W}^2\right)$$
(B.1)
to the hamiltonian with $`ϵ=\delta =\beta =0`$ in the limit where $`A`$ is large. We may work within subspaces of constant $`l`$ so that there are no degeneracies.
We denote the basis from the $`O(3)`$ chain by $`|[N],\omega ,l`$ where $`\omega (\omega +1)`$ is the eigenvalue of $`\widehat{W}^2`$. The $`U(2)`$ chain basis is denoted by $`|[N],n,l`$ analogously. We denote the transformation matrix between them by $`\zeta `$, i.e. $`|[N],\omega =_n\zeta _n^\omega |[N],n`$, where we have suppressed the implicit $`l`$ dependance.
The wavefunction for the ground state may be perturbatively calculated to first order:
$$|[N],E_{\mathrm{low}}|[N],\omega =N+\underset{\omega ^{}N}{}|[N],\omega ^{}\frac{\underset{n}{}(\zeta _n^\omega ^{})^{}\zeta _n^{\omega =N}\left[(ϵ+2\delta )n+\delta n^2\right]}{AN(N+1)A\omega ^{}(\omega ^{}+1)}.$$
(B.2)
As always in perturbation theory, to first order in $`\frac{1}{A}`$ the wavefunctions only have corrections orthogonal to the unpeturbed wavefunction. Consequently the overlap from a slight deformation off the $`O(3)`$ chain to the $`O(3)`$ chain is identically $`1`$ to this order.
We now concentrate on the overlap between the $`O(3)`$ and $`U(2)`$ chains. In this limit we have:
$`I_{0,0}=[N],n=|l|,l|[N],\omega =N,l=\zeta _{n=|l|}^{\omega =N}.`$ (B.3)
These coefficients are calculated in \[References\]:
$`\zeta _{n=|l|}^{\omega =N}=\left[{\displaystyle \frac{(N+|l|)!}{2^{|l|}|l|!(2N1)!!}}\right]^{\frac{1}{2}}.`$ (B.4)
Using $`(2N1)!!=2^N\mathrm{\Gamma }(N+\frac{1}{2})/\mathrm{\Gamma }(\frac{1}{2})`$ and the Stirling approximation $`\mathrm{\Gamma }(z+b)\sqrt{2\pi }e^zz^{z+b\frac{1}{2}}`$ \[References\] this formula may be approximated for large $`N`$ by:
$`I_{0,0}{\displaystyle \frac{\pi ^{\frac{1}{4}}}{2^{\frac{|l|}{2}}\sqrt{|l|!}}}N^{\frac{|l|}{2}\frac{1}{4}}e^{\frac{\mathrm{log}\mathrm{\hspace{0.17em}2}}{2}N}.`$ (B.5)
If one wishes to algebraically calculate FC transitions to higher energy $`U(2)`$ bound states one may proceed in the exact same fashion to find
$`I_{n,0}{\displaystyle \frac{\pi ^{\frac{1}{4}}N^{\frac{n}{2}\frac{1}{4}}e^{\frac{\mathrm{log}\mathrm{\hspace{0.17em}2}}{2}N}}{2^{\frac{n}{2}}\sqrt{(\frac{n+l}{2})!(\frac{nl}{2})!}}}.`$ (B.6)
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# SELF-SIMILAR STATIC SOLUTIONS ADMITTING A TWO-SPACE OF CONSTANT CURVATURE
Departament de Física
Universitat de les Illes Balears
E-07071 Palma de Mallorca. SPAIN
Abstract
A recent result by Haggag and Hajj-Boutros is reviewed within the framework of self-similar space-times, extending, in some sense, their results and presenting a family of metrics consisting of all the static spherically symmetric perfect fluid solutions admitting a homothety.
PACS numbers: 04.20Jb, 02.40+m, 98.80Dr
In a recent paper, Haggag and Hajj-Boutros presented a static, spherically symmetric perfect fluid solution with a stiff-matter type equation of state (i.e.: $`p=\mu `$). By means of a few clever changes of coordinates, the authors reduce the problem to that of solving a non-linear, second order differential equation, whose polynomic solutions they investigate showing that only three such solutions exist, two of them being vacuum (flat Minkowski space-time and Schwarzschild solution) and the third one being that leading to the new metric referred to above, henceforth called HHB solution.
The purpose of this letter is to give all the static, spherically symmetric perfect fluid solutions admitting a homothety. This family can be completely characterized by means of a real parameter $`\gamma `$ (arising quite naturally from the equation of state for these fluids, see below), which must be in the interval $`[1,2]`$ in order to satisfy energy conditions. The two limiting values of $`\gamma `$, namely $`\gamma =1`$ and $`\gamma =2`$ correspond to Minkowski flat space-time and to the HHB solution respectively.
A few remarks concerning the similarity group and its action are in order here. It is a well known fact that an $`r`$-parameter group of homotheties $`H_r`$ (in which at least one proper homothety exists) always admits an $`(r1)`$-parameter subgroup of isometries $`G_{r1}`$. Now, the maximal dimension of the group of homotheties that a perfect fluid space-time may admit is $`r=7`$, in which case it is one of the special Robertson-Walker space-times , and therefore they are all known. The case $`r=6`$ is not compatible with an energy-momentum tensor of the perfect fluid type; thus, apart from the special Robertson-Walker solutions mentioned above, the highest dimension of the group of homotheties that a perfect fluid space-time may admit is $`r=5`$. In such case, the associated isometry subgroup $`G_4`$ has necessarily 3-dimensional non-null orbits . Notice that this is precisely the case we are interested in. We shall not treat here the case in full generality, namely; studying all perfect fluid space-times admitting an $`H_5`$ of homotheties, since this would be beyond the purpose of this letter, but we shall restrict ourselves to the case when the subgroup $`G_4`$ has timelike orbits $`T_3`$ and the subgroup $`G_3`$ that it necessarily contains has two-dimensional orbits. Everything else follows from these assumptions and the field equations. For further information on groups of homotheties and related issues, we refer the reader to and .
We start with a space-time that contains a non-null two-space of constant curvature (i.e.: there exists a three-parameter isometry group $`G_3`$ acting on this two-space). In this case the orbits $`V_2`$ admit orthogonal surfaces in $`M`$ . By performing a coordinate transformation in the two-spaces orthogonal to the Killing orbits the space-time metric can be put into diagonal form:
$$ds^2=A^2(r,t)(dt^2+dr^2)+B^2(r,t)(d\theta ^2+f^2(\theta ,k)d\varphi ^2)$$
(1)
$$f(\theta ,k)=\{\begin{array}{cc}\mathrm{sin}\theta \hfill & k=+1\hfill \\ \theta \hfill & k=0\hfill \\ \mathrm{sinh}\theta \hfill & k=1.\hfill \end{array}$$
(2)
where we have restricted ourselves to the case of spacelike Killing orbits, since perfect fluid and dust solutions cannot admit a group $`G_3`$ on two-dimensional timelike orbits .
Using the Jacobi identities and the fact that the Lie bracket of a proper homothetic vector field (HVF) and a Killing vector (KV) is a KV it can be easily shown that the HVF $`X`$ must be of either one of the following forms:
$`(I)X`$ $`=`$ $`X^t(r,t)_t+X^r(r,t)_r,k=1,0,1,`$ (3)
$`(II)X`$ $`=`$ $`X^t(r,t)_t+X^r(r,t)_r\theta _\theta ,k=0.`$ (4)
Now by using isotropic coordinates, one finds that static metrics can be expressed as :
$$ds^2=A^2(r)dt^2+B^2(r)[dr^2+r^2(d\theta ^2+f^2(\theta ,k)d\varphi ^2)],$$
(5)
where $`_t`$ is the hypersurface orthogonal timelike KV.
In this coordinate chart, the HVF in (3) and (4) takes the following forms:
$`(I)X`$ $`=`$ $`nt_t+R(r)_r,k=1,0,1,`$ (6)
$`(II)X`$ $`=`$ $`nt_t+R(r)_r\theta _\theta ,k=0,`$ (7)
where $`n`$ is a constant.
The homothetic equation $`_Xg_{ab}=2g_{ab}`$ specified to the components $`rr`$ and $`\theta \theta `$ of the metric (5), gives:
$$R_{,r}\frac{R}{r}X_{}^{\theta }{}_{,\theta }{}^{}=0,$$
(8)
and integrating, one gets:
$`(I)X`$ $`=`$ $`nt_t+qr_r,k=1,0,1,`$ (9)
$`(II)X`$ $`=`$ $`nt_t+(r\mathrm{ln}r+cr)_r\theta _\theta ,k=0,`$ (10)
where $`c`$ and $`q`$ ($`0`$) are constants.
Case ($`I`$)
By means of the coordinate transformation $`\widehat{r}=r^{1/q}`$, the HVF and the metric can be written as
$$X=nt_t+\widehat{r}_{\widehat{r}},$$
(11)
$$ds^2=\widehat{A}^2(\widehat{r})dt^2+\widehat{B}^2(\widehat{r})[q^2d\widehat{r}^2+\widehat{r}^2(d\theta ^2+f^2(\theta ,k)d\varphi ^2)].$$
(12)
The metric functions can be determined via the homothetic equations, that gives:
$$\widehat{B}=\mathrm{constant},\widehat{A}\widehat{r}^{1n}.$$
(13)
Defining a new radial coordinate $`r`$ as $`r=\widehat{r}\widehat{B}`$, one can come to the following simple forms for $`X`$ and the metric:
$$X=nt_t+r_r,$$
(14)
$$ds^2=r^{2(1n)}dt^2+q^2dr^2+r^2(d\theta ^2+f^2(\theta ,k)d\varphi ^2).$$
(15)
Case ($`II`$)
Imposition of the homothetic equations specified to the metric (5) and to the HVF (10), leads directly to:
$$ds^2=(\mathrm{ln}r+c)^{2(n1)}dt^2+\frac{b^2}{r^2(\mathrm{ln}r+c)^4}\left[dr^2+r^2(d\theta ^2+\theta ^2d\varphi ^2)\right],$$
(16)
$`b`$ and $`c`$ being constants.
From the expressions (15) and (16) of the metric, it is immediate to see that the components $`tt`$ of their respective Einstein tensors are negative for $`k=1`$ and $`0`$ (i.e.: hyperbolic and flat two-spaces) and therefore cannot verify energy conditions. Thus, it only remains to study the spherically symmetric case.
In this latter case, the field equations for a perfect fluid matter content lead to the metric:
$$ds^2=r^{44/\gamma }dt^2+\left(\frac{\gamma ^2+4\gamma 4}{\gamma ^2}\right)dr^2+r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2),$$
(17)
These metrics were already found, following a completely different approach, by Ibañez et al., and particular cases of them can be also found in Misner et al. (which are particular cases of Tolman class $`VI`$ solutions). Some particular cases (when the HVF is orthogonal to the fluid flow) were also studied by Herrera et al..
The matter variables being
$$\mu =\frac{1}{r^2}\left(\frac{4\gamma 4}{\gamma ^2+4\gamma 4}\right),$$
(18)
$$p=(\gamma 1)\mu $$
(19)
as one would have expected from $`p`$ and $`\mu `$ being functions of $`r`$ alone (and therefore, by the implicit function theorem, the fluid has a barotropic equation of state) and the space-time being self-similar . The HVF takes then the form:
$$X=\frac{2\gamma }{\gamma }t_t+r_r.$$
(20)
These are all the static, spherically-symmetric self-similar perfect fluid solutions. They are shear-free and have null volume expansion since the four-velocity $`u`$ of the fluid is parallel to the timelike KV, the vorticity is also zero (since $`u`$ is orthogonal to the orbits $`S_2`$ of the $`G_3`$ they contain); and the fluid has non-geodesic flow.
The particular case, $`\gamma =2`$, is the HHG solution and in this case the HVF $`X`$ becomes orthogonal to the fluid four-velocity, and for $`\gamma =1`$ the space-time is obviously flat.
It is interesting to notice that static and self-similar solutions admitting a two-space of constant curvature can only be spherically symmetric (irrespectively of the matter content) and that an $`H_5`$ static space-time with a $`G_3`$ acting on spacelike orbits, necessarily contains an $`H_4`$.
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# 1 Introduction
## 1 Introduction
Last two decades have witnessed an increase of interest to the multidimensional cosmology (see, for instance, \[5-7,9,11-18,20-25,27-29,31,32,39,41,42,45-48\] and references therein). According to this theory one assumes that the Universe had a higher dimension at a very early stage of its evolution, and that quantum processes have been responsible for the (topological) partition of the space, which provides us at present with the usual $`3`$-dimensional (external) space, in addition to internal space(s). The manifold which accounts for such a multidimensional spacetime has the following topology
$$M=R\times M_1\times \mathrm{}\times M_n,$$
(1.1)
where $`R`$ stands for the cosmic time axis, and the product with one part of manifolds $`M_1,\mathrm{},M_n`$ gives the external space, when the remaining part stands for internal spaces. The classical stage of the evolution is governed by the multidimensional version of Einstein’s equations. According to the classical description, the Riemann curvatures of spaces $`M_1,\mathrm{},M_n`$ are assumed to be constant (Einstein spaces). Such a model is the simplest multidimensional generalization of the space-time upon which the Friedman-Robertson-Walker (FRW) world model is based on. One easily understands that the use of extra dimensions (for the physical space-time) can be a sensible scenario only at primordial epochs, since the standard FRW world model is known to be in a sufficiently good agreement with the observational constraints down to a quite primordial epoch such as the nucleosynthesis era. Hence, it is clear that a reduction process (called dynamical compactification of additional dimensions) is required before such an epoch, to make the internal spaces contracting themself down to unobservable sizes.
Herein, we assume that the cosmic fluid (the source of the gravitational field at early stages) is viscous, which might simulate high energy physics processes (such as the particles creation). The effects related to viscosity in $`4`$-dimensional Universe were studied through different viewpoints (see e.g., \[2-4,8,10,26,30,33-38,44\]). Before developing the multidimensional model let us briefly discuss (extensive review of the subject was given by Gron ) the main trends in $`4`$-dimensional cosmology with viscous fluid as a source.
First, Misner considered neutrino viscosity as a mechanism of reducing the anisotropy in the Early Universe. Stewart and Collins and Stewart proved that it is possible only if initial anisotropies are small enough. Another series of papers which concerns the production of entropy in the viscous Universe was started by Weinberg . Both isotropization and production of entropy during lepton era in models of Bianchi types I,V were considered by Klimek . Caderni and Fabbri calculated coefficients of shear and bulk viscosity in plasma and lepton eras within the model of Bianchi type I. The next approach is connected with obtaining singularity free viscous solutions. The first nonsingular solution was obtained by Murphy within the flat FRW model with fluid possessing a bulk viscosity. However, Belinsky and Khalatnikov showed that this solution corresponds to the very peculiar choice of parameters and is unstable with respect to the anisotropy perturbations. Other nonsingular solutions with bulk viscosity were obtained by Novello and Araujo , Romero , Oliveira and Salim .
The crucial feature of each viscous cosmological model is assuming of the so called ”second equation of state”, which provides us with the viscosity coefficients dependence on time. Further we denote by $`\zeta `$ and $`\eta `$ the bulk and shear viscosity coefficients, correspondingly. Murphy integrated the 4-dimensional flat FRW model with bulk viscosity by assuming $`\zeta \rho `$ (as second equation of state), where $`\rho `$ is the density of the viscous fluid. Belinsky and Khalatnikov studied the behavior of this model as well as homogeneous anisotropic models of Bianchi types I and IX by means of qualitative methods with more general second equations of state $`\zeta ,\eta \rho ^\nu `$, where $`\nu `$ is constant. Lukacs integrated the homogeneous and isotropic $`4`$-dimensional model with a viscous pressureless fluid and a second equation of state given by $`\zeta [\mathrm{scale}\mathrm{factor}]^1`$. A curvature-dependent bulk viscosity was studied in multidimensional cosmology by Wolf . Recently, Motta and Tomimura studied a $`4`$-dimensional inhomogeneous cosmology with a bulk viscosity coefficient which depends on the metric. In our previous papers exact solutions in multidimensional models with bulk viscosity were obtained and their properties were studied for the following type of the second equation of state: $`\zeta [\mathrm{volume}\mathrm{of}\mathrm{the}\mathrm{Universe}]^1`$.
The aim of the present investigation is to integrate the Einstein equations for a multidimensional cosmological model formed by a chain of Ricci-flat spaces and a cosmic fluid possessing both shear and bulk viscosity. The second equations of state are chosen in the following form of metrical dependence of the bulk and shear viscosity coefficients: $`\zeta ,\eta [\mathrm{volume}\mathrm{of}\mathrm{the}\mathrm{Universe}]^1`$.
The paper is organized as follows. In Sec. 2 we describe the model and get basic equations. To integrate them, we develop some vector formalism proposed in our previous papers . Thermodynamical concepts in multidimensional cosmologies are defined in Sec. 3, where a formula which provides us with the variation rate of the entropy is derived. The equations of motion for the special set of parameters in the first and the second equations of state are integrated in Sec. 4. Exact solutions are given in a Kasner-like form, their physical properties are investigated in Sec. 5, where the process of dynamical compactification and the entropy production are also defined.
## 2 The model
Let us have the following metric
$$\mathrm{d}s^2=\mathrm{e}^{2\gamma (t)}\mathrm{d}t^2+\underset{i=1}{\overset{n}{}}\mathrm{exp}[2x^i(t)]\mathrm{d}s_i^2,$$
(2.1)
on the manifold defined in Eq. (1.1), where $`\mathrm{d}s_i^2`$ is the metric of the Einstein space $`M_i`$, $`\gamma (t)`$ and $`x^i(t)`$ are scalar functions of the cosmic time $`t`$. The dimension of this manifold is given by $`D=1+_{i=1}^nN_i`$, where $`N_i=\mathrm{dim}M_i`$. Herein, for reason of simplicity, only Ricci-flat spaces $`M_1,\mathrm{},M_n`$ are assumed (i.e., the components of the Ricci tensor for the metrics $`\mathrm{d}s_i^2`$ are zero). One easily obtains the non-zero components of the Ricci-tensor for the metric defined in Eq. (2.1) (see ):
$`R_0^0`$ $`=`$ $`\mathrm{e}^{2\gamma (t)}\left({\displaystyle \underset{i=1}{\overset{n}{}}}N_i(\dot{x}^i)^2+\ddot{\gamma }_0\dot{\gamma }\dot{\gamma }_0\right),`$ (2.2)
$`R_{n_i}^{m_i}`$ $`=`$ $`\mathrm{e}^{2\gamma (t)}\left(\ddot{x}^i+(\dot{\gamma }_0\dot{\gamma })\dot{x}^i\right)\delta _{n_i}^{m_i},`$ (2.3)
for $`i=1,\mathrm{},n`$, where the indices $`m_i`$ and $`n_i`$ run from $`D_{j=i}^nN_j`$ to $`D_{j=i}^nN_j+N_i`$ and $`\gamma _0=_{i=1}^nN_ix^i`$.
A viscous fluid is characterized by a density $`\rho `$, a pressure $`p`$, a bulk viscosity coefficient $`\zeta `$, a shear viscosity coefficient $`\eta `$, so that the (standard form of the) energy-momentum tensor reads
$$T_\nu ^\mu =\rho u^\mu u_\nu +(p\zeta \theta )P_\nu ^\mu 2\eta \sigma _\nu ^\mu ,$$
(2.4)
where $`u^\mu `$ is the $`D`$-dimensional velocity of the fluid, $`\theta =u_{;\mu }^\mu `$ denotes the scalar expansion, $`P_\nu ^\mu =\delta _\nu ^\mu +u^\mu u_\nu `$ is the projector on the $`(D1)`$-dimensional space orthogonal to $`u^\mu `$, and $`\sigma _\nu ^\mu =\frac{1}{2}\left(u_{\alpha ;\beta }+u_{\beta ;\alpha }\right)P^{\alpha \mu }P_\nu ^\beta (D1)^1\theta P_\nu ^\mu `$ is the traceless shear tensor (defined as usual).
By choosing the $`D`$-dimensional velocity so that $`u^\mu =\delta _0^\mu \mathrm{e}^{\gamma (t)}`$ (the comoving observer condition), we obtain
$`\theta `$ $`=`$ $`\dot{\gamma }_0\mathrm{e}^{\gamma (t)},`$ (2.5)
$`(u^\mu u_\nu )`$ $`=`$ $`\mathrm{diag}(1,0,\mathrm{},0),`$ (2.6)
$`(P_\nu ^\mu )`$ $`=`$ $`\mathrm{diag}(0,1,\mathrm{},1),`$ (2.7)
$`(\sigma _\nu ^\mu )`$ $`=`$ $`\mathrm{e}^\gamma \mathrm{diag}(0,\left(\dot{x}^1{\displaystyle \frac{\dot{\gamma }_0}{D1}}\right)\delta _{l_1}^{k_1},\mathrm{},\left(\dot{x}^n{\displaystyle \frac{\dot{\gamma }_0}{D1}}\right)\delta _{l_n}^{k_n}),`$ (2.8)
where $`k_i`$, $`l_i=1,\mathrm{},N_i`$ for $`i=1,\mathrm{},n`$. The function $`\gamma (t)`$ determines a time gauge (a harmonic time gauge for $`\gamma (t)=\gamma _0`$ and a synchronous time gauge for $`\gamma (t)=0`$), see Eq. (2.1); note that the harmonic time $`t`$ and the synchronous time $`t_s`$ are related by $`\mathrm{d}t_s=\mathrm{exp}[\gamma _0]\mathrm{d}t`$. By assuming anisotropy properties for the pressure and the bulk viscosity, with respect to the whole space $`M_1\times \mathrm{}\times M_n`$, one has
$$(T_\nu ^\mu )=\mathrm{diag}(\rho ,p_1^{}\delta _{l_1}^{k_1},\mathrm{},p_n^{}\delta _{l_n}^{k_n}),$$
(2.9)
where
$$p_i^{}=p_i\mathrm{e}^\gamma \left[\zeta _i\dot{\gamma }_0+2\eta \left(\dot{x}^i\frac{\dot{\gamma }_0}{D1}\right)\right],$$
(2.10)
and $`p_i`$, resp. $`\zeta _i`$, is the pressure, resp. the bulk viscosity coefficient, in the space described by the manifold $`M_i`$. Furthermore, we assume that the barotropic equations of state hold
$$p_i=(1h_i)\rho (t),$$
(2.11)
where the $`h_i`$ are constants. One easily shows that the form of the equation of motion ($`_MT_0^M=0`$) for a viscous fluid described by a tensor given by Eq. (2.9), is given by
$$\dot{\rho }+\underset{i=1}{\overset{n}{}}N_i\dot{x}^i(\rho +p_i^{})=0.$$
(2.12)
The Einstein equations $`R_\nu ^\mu \frac{1}{2}\delta _\nu ^\mu R=\kappa ^2T_\nu ^\mu `$, where $`\kappa ^2`$ is the gravitational constant, can be written as $`R_\nu ^\mu =\kappa ^2(T_\nu ^\mu \frac{T}{D2}\delta _\nu ^\mu )`$. Further, by using the equations $`R_0^0\frac{1}{2}\delta _0^0R=\kappa ^2T_0^0`$ and $`R_{n_i}^{m_i}=\kappa ^2(T_{n_i}^{m_i}\frac{T}{D2}\delta _{k_i}^{m_i})`$, Eqs.(2.2,2.3,2.9) give the following equations of motion
$$\underset{i=1}{\overset{n}{}}N_i(\dot{x}^i)^2\dot{\gamma }_0^2=2\kappa ^2\mathrm{e}^{2\gamma }\rho ,$$
(2.13)
$`\ddot{x}^i+(\dot{\gamma }_0\dot{\gamma })\dot{x}^i`$ $`=`$ $`\kappa ^2\mathrm{e}^\gamma [\mathrm{e}^\gamma \rho (h_i+{\displaystyle \frac{_{k=1}^nN_kh_k}{D2}})`$ (2.14)
$`+`$ $`\dot{\gamma }_0(\zeta _i+{\displaystyle \frac{_{k=1}^nN_k\zeta _k}{D2}})]2\kappa ^2\mathrm{e}^\gamma \eta (\dot{x}^i{\displaystyle \frac{\dot{\gamma }_0}{D1}}).`$
We use an integration procedure which is based on the $`n`$-dimensional real vector space $`R^n`$. Let $`e_1,\mathrm{},e_n`$ be the canonical basis in $`R^n`$ (i.e. $`e_1=(1,0,\mathrm{},0)`$ etc…), and $`,`$ denote a symmetrical bilinear form defined on $`R^n`$ by
$$e_i,e_j=\delta _{ij}N_jN_iN_jG_{ij}.$$
(2.15)
It has been used as a mini-super-space metric for cosmological models (see \[20-25\]). Such a form is non-generate and has the pseudo-Euclidean signature $`(,+,\mathrm{},+)`$. With this in mind, a vector $`yR^n`$ is time-like, resp. space-like or isotropic, if $`y,y`$ takes negative, resp. positive or null values; and two vectors $`y`$ and $`z`$ are orthogonal if $`y,z=0`$. Hereafter, we use the following vectors
$`x`$ $`=`$ $`x^1e_1+\mathrm{}+x^ne_n,`$ (2.16)
$`u`$ $`=`$ $`u^1e_1+\mathrm{}+u^ne_n,u^i=h_i{\displaystyle \frac{_{k=1}^nN_kh_k}{D2}},u_i=N_ih_i,`$ (2.17)
$`\xi `$ $`=`$ $`\xi ^1e_1+\mathrm{}+\xi ^ne_n,\xi ^i=\zeta _i{\displaystyle \frac{_{k=1}^nN_k\zeta _k}{D2}},\xi _i=N_i\zeta _i,`$ (2.18)
where covariant coordinates of the vectors are introduced by the usual way. Moreover, let us denote $`u_d`$ the particular vector given by Eq. (2.17) with $`h_{i=1,\mathrm{},n}=1`$ (it is related to dust in the whole space, see Eq. (2.11)). One has
$$(u_d)_i=N_i,u_d^i=\frac{1}{D2},u_d,u_d=\frac{D1}{D2},u_d,x=\gamma _0.$$
(2.19)
Thus, using Eqs. (2.16-2.19) we rewrite the Einstein equations (2.13),(2.14) in the form
$`\dot{x},\dot{x}`$ $`=`$ $`2\kappa ^2\mathrm{e}^{2\gamma }\rho ,`$ (2.20)
$`\ddot{x}+\left(u_d,\dot{x}\dot{\gamma }+2\eta \kappa ^2\mathrm{e}^\gamma \right)\dot{x}`$ $`=`$ $`{\displaystyle \frac{\dot{x},\dot{x}}{2}}u`$ (2.21)
$``$ $`\kappa ^2\mathrm{e}^\gamma u_d,\dot{x}\left(\xi {\displaystyle \frac{2\eta }{u_d,u_d}}u_d\right),`$
where the formal dependence on $`\rho `$ in Eq. (2.21) has been canceled, according to Eq. (2.20). Moreover, Eq. (2.12) can be written as
$$\dot{\rho }+2u_du,\dot{x}\rho \mathrm{e}^\gamma \left(2\eta \dot{x},\dot{x}+u_d,\dot{x}\xi \frac{2\eta u_d}{u_d,u_d},\dot{x}\right)=0.$$
(2.22)
To integrate Eq.(2.21) one needs second equations of state, involving the bulk viscosity coefficients $`\zeta _1,\mathrm{},\zeta _n`$ and the shear viscosity coefficient $`\eta `$. Let us assume that these coefficients are proportional to $`\mathrm{exp}[\gamma _0]`$ (or inversely proportional to the volume of the Universe), i.e.
$$\eta ,\zeta _i[\mathrm{scale}\mathrm{factor}\mathrm{of}M_1]^{\mathrm{dim}(\mathrm{M}_1)}\mathrm{}[\mathrm{scale}\mathrm{factor}\mathrm{of}M_n]^{\mathrm{dim}(\mathrm{M}_\mathrm{n})},$$
(2.23)
which means (from a physical viewpoint) that the viscosity decreases when the space $`M_1\times \mathrm{}\times M_n`$ expands. The integrability of the basic equation (provided the second equations of state) is ensured when the vectors $`u,\xi ,u_d`$ are either colinear or orthogonal (with respect to the mini-super-space metric) in some combination . Herein, we suppose that these vectors are colinear, which means that the viscous fluid has identical properties in the internal space(s) and the external space. Hence, all these assumptions, for the pressures and the viscosity coefficients, allow us to write
$`p_i`$ $`=`$ $`(1h)\rho \mathrm{or}u=hu_d,`$ (2.24)
$`\zeta _i`$ $`=`$ $`{\displaystyle \frac{\zeta _0}{\kappa ^2}}\mathrm{e}^{\gamma _0}\zeta \mathrm{or}\xi =\zeta u_d={\displaystyle \frac{\zeta _0}{\kappa ^2}}\mathrm{e}^{\gamma _0}u_d,(i=1,\mathrm{},n),`$ (2.25)
$`\eta `$ $`=`$ $`{\displaystyle \frac{\eta _0}{2\kappa ^2}}\mathrm{e}^{\gamma _0},`$ (2.26)
where $`\zeta _0,\eta _0`$ and $`h`$ are constants.
## 3 Multidimensional Thermodynamics
According to , let us summarize thermodynamics principles in such a multidimensional cosmology. The first law of thermodynamics reads
$$T\mathrm{d}S=\mathrm{d}(\rho V)+V\underset{i=1}{\overset{n}{}}p_i\frac{\mathrm{d}V_i}{V_i},$$
(3.1)
where $`V_i`$ stands for a fluid volume in the space $`M_i`$, when $`V=V_1\mathrm{}V_n`$ is a fluid volume in the whole space, and $`S`$ is an entropy in the volume $`V`$. By assuming that the baryon particle number $`N_\mathrm{B}`$ in the volume $`V`$ is conserved, Eq. (3.1) transforms to
$$nT\dot{s}=\dot{\rho }+\rho \underset{i=1}{\overset{n}{}}N_i\dot{x}^i+\underset{i=1}{\overset{n}{}}p_iN_i\dot{x}^i,$$
(3.2)
where $`s=S/N_\mathrm{B}`$, resp. $`n=N_\mathrm{B}/V`$, stands for the entropy per baryon, resp. the baryon number density. Let us remind that $`\mathrm{exp}[x^i]`$ is the scale factor of the space $`M_i`$ (of dimension $`N_i`$). For a perfect fluid ($`\zeta _i=0`$, $`\eta =0`$), the comparison between Eqs. (2.12,3.2) gives the entropy conservation ( i.e., $`s`$ is constant). Similarly, we obtain also the temperature (see ). From Eq. (3.2) we have
$$\left(\frac{\rho }{x^i}\right)_{s,x^j}=\rho N_ip_iN_i=(h_i2)N_i\rho ,(ji),$$
(3.3)
and then
$$\rho =K(s)\mathrm{exp}\left[\underset{i=1}{\overset{n}{}}(h_i2)N_ix^i\right],$$
(3.4)
where $`K(s)`$ is an unknown function (which reads in term of the entropy per baryon $`s`$). Using Eqs. (3.2),(3.4) we get
$$\left(\frac{\rho }{s}\right)_{x^i}=nT=K^{}(s)\mathrm{exp}\left[\underset{i=1}{\overset{n}{}}(2h_i)N_ix^i\right].$$
(3.5)
For a perfect fluid, we have $`K^{}(s)=1/B`$ where $`B`$ is a constant, then
$$nT=\frac{1}{B}\mathrm{exp}\left[\underset{i=1}{\overset{n}{}}(h_i2)N_ix^i\right]=\frac{1}{B}\mathrm{exp}[u2u_d,x].$$
(3.6)
For a fluid with a bulk and shear viscosity, the comparison between Eqs. (2.12,3.2) provides us with
$$nT\dot{s}=\mathrm{e}^\gamma \left(2\eta \dot{x},\dot{x}+u_d,\dot{x}\xi 2\frac{2\eta }{u_d,u_d}u_d,\dot{x}\right).$$
(3.7)
Such a formula gives the variation rate of entropy per baryon in multidimensional cosmology on the manifold $`M=R\times M_1\times \mathrm{}\times M_n`$ with viscosity. The entropy production can be calculated if the temperature of the fluid is known. Herein, we suppose that the temperature is given by the perfect fluid formula Eq. (3.6), which is valid with sufficient accuracy when effects of viscosity are small. Hence, Eqs. (3.6),(3.7) give
$$\dot{s}=B\mathrm{e}^{2u_du,x\gamma }\left(2\eta \dot{x},\dot{x}+u_d,\dot{x}\xi 2\frac{2\eta u_d}{u_d,u_d},\dot{x}\right).$$
(3.8)
## 4 Exact solutions
According to assumptions given in Eqs. (2.24-2.26), the basic vector equation, see Eq. (2.21), reads
$`\ddot{x}+\left(u_d,\dot{x}\dot{\gamma }+2\eta \kappa ^2\mathrm{e}^\gamma \right)\dot{x}=\left[{\displaystyle \frac{h}{2}}\dot{x},\dot{x}\kappa ^2\mathrm{e}^\gamma u_d,\dot{x}\left(\zeta {\displaystyle \frac{2\eta }{u_d,u_d}}\right)\right]u_d.`$ (4.1)
In order to integrate such a (vector) equation, we use the orthogonal basis
$$\frac{u_d}{u_d,u_d},f_2,\mathrm{},f_nR^n,$$
(4.2)
where the orthogonality property reads
$$u_d,f_j=0,f_j,f_k=\delta _{jk},(j,k=2,\mathrm{},n).$$
(4.3)
Let us note that the basis vectors $`f_2,\mathrm{},f_n`$ are space-like, since they are orthogonal to the time-like vector $`u_d`$. The vector $`xR^n`$ decomposes as follows
$$x=u_d,x\frac{u_d}{u_d,u_d}+\underset{j=2}{\overset{n}{}}f_j,xf_j.$$
(4.4)
Hence, the basic vector equation given in Eq. (4.1) reads in term of coordinates in such a basis as follows
$$f_j,\ddot{x}+\left(u_d,\dot{x}\dot{\gamma }+2\eta \kappa ^2\mathrm{e}^\gamma \right)f_j,\dot{x}=0$$
(4.5)
$$\frac{u_d,\ddot{x}}{u_d,u_d}+\left(\frac{u_d,\dot{x}\dot{\gamma }}{u_d,u_d}+\zeta \kappa ^2\mathrm{e}^\gamma \right)u_d,\dot{x}=\frac{h}{2}\left[\frac{u_d,\dot{x}^2}{u_d,u_d}+\underset{j=2}{\overset{n}{}}f_j,\dot{x}^2\right],$$
(4.6)
which provides us with a set of equations (for $`j=2,\mathrm{},n`$). For the metric dependence of the viscosity coefficients, see Eqs. (2.25),(2.26) Eqs. (4.5,4.6) read
$$f_j,\ddot{x}+\eta _0f_j,\dot{x}=0,(j=2,\mathrm{},n)$$
(4.7)
$$u_d,\ddot{x}\frac{h}{2}u_d,\dot{x}^2+u_d,u_d\left(\zeta _0u_d,\dot{x}\frac{h}{2}\underset{j=2}{\overset{n}{}}f_j,\dot{x}^2\right)=0$$
(4.8)
in the harmonic time gauge
$$\gamma =\gamma _0=u_d,x=\underset{i=1}{\overset{n}{}}N_ix^i.$$
(4.9)
Such a set of equations is integrable for any values of the constant parameters $`h`$, $`\eta _0`$ and $`\zeta _0`$.
Let us first assume models with
$$h0.$$
(4.10)
The integration of Eq. (4.7) gives
$$f_j,x=\{\begin{array}{cc}tp^j+q^j\hfill & \text{if }\eta _0=0\text{,}\hfill \\ \mathrm{e}^{\eta _0t}p^j+q^j\hfill & \text{if }\eta _00\text{,}\hfill \end{array}$$
(4.11)
where $`p^j`$ and $`q^j`$ are arbitrary constants. The Kasner-like form solution can be written in term of vectors $`\alpha ,\beta R^n`$, defined as follows
$$\alpha =\underset{j=2}{\overset{n}{}}p^jf_j\underset{i=1}{\overset{n}{}}\alpha ^ie_i,\beta =\underset{j=2}{\overset{n}{}}q^jf_j\underset{i=1}{\overset{n}{}}\beta ^ie_i,$$
(4.12)
where $`\alpha ^i`$ and $`\beta ^i`$ are their coordinates in the canonical basis $`e_1,\mathrm{},e_n`$. By using the orthogonality conditions, we obtain
$$\alpha ,u_d=\underset{i=1}{\overset{n}{}}\alpha ^iN_i=0,\beta ,u_d=\underset{i=1}{\overset{n}{}}\beta ^iN_i=0,$$
(4.13)
$$\alpha ,\alpha =\underset{i=1}{\overset{n}{}}\left(\alpha ^i\right)^2N_i=\underset{j=2}{\overset{n}{}}\left(p^j\right)^20$$
(4.14)
$$\beta ,\beta =\underset{i=1}{\overset{n}{}}\left(\beta ^i\right)^2N_i=\underset{j=2}{\overset{n}{}}\left(q^j\right)^20,$$
(4.15)
where the constants $`\alpha ^i`$ and $`\beta ^i`$ may be called Kasner-like parameters, because of the existence of these constraints. By using Eqs. (4.4,4.11,4.12), we obtain
$$x=u_d,x\frac{u_d}{u_d,u_d}+a(t)\alpha +\beta ,$$
(4.16)
where the function
$$a(t)=\{\begin{array}{cc}t\hfill & \text{if }\eta _0=0\text{,}\hfill \\ \mathrm{e}^{\eta _0t}\hfill & \text{if }\eta _00\text{,}\hfill \end{array}$$
(4.17)
By substituting the functions $`f_j,x`$ into Eq. (4.8) we obtain the following equation for the unknown function $`u_d,x`$
$$u_d,\ddot{x}\frac{h}{2}u_d,\dot{x}^2+u_d,u_d\zeta _0u_d,\dot{x}=\frac{h}{2}A^2\dot{a}^2(t),$$
(4.18)
where
$$A=\sqrt{\frac{D1}{D2}\alpha ,\alpha }.$$
(4.19)
The equation Eq. (4.18) has been integrated for the non-viscous model $`\eta _0=\zeta _0=0`$ (see, e.g., ) and for the model with $`\zeta _00`$ and $`\eta _0=0`$ in . For the model with $`\eta _00`$, it can be reduced to the modified Bessel equation
$$\tau ^2\frac{\mathrm{d}^2z}{\mathrm{d}\tau ^2}+\tau \frac{\mathrm{d}z}{\mathrm{d}\tau }(\tau ^2+\nu ^2)z=0$$
(4.20)
by means of the transformation
$$\tau =\frac{1}{2}|h|A\mathrm{e}^{\eta _0t},$$
(4.21)
$$u_d,\dot{x}=\frac{2\eta _0}{h}\tau \frac{\mathrm{d}}{\mathrm{d}\tau }\mathrm{ln}|\tau ^\nu z(\tau )|$$
(4.22)
in the non-trivial case $`\alpha ,\alpha 0`$, where the constant
$$\nu =\frac{D1}{D2}\frac{\zeta _0}{2\eta _0}.$$
(4.23)
The general solution of the modified Bessel equation is given by
$$z(\tau )=C_1I_{|\nu |}(\tau )+C_2K_{|\nu |}(\tau ),$$
(4.24)
see , where $`I_{|\nu |}(\tau )`$, resp. $`K_{|\nu |}(\tau )`$, is the related modified Bessel function, resp. Mac-Donald function.
Finally, the results of Eqs. (4.7,4.8) integration for various constants $`\zeta _0`$ and $`\eta _0`$ can be presented in Tab. 1
Table 1
| | $`\zeta _0=0`$ | $`\zeta _00`$ |
| --- | --- | --- |
| $`\eta _0=0`$ | Solution II with $`\alpha ,\alpha 0`$ and Solution I | Solution II |
| $`\eta _00`$ | Solution III with $`\alpha ,\alpha 0`$ and Solution I | Solution III with $`\alpha ,\alpha 0`$ and Solution II with $`\alpha ,\alpha =0`$ |
The solutions I,II,III in term of scale factors are the following
* Solution I :
$`\mathrm{e}^{x^i}`$ $`=`$ $`\mathrm{e}^{\beta ^i}\left|C_1+C_2t\right|^{2/[h(D1)]}.`$ (4.25)
* Solution II :
$`\mathrm{e}^{x^i}`$ $`=`$ $`\mathrm{e}^{\alpha ^it+\beta ^i}\left|C_1\mathrm{e}^{(\stackrel{~}{A}\stackrel{~}{\zeta }_0)ht/2}+C_2\mathrm{e}^{(\stackrel{~}{A}+\stackrel{~}{\zeta }_0)ht/2}\right|^{2/[h(D1)]},`$ (4.26)
where
$$\stackrel{~}{\zeta }_0=\frac{D1}{D2}\frac{\zeta _0}{h},\stackrel{~}{A}=\sqrt{\frac{D1}{D2}\alpha ,\alpha +\left(\frac{D1}{D2}\frac{\zeta _0}{h}\right)^2}=\sqrt{A^2+\stackrel{~}{\zeta }_0^2}.$$
(4.27)
* Solution III :
$`\mathrm{e}^{x^i}`$ $`=`$ $`\mathrm{exp}\left(\alpha ^i\mathrm{e}^{\eta _0t}+\beta ^i\right)\left(\tau ^\nu \left|C_1I_{|\nu |}(\tau )+C_2K_{|\nu |}(\tau )\right|\right)^{2/[h(D1)]},`$ (4.28)
where the variable $`\tau >0`$ is given in Eq. (4.21), the constants $`A`$ in Eq. (4.19) and $`\nu `$ in Eq. (4.23).
In formulas (4.25),(4.27),(4.28) $`C_{i=1,2}`$ are integration constants such that $`C_1^2+C_2^2>0`$ and $`i=1,\mathrm{},n`$. The Kasner-like parameters $`\alpha ^i`$ and $`\beta ^i`$ obey the relations given in Eq. (4.13).
The set of equations given in Eqs. (4.7,4.8) is easily integrable in the case
$$h=0,$$
(4.29)
which, with the barotropic equation of state in mind, relates to Zeldovich or stiff matter. The results are given as follows :
* Solution IV : for $`i=1,\mathrm{},n`$
$`\mathrm{e}^{x^i}`$ $`=`$ $`\mathrm{exp}[\alpha ^ia(t)+\beta ^i]\times \{\begin{array}{cc}\mathrm{exp}(C_1+C_2t)\hfill & \text{if }\zeta _0=0\hfill \\ \mathrm{exp}(C_1+C_2\mathrm{exp}(\zeta _0\frac{D1}{D2}t\left)\right)\hfill & \text{if }\zeta _00\text{,}\hfill \end{array}`$ (4.32)
where $`C_{i=1,2}`$ are arbitrary constants, and the function $`a(t)`$ is given in Eq. (4.17).
## 5 Discussion
Let us remind the multidimensional generalization of the well-known Kasner solution , it reads (for the synchronous time $`t_s`$) as follows
$$\mathrm{d}s^2=\mathrm{d}t_s^2+\underset{i=1}{\overset{n}{}}A_it_s^{2\epsilon ^i}\mathrm{d}s_i^2.$$
(5.1)
Such a metric describes the evolution of a vacuum model defined on the manifold $`R\times M_1\times \mathrm{}\times M_n`$, where the $`M_i`$ are Ricci-flat spaces of dimension $`N_i`$ with the metric $`\mathrm{d}s_i^2`$, $`A_i`$ are arbitrary constants and $`\epsilon _i`$ are the Kasner parameters, which satisfy the relations
$$\underset{i=1}{\overset{n}{}}N_i\epsilon ^i=1,\underset{i=1}{\overset{n}{}}N_i(\epsilon ^i)^2=1.$$
(5.2)
The $`\epsilon ^i`$ and the Kasner-like parameters $`\alpha ^i`$ (used in the above formulas for the exact solutions) are related as
$$\epsilon ^i=\pm \frac{\alpha ^i}{A}+\frac{1}{D1},\alpha ,\alpha 0.$$
(5.3)
By using Eq. (4.16) (i.e., a general decomposition of the vector $`xR^n`$) and the result of Sec.4, see Eq. (3.8), we obtain the variation rate of entropy
$$\dot{s}=\frac{B}{\kappa ^2}\mathrm{e}^{h\gamma _0}\left(\zeta _0\dot{\gamma }_0^2+\eta _0\alpha ,\alpha \dot{a}^2(t)\right),$$
(5.4)
which shows that the entropy increases when $`\zeta _0dt>0`$ and $`\eta _0dt>0`$. Further we assume
$$\zeta _00,\eta _00,$$
(5.5)
so harmonic time $`t`$ (as well as synchronous time $`t_s`$) increases during the evolution.
One easily shows that the weak energy condition $`T_\nu ^\mu v^\nu v_\mu 0`$, for any D-dimensional non space-like vector $`v^\nu `$ (and thus as well as for the 4-dimensional case), applied to the stress-energy tensor given in Eq. (2.9), can be written as inequalities
$$\rho 0,\rho +p_i^{}0,(i=1,\mathrm{},n),$$
(5.6)
where $`\rho `$, resp. $`p_i^{}`$, is the density, resp. the effective pressure, of the fluid.
The dominant energy condition $`T_\nu ^\mu v^\nu v_\mu 0`$ and $`T_\mu ^\nu T_\lambda ^\mu v_\nu v^\lambda 0`$, for any non space-like vector $`v^\mu `$, applied to the stress-energy tensor given in Eq. (2.9), reduces to inequalities defined in Eq. (5.6), with the following additional condition
$$\rho p_i^{}0,(i=1,\mathrm{},n)$$
(5.7)
Notice that due to the weak energy condition the following restriction on the constant $`h`$ (taken from the barotropic equation of state (2.11)) arises in the nonviscous case: $`h2`$. The dominant energy condition for the nonviscous stress-energy tensor implies: $`0h2`$.
It is important to note that Solution I and Solutions II,IV for zero Kasner-like parameters ($`\alpha ^i=0,i=1,\mathrm{},n`$) are isotropic, since the spaces $`M_1,\mathrm{},M_n`$ have identical scale factors.
Further we discuss the solutions obtained, which are of interest within multidimensional or $`4`$-dimensional cosmology.
### 5.1 Non Viscous Models
For a better understanding of the viscosity effect on the dynamics, we first outline the properties of non viscous models.
#### 5.1.1
The isotropic non viscous model is described by Solution I ($`h0`$) and Solution IV ($`h=0`$) with $`\alpha ^i=0`$ (for $`i=1,\mathrm{},n`$), which represents the multidimensional generalization of the flat FRW model. It is the steady-state model
$$\mathrm{e}^{x^i}\mathrm{exp}\left[\frac{C_2}{D1}t_s\right],\rho =\mathrm{const}$$
(5.8)
for $`h=2`$ and shows a power-law behavior
$$\mathrm{e}^{x^i}t_s^{2/[(2h)(D1)]},\rho t_s^2$$
(5.9)
for $`h2`$, where $`t_s`$ is the synchronous time (stationary solution with zero density is also possible); let us call it as Friedman-like behavior.
#### 5.1.2
The anisotropic non viscous model for $`h(0,2]`$ is described by Solution II with $`\zeta _0=0`$ and $`\alpha ,\alpha 0`$. One easily shows that the integration constants $`C_{i=1,2}`$ have to satisfy the condition $`C_1C_2<0`$, otherwise $`\rho 0,t`$. By using this condition, we obtain from Eq. (4.26)
$$\mathrm{e}^{\gamma _0}\left|\mathrm{sinh}[Ah(tt_0)/2]\right|^{2/h}$$
(5.10)
where $`t_0=0`$ can be chosen (with no loss of generality). Then, for a suitable choice of the integration constant of equation $`\mathrm{d}t_s=\mathrm{exp}(\gamma _0)\mathrm{d}t`$, one has the following correspondences $`t(\mathrm{},0)t_s(0,+\mathrm{})`$ and $`t(0,+\mathrm{})t_s(\mathrm{},0)`$. Hence, we solely investigate the solution $`t_s(\mathrm{},0)`$, since the evolution of the non-viscous fluid is reversible.
From Eq. (5.10), we obtain in the main order $`\mathrm{exp}[\gamma _0]t^{2/h}`$ when $`t+0`$ ($`t_s+\mathrm{}`$), then one has $`t_st^{12/h}`$ for $`h(0,2)`$ and $`t_s\mathrm{ln}t`$ for $`h=2`$. By using these relations and Eq. (4.26), we easily see that the multidimensional Universe shows an isotropical Friedman-like contraction, as defined by Eq. (5.9), in the (infinite) past. Such a conclusion is also valid for Zeldovich matter ($`h=0`$).
Let us now investigate the behavior of the non-viscous anisotropic model at $`t_s=0`$. For Solution II ($`h0`$), by using Eq. (5.10), we obtain in the main order $`\mathrm{exp}[\gamma _0]\mathrm{exp}[At]`$ when $`t+\mathrm{}`$ ($`t_s+0`$), then $`t_s\mathrm{exp}[At]`$. By substituting the latter relation into Eqs. (4.26), we obtain in the main order
$$\mathrm{e}^{x^i}|t_s|^{\alpha ^i/A+1/(D1)},\rho |t_s|^{h2}(t_s+0).$$
(5.11)
According to Eq. (5.3), the model for $`h(0,2]`$ has a Kasner-like behavior near the singularity (at $`t_s=0`$). Such a behavior describes the contraction of some spaces $`M_1,\mathrm{},M_n`$ and the expansion for the other ones. According to Eq. (5.2), the number of either contracting or expanding spaces depends on $`n`$ (the total number of spaces) and $`N_{i=1,n}`$ (their dimensions), but there is at least a contracting manifold and expanding one. Such dynamics is a mechanism of extra dimensions compactification (within the multidimensional cosmology).
One can easily show that the non-viscous anisotropic model for Zeldovich matter ($`h=0`$) described by Solution IV has almost the same properties but with the second constraint given Eq. (5.2) for the Kasner parameters substituted by $`_{i=1}^nN_i(\epsilon ^i)^2=\epsilon `$, where $`\epsilon `$ is a constant such that $`1/(D1)<\epsilon <1`$.
### 5.2 Models with bulk viscosity
Let us now investigate the models with bulk viscosity.
#### 5.2.1
The viscous isotropic model shows interesting features. The shear viscosity is not significant in this case and the model is described for $`h0`$ by Solution II with $`\alpha ^i=0`$ for $`i=1,\mathrm{},n`$. If $`C_1C_2<0`$ and $`h>0`$ then the solution can be written as follows : for $`i=1,\mathrm{},n`$
$`\mathrm{e}^{x^i}`$ $`=`$ $`R_i\left(1\mathrm{exp}\left[{\displaystyle \frac{D1}{D2}}\zeta _0(tt_0)\right]\right)^{2/[h(D1)]},`$ (5.12)
$`p_i^{}`$ $`=`$ $`\left(1h\mathrm{exp}\left[{\displaystyle \frac{D1}{D2}}\zeta _0(tt_0)\right]\right)\rho ,`$ (5.13)
$`\rho `$ $`=`$ $`{\displaystyle \frac{D1}{D2}}{\displaystyle \frac{2\zeta _0^2}{\kappa ^2h^2}}{\displaystyle \underset{i=1}{\overset{n}{}}}R_i^{2N_i}\mathrm{exp}\left[2{\displaystyle \frac{D1}{D2}}\zeta _0(tt_0)\right]`$ (5.14)
$`\times \left(1\mathrm{exp}\left[{\displaystyle \frac{D1}{D2}}\zeta _0(tt_0)\right]\right)^{2(2h)/h]},`$
$`s`$ $`=`$ $`{\displaystyle \frac{D1}{D2}}{\displaystyle \frac{2B\zeta _0^2}{\kappa ^2h^2}}{\displaystyle \underset{i=1}{\overset{n}{}}}R_i^{hN_i}\mathrm{exp}\left[2{\displaystyle \frac{D1}{D2}}\zeta _0(tt_0)\right]+s(\mathrm{}),`$ (5.15)
where we use the set of independent constants $`t_0,R_1,\mathrm{},R_n`$, defined such that
$$\frac{C_1}{C_2}=\mathrm{exp}\left[\frac{D1}{D2}\zeta _0t_0\right],|C_1|^{2/[h(D1)]}\mathrm{exp}[\beta ^i]=R_i.$$
(5.16)
Let us consider this solution on the interval $`(\mathrm{},t_0)`$ for $`h(0,2)`$. Then the synchronous time $`t_s`$ changes during the evolution on the interval $`(\mathrm{},+\mathrm{})`$. Such a solution is non-singular and describes a monotonic isotropic expansion of the $`D`$-dimensional Universe with Friedman-like stage as $`t_s+\mathrm{}`$. The density $`\rho `$ increases from zero in the infinite past to some maximum value and then decreases to zero at the Friedman-like stage. The maximum of the density is reached at $`tt_m=(D2)\mathrm{ln}[h/2]/[\zeta _0(D1)]+t_0`$, when $`p_i^{}=\rho `$, see Eq. (5.13). We have $`p_i^{}+\rho <0`$ for $`t(\mathrm{},t_m)`$ and $`p_i^{}+\rho >0`$ for $`t(t_m,t_0)`$, so the weak energy condition, see Eq. (5.6), is not satisfied on the time interval $`(\mathrm{},t_m)`$. The entropy per baryon monotonically increases during the evolution and tends to some constant value in the infinite future.
The nonsingular solution obtained by Murphy within flat FRW model with bulk viscosity for another second equation of state: $`\zeta =\mathrm{const}\rho `$ exhibits the similar properties except the violating of the weak energy condition. The scale factor monotonically increases from zero in the infinite past and tends to the infinity at the Friedman stage of the evolution. The density monotonically decreases from some constant value in the infinite past and tends to zero according to Eq. (5.9) in the infinite future. The weak energy condition is valid for $`t_s(\mathrm{},+\mathrm{})`$ (but the strong energy condition $`\rho +3p^{}0`$ is not satisfied on the interval $`(\mathrm{},t_s^{})`$, where $`t_s^{}`$ is some constant), so from this point of view the Murphy solution is more attractive.
#### 5.2.2
Now, let us study the properties of the anisotropic viscous model by taking into account only the bulk viscosity. Such model is described by Solution II with non-zero Kasner-like parameters $`\alpha ^i`$ and has been previously studied . If $`C_1C_2>0`$ in Solution II with $`\alpha ^i0`$ then the density has negative values at some stage of the evolution, which means that the weak energy condition is not satisfied. Hence, hereafter such solutions are not considered. If $`C_1C_2<0`$, then Solution II can be written as follows : for $`i=1,\mathrm{},n`$
$`\mathrm{e}^{x^i}`$ $`=`$ $`R_i\left|\mathrm{sinh}[\stackrel{~}{A}h(tt_0)/2]\right|^{2/[h(D1)]}\mathrm{exp}\left[\left(\alpha ^i{\displaystyle \frac{\zeta _0}{h(D2)}}\right)t\right],`$ (5.17)
$`p_i^{}`$ $`=`$ $`\left(1{\displaystyle \frac{h\mathrm{cosh}^2\stackrel{~}{a}}{1+\mathrm{sinh}\stackrel{~}{a}\mathrm{sinh}[\stackrel{~}{A}h(tt_0)+\stackrel{~}{a}]}}\right)\rho ,`$ (5.18)
$`\rho `$ $`=`$ $`{\displaystyle \frac{\alpha ,\alpha }{2\kappa ^2}}{\displaystyle \underset{i=1}{\overset{n}{}}}R_i^{2N_i}\left|\mathrm{sinh}[\stackrel{~}{A}h(tt_0)/2]\right|^{2(2h)/h}\mathrm{exp}[2\stackrel{~}{\zeta }_0t]`$ (5.19)
$`\times \left(1+\mathrm{sinh}\stackrel{~}{a}\mathrm{sinh}[\stackrel{~}{A}h(tt_0)+\stackrel{~}{a}]\right),`$
$`\dot{s}`$ $`=`$ $`{\displaystyle \frac{B\zeta _0A^2}{\kappa ^2}}{\displaystyle \underset{i=1}{\overset{n}{}}}R_i^{hN_i}\mathrm{exp}\left[{\displaystyle \frac{D1}{D2}}\zeta _0\left(1{\displaystyle \frac{2}{h}}\right)t\right]\mathrm{sinh}^2[\stackrel{~}{A}h(tt_0)]`$ (5.20)
$`\times \mathrm{sinh}^2[\stackrel{~}{A}h(tt_0)+\stackrel{~}{a}],`$
where we use independent constants $`t_0,R_1,\mathrm{},R_n`$ defined by
$`C_1={\displaystyle \frac{C}{2}}\mathrm{exp}[\stackrel{~}{A}ht_0/2],C_2={\displaystyle \frac{C}{2}}\mathrm{exp}[\stackrel{~}{A}ht_0/2],`$ (5.21)
$`R_i=|C|^{2/[h(D1)]}\mathrm{exp}[\beta ^i],`$ (5.22)
and $`\stackrel{~}{a}`$ is defined by
$$\mathrm{sinh}\stackrel{~}{a}=\frac{\stackrel{~}{\zeta }_0}{A},\mathrm{cosh}\stackrel{~}{a}=\frac{\stackrel{~}{A}}{A}.$$
(5.23)
Let us note that for $`t(\mathrm{},t_0)`$ the density, given by Eq. (5.19), has negative values. In the following, this solution is only used on the interval $`(t_0,+\mathrm{})`$.
If $`h(0,2)`$ then the harmonic time interval $`t:(t_0,+\mathrm{})`$ corresponds to the following synchronous time interval $`t_s:(\mathrm{},t_s^0)`$, where we choose the integration constant $`t_s^0=0`$. We can easily prove that such a model has the stage of isotropic contraction by the Friedman-like law defined in Eq. (5.9) in the infinite past. Near the final point of the evolution $`t_s=0`$ we obtain in the main order
$$\mathrm{e}^{x^i}t_s^{\alpha ^i/(\stackrel{~}{A}+\stackrel{~}{\zeta }_0)+1/(D1)},i=1,\mathrm{},n.$$
(5.24)
The final point of the evolution $`t_s=0`$ is singular ($`\rho +\mathrm{}`$ as $`t_s0`$), and the characteristic of the singularity depends on the parameter $`\stackrel{~}{\zeta }_0/A`$ (which determines the ratio between the viscosity parameter $`\stackrel{~}{\zeta }_0`$ and the anisotropy parameter $`A`$). If the ratio $`\stackrel{~}{\zeta }_0/A1`$ then we obtain from Eq. (5.24) $`\mathrm{exp}[x^i]t_s^{1/(D1)}`$ near $`t_s=0`$. In such a case, the model describes the contraction of all spaces $`M_1,\mathrm{},M_n`$ in the vicinity of the singularity. If $`\stackrel{~}{\zeta }_0/A1`$ then we obtain ¿from Eq. (5.24) $`\mathrm{exp}[x^i]t_s^{\epsilon ^i}`$ (i.e. the singularity is of the Kasner type). According to Eq. (5.20), in both cases the model describes the unbounded production of the entropy at the final stage of the evolution.
Therefore, anisotropic Solution II for $`h(0,2)`$ and $`C_1C_2<0`$ describes the model with Friedman-like isotropic contraction in the infinite past and the anisotropic Kasner-like behavior near the final singularity if the parameter $`\stackrel{~}{\zeta }_0/A`$ is small enough. Under such a condition, the behavior of scale factors remains (qualitatively) the same as for the anisotropic non-viscous model for $`t_s<0`$. Let us note that this model satisfies the dominant energy condition during the evolution. According to Eq. (5.18), the ratio $`p_i^{}/\rho `$ increases monotonically from the value $`(1h)`$ at the Friedman-like stage in the infinite past, and tends to $`1`$ in the vicinity of the final singularity. One may also consider Solution II for $`h>0`$ and $`C_1=0`$, $`C_20`$. This partial solution has the similar dynamical behavior at the final stage of the evolution and satisfies the dominant energy condition as $`p_i^{}=\rho `$, $`t_s`$.
If the dominant energy condition is ignored then we may also consider anisotropic Solution II for $`h<0`$ and $`C_1C_2<0`$. In such a case the harmonic time interval $`t:(t_0,+\mathrm{})`$ corresponds to the synchronous time interval $`t_s:(0,+\mathrm{})`$ . We easily show that such a model has the Friedman-like singularity, defined in Eq. (5.9) at $`t_s=0`$. The anisotropic behavior is possible far from the singularity, if the parameter $`|\stackrel{~}{\zeta }_0|/A`$ is small enough. This solution satisfies the weak energy condition as the ratio $`p_i^{}/\rho `$ monotonically decreases during the evolution from the value $`(1h)`$ at the Friedman-like stage, and tends to $`1`$ in the infinite future. Let us note that the asymptotical behavior of this solution far from the singularity ($`t_s+\mathrm{}`$) is given by the Solution II for $`h<0`$ and $`C_2=0`$, $`C_10`$. To investigate the possible anisotropic behavior far from the singularity let us consider this partial solution. It may be written in the synchronous time as following : for $`i=1,\mathrm{},n`$
$`\mathrm{e}^{x^i}`$ $`=`$ $`R_it_s^{\alpha ^i/(\stackrel{~}{A}+|\stackrel{~}{\zeta }_0|)+1/(D1)},t_s>0,`$ (5.25)
$`p_i^{}=\rho `$ $`=`$ $`{\displaystyle \frac{D2}{D1}}{\displaystyle \frac{|\stackrel{~}{\zeta }_0|}{\kappa ^2(\stackrel{~}{A}+|\stackrel{~}{\zeta }_0|)}}t_s^2,`$ (5.26)
$`s`$ $`=`$ $`{\displaystyle \frac{(D2)B|\zeta _0|(\stackrel{~}{A}+|\stackrel{~}{\zeta }_0|)}{(D1)\kappa ^2}}t_s^h+s(0),`$ (5.27)
where the constants $`R_1,\mathrm{},R_n`$ are such that
$$\underset{i=1}{\overset{n}{}}R_i^{N_i}=\stackrel{~}{A}+|\stackrel{~}{\zeta }_0|.$$
(5.28)
It is evident from Eq. (5.25) that $`\alpha ^i/(\stackrel{~}{A}+|\stackrel{~}{\zeta }_0|)+1/(D1)\epsilon ^i`$ (Kasner parameter) as $`|\stackrel{~}{\zeta }_0|/A0`$. Therefore, if the parameter $`|\stackrel{~}{\zeta }_0|/A`$ is small enough, then the model describes the contraction of a part of the spaces $`M_1,\mathrm{},M_n`$ and the expansion of another part. According to Eq. (5.27), the solution describes the unbounded entropy production as $`h<0`$.
We note, that the anisotropic viscous model described by Solution IV ($`h=0`$) does not satisfy the weak energy condition as the density has negative values at some stage of the evolution.
### 5.3 Models with bulk and shear viscosity
Now let us investigate the model with both bulk and shear viscosity. Such a model for $`h0`$ is described by Solution III with non-zero Kasner like parameters $`\alpha ^i`$ and $`\nu >0`$. By using the following properties of the modified Bessel functions
$$I_{\nu +1}(\tau )<I_\nu (\tau ),K_{\nu +1}(\tau )>K_\nu (\tau )\tau (0,+\mathrm{}),\nu 0,$$
(5.29)
it can be proved that the density given has negative values at some interval of time if $`C_2=0`$. If $`C_20`$ then the Solution III can be written as follows: for $`i=1,\mathrm{},n`$
$`\mathrm{e}^{x^i}`$ $`=`$ $`R_i\left(\tau ^\nu \left|CI_\nu (\tau )+K_\nu (\tau )\right|\right)^{2/[h(D1)]}\mathrm{exp}\left[\alpha ^i\mathrm{e}^{\eta _0t}\right],`$ (5.30)
$`\dot{x}^i`$ $`=`$ $`{\displaystyle \frac{2\eta _0\tau }{h}}\left[\left({\displaystyle \frac{CI_{\nu +1}(\tau )K_{\nu +1}(\tau )}{CI_\nu (\tau )+K_\nu (\tau )}}+1\right){\displaystyle \frac{1}{D1}}\epsilon ^i\right],`$ (5.31)
$`p_i^{}`$ $`=`$ $`\left(1h+{\displaystyle \frac{h}{\tau }}{\displaystyle \frac{(D1)\epsilon ^i1\nu \frac{CI_{\nu +1}(\tau )K_{\nu +1}(\tau )}{CI_\nu (\tau )+K_\nu (\tau )}}{\left(\frac{CI_{\nu +1}(\tau )K_{\nu +1}(\tau )}{CI_\nu (\tau )+K_\nu (\tau )}\right)^21}}\right)\rho ,`$ (5.32)
$`\rho `$ $`=`$ $`{\displaystyle \frac{2\eta _0^2}{h^2\kappa ^2}}{\displaystyle \frac{D2}{D1}}{\displaystyle \underset{i=1}{\overset{n}{}}}R_i^{2N_i}\tau ^{2(12\frac{\nu }{h})}\left|CI_\nu (\tau )+K_\nu (\tau )\right|^{\frac{4}{h}}`$ (5.33)
$`\times \left[\left({\displaystyle \frac{CI_{\nu +1}(\tau )K_{\nu +1}(\tau )}{CI_\nu (\tau )+K_\nu (\tau )}}\right)^21\right],`$
$`\dot{s}`$ $`=`$ $`{\displaystyle \frac{4\eta _0^3B}{h^2\kappa ^2}}{\displaystyle \frac{D1}{D2}}{\displaystyle \underset{i=1}{\overset{n}{}}}R_i^{hN_i}\tau ^{2(1\nu )}`$ (5.34)
$`\times \left(\nu [CI_{\nu +1}(\tau )K_{\nu +1}(\tau )]^2+[CI_\nu (\tau )+K_\nu (\tau )]^2\right),`$
where the independent constants $`R_1,\mathrm{},R_n,C=C_1/C_2`$ are used. The variable $`\tau `$ was determined in Eq. (4.21). Also we introduced Kasner parameters $`\epsilon ^i`$ by
$$\epsilon ^i=\mathrm{sgn}[h]\frac{\alpha ^i}{A}+\frac{1}{D1}.$$
(5.35)
By using the properties given in Eq. (4.29) and the asymptotical behavior of the modified Bessel functions , one may prove that the density given in Eq. (5.33) has no negative values during the evolution only for :
1. partial solution with $`C=0`$ on the interval $`(\mathrm{},+\mathrm{})`$ of the harmonic time;
2. solution with $`C<0`$ on the interval $`(t_0(\tau _0),+\mathrm{})`$, where $`\tau _0`$ is the root of the equation
$$CI_\nu (\tau _0)+K_\nu (\tau _0)=0.$$
(5.36)
We note that the asymptotical behavior of the solution with $`C<0`$ as $`t+\mathrm{}`$ is given by the partial solution with $`C=0`$.
Let us now investigate the models having a singularity at the beginning of the evolution. Such solutions arise when $`h<0`$. In this case for the solutions with $`C<0`$ one obtain the following correspondences $`(\tau )(\tau _0,0)t(t_0(\tau _o),+\mathrm{})t_s(0,+\mathrm{})`$. The singularity at $`t_s=0`$ is of Friedman type. The shear viscosity leads to the isotropization at the final stage of evolution. We obtain from Eq. (5.30-5.33) $`\mathrm{exp}[x^i]t_s^{1/(D1)}`$ and $`p_i^{}=\rho t_s^2`$ as $`t_s+\mathrm{}`$, i.e. the model describes the isotropic Friedman-like expansion corresponding to Zeldovich matter ($`h=0`$). The anisotropic behavior is possible on some interval of time after the Friedman-like behavior. One can prove that if the constant $`|C|0`$ is small enough then on some interval the sign of the Hubble parameter $`\dot{x}^i`$ coincides with the sign of the Kasner parameter $`\epsilon ^i`$ (see Eq. (5.31) for $`h<0`$). However, such regime is possible only on limited interval of time, because the final stage of the evolution exhibits the isotropic expansion due to the shear viscosity.
Acknowledgments
This work was supported in part by the Russian State Committee for Science and Technology, the Russian Fund of Basic Sciences and by Université de Provence (for V.N.M.).
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# Scaling behavior of developing and decaying networks
## Abstract
We find that a wide class of developing and decaying networks has scaling properties similar to those that were recently observed by Barabási and Albert in the particular case of growing networks. The networks considered here evolve according to the following rules: (i) Each instant a new site is added, the probability of its connection to old sites is proportional to their connectivities. (ii) In addition, (a) new links between some old sites appear with probability proportional to the product of their connectivities or (b) some links between old sites are removed with equal probability.
PACS numbers: 05.10.-a, 05.40.-a, 05.50.+q, 87.18.Sn
Recently observed scaling behaviour of a number of networks $`[114]`$ again sharply increased the interest of these exciting objects widely studied for a long time $`[1520]`$. In fact the scaling properties were found only in a few of a great number of growing networks \[2, 4, 7, 22, dmj00\] but one of them – Web – is so significant for everybody that the topic turned to be really hot.
The simplest model of a scale-free growing network was proposed by Barabási and Albert . In this model, each new site is connected with some old site with probability proportional to its connectivity $`k`$, i.e. to the number of connections with this site. Then the distribution of the connectivities in the large network has a power-law dependence $`P(k)k^\gamma `$ with the exponent $`\gamma =3`$ . In fact, such a network is self-organized into a scale-free structure. A full form of the distribution of the connectivities and some other related properties of the model were calculated exactly . Introduction of aging of sites proportional to $`\tau ^\alpha `$, where $`\tau `$ is the age of a site, does not change scaling properties crucially for $`\alpha <1`$, but scaling breaks at higher values of the aging exponent . Several examples of real networks with aging of sites are described in . The simplicity of the Barabási–Albert’s (BA) model makes it a convenient object to study evolution of networks.
Nevertheless, the BA model describes only a particular type of evolving networks. Of course, reality is much reacher. In real networks, e.g. in Internet, links are not only added but may break from time to time. That certainly changes the structure of such networks. Note also, that new links between old sites may appear in a different way than links between new and old sites, according to different rules. Therefore, it is tempting to find out whether the observed behaviour is usual for a vast variety of networks or applies only to a very restricted number of invented objects. With that purpose, we extend a set of models of evolving networks starting from the BA model.
In the present letter, we propose models of developing and decaying networks with undirected links which show scaling behaviour. We consider structures which evolve due to the following reasons. First, they grow like in the BA model, i.e. in each instant one new site is added and is connected with an old site by an undirected link with a probability proportional to its connectivity $`k`$ (one may check that the general results – the existence of the scaling and the values of the scaling exponents – do not depend on the number of the connections with a new site). In addition, we introduce a new parallel component of the evolution – the permanent addition of new undirected links between old sites or, on the contrary, the permanent removal of some old links. We consider two different cases. (a) A developing network: Each instant, new $`c`$ links are added between unconnected pairs of old sites $`i`$ and $`j`$ with probability proportional to the product of their connectivities $`k_ik_j`$. ($`c`$ may be also non integer. For that one can introduce probability of addition of a link.) $`c0`$. (b) A decaying structure (in fact, it is a set of clusters): Each instant, some links between old sites are removed with equal probability. In this case $`c0`$.
Note that both processes – the addition of new sites with new links and the addition of new links between old sites (or removing old links) proceed in parallel, so the resulting structures differ from the original BA model all the time.
We study the following one-site quantities of the structures: the total distribution of connectivities at long times, $`P(k)`$, and the average connectivity of a site of an age $`s`$ at long time $`t`$, $`\overline{k}(s,t)`$, and their scaling exponents $`P(k)k^\gamma `$ and $`\overline{k}(s,t)(s/t)^\beta `$.
Below, we demonstrate both analytically and by simulation that the introduced evolving networks show scaling behaviour in a wide range of values of $`c`$. Nevertheless, while both $`P(k)`$ and $`\overline{k}(s,t)`$, for the developing networks, are power-law functions for all $`c0`$, only $`\overline{k}(s,t)`$ demonstrates the power-law behaviour in the whole range $`1<c<0`$ for the decaying structures. In this case, the power-law dependence of the distribution $`P(k)`$ is observed only close to $`c=0`$.
In order to study scaling properties of the evolving networks we performed numerical simulations according to the introduced above rules. Each instant, we add one new site with one link and, in addition, may remove some of the old links (decaying network) or, on the contrary, may add some new links between unconnected directly old sites (developing network) with the relative rate $`c`$. Therefore, to study even one-site properties of the structures, one has to keep in memory information about all connections among them. We performed simulations with a total number of sites (e. g. time) $`t=1000`$ with $`100000`$ averages for decaying networks, $`1<c<0`$, and $`t=10000`$ with $`10000`$ averages for developing networks, $`c0`$. In Fig. 1, we present the dependences of the average connectivity $`\overline{k}(s,t)`$ on the number of a site $`s`$ at different values of $`c`$ for both structures, i.e. for the decaying network, Fig. 1 (a), and for the developing one, Fig. 1 (b).
For both models, in the whole range of $`c`$, $`\overline{k}(s,t)(s/t)^\beta `$. The change of the sign of the exponent $`\beta `$ in the developing network at $`c=1/2`$ was unexpected (compare with the behaviour of $`\beta `$ vs. an aging exponent in networks with aging of sites ), see Fig. 1 (a) and the dependence $`\beta (c)`$ in Fig. 2. At this point, the average connectivity turns to be independent of the site age, $`\overline{k}(s,t)=1`$. We studied also the distribution $`P(k)`$. It behaves as $`k^\gamma `$ for all $`c0`$ for the developing network but, for the decaying network, the power-law dependence is found only in a narrow region of $`c`$ near zero (see Fig. 2). The range of the values of $`k`$ for which we observe such behaviour diminishes with decrease of $`c`$ and then disappears. Note that the finite size effects are strong in this region ($`c<0`$).
We studied also the models in which, unlike the structures considered above, old links between sites are permanently removed with probability proportional to the product of the connectivities of the sites, or new links between old sites are permanently added with equal probability. The simulation demonstrates that the scaling breaks in both cases.
Let us describe the obtained results analytically. We start from the case of the developing network. One may use the simple continuous approach that gives exact results for the scaling exponents as it was demonstrated in . Since one site is added per unit of time, then the total number of sites is $`t`$ and each site is labeled by the time of its birth $`st`$. Then the equation for the average connectivity of the site $`s`$ at time $`t`$, $`\overline{k}(s,t)`$, in the continuous limit may be written in the following form:
$$\frac{\overline{k}(s,t)}{t}=\frac{\overline{k}(s,t)}{\underset{0}{\overset{t}{}}𝑑u\overline{k}(u,t)}+2c\frac{\overline{k}(s,t)\left[\underset{0}{\overset{t}{}}𝑑u\overline{k}(u,t)\overline{k}(s,t)\right]}{\left[\underset{0}{\overset{t}{}}𝑑u\overline{k}(u,t)\right]^2\underset{0}{\overset{t}{}}𝑑u\overline{k}^2(u,t)},\overline{k}(t,t)=1$$
(1)
(we wrote also the boundary condition – one link connected with a new site is added at each time step). The first term is the same as in the BA model, and the second one describes the increase of the connectivity due to the addition of new links between old sites with probability proportional to the product of connectivities of the connected sites. Note that the new links between the old sites can appear only if there is still no links between them. Hence, to write the last term in the present form, we make a strong assumption: we assume that the effect of multiplying of links is not essential at long times. We checked the validity of this non obvious assumption by simulation.
Eq. (1) is simplified at long times:
$$\frac{\overline{k}(s,t)}{t}=(1+2c)\frac{\overline{k}(s,t)}{\underset{0}{\overset{t}{}}𝑑u\overline{k}(u,t)}.$$
(2)
Applying $`_0^t𝑑s`$ to Eq. (2) one gets
$$_0^t𝑑s\frac{\overline{k}(s,t)}{t}=\frac{}{t}_0^t𝑑s\overline{k}(s,t)\overline{k}(t,t)=1+2c,$$
(3)
so we obtain the obvious relation $`_0^t𝑑s\overline{k}(s,t)=2(1+c)t`$. Now Eq. (1) is of the following simple form:
$$\frac{\overline{k}(s,t)}{t}=\frac{1+2c}{2(1+c)}\frac{\overline{k}(s,t)}{t}.$$
(4)
It solution is $`\overline{k}(s,t)=(s/t)^\beta `$ with the exponent
$$\beta =\frac{1+2c}{2(1+c)}.$$
(5)
To obtain the exponent of the distribution of connectivities, $`\gamma `$, one uses the general relation between the scaling exponents of growing networks :
$$\beta (\gamma 1)=1$$
(6)
that was obtained on the assumption that both $`\overline{k}(s,t)`$ and $`P(k)`$ show scaling behaviour. From our simulation, we know that this condition is fulfilled for $`c>0`$. Therefore,
$$\gamma =2+\frac{1}{1+2c},$$
(7)
so we get both scaling exponents for the developing network.
Let us consider now the decaying network. Again we apply the continuous approach. In fact, the removal of old links with equal probability seems to be equivalent to the decrease of the connectivities of old sites with probability proportional to their particular values, so Eq. (2) may also be applicable to this case. Nevertheless, one should account for the fact that only existing links may be removed. Therefore, we prefer to make the calculations more thoroughly.
One introduces the average number of links between the sites $`s`$ and $`s^{}`$ at time $`t`$, $`\overline{n}(s,s^{},t)`$, where $`0ss^{}t`$. The average connectivity may be expressed in terms of this quantity:
$$\overline{k}(s,t)=\underset{0}{\overset{s}{}}𝑑u\overline{n}(u,s,t)+\underset{s}{\overset{t}{}}𝑑w\overline{n}(s,w,t).$$
(8)
The set of equations for $`\overline{n}(s,s^{},t)`$ is
$`\overline{n}(s,t,t)={\displaystyle \frac{\underset{0}{\overset{s}{}}𝑑u\overline{n}(u,s,t)+\underset{s}{\overset{t}{}}𝑑w\overline{n}(s,w,t)}{\underset{0}{\overset{t}{}}𝑑s\left[\underset{0}{\overset{s}{}}𝑑u\overline{n}(u,s,t)+\underset{s}{\overset{t}{}}𝑑w\overline{n}(s,w,t)\right]}},`$ (10)
$`{\displaystyle \frac{\overline{n}(s,s^{},t)}{t}}=c{\displaystyle \frac{\overline{n}(s,s^{},t)}{\underset{0}{\overset{t}{}}𝑑s\underset{s}{\overset{t}{}}𝑑s^{}\overline{n}(s,s^{},t)}}.`$
(Note that $`c`$ is negative now!) The first equality of Eq. (10) describes the links added to the network together with new sites as in the BA model. We again set the number of links connected with each new site to be unit. Applying $`_0^t𝑑s`$ to this equality we get
$$_0^s^{}𝑑s\overline{n}(s,s^{},s^{})=1.$$
(11)
The second equality of Eq. (10) shows how $`\overline{n}(s,s^{},t)`$ changes due to the removing of links between the old sites. Application of $`_0^t𝑑s_s^t𝑑s^{}`$ to this equality leads to the other obvious relation, $`_0^t𝑑s_s^t𝑑s^{}\overline{n}(s,s^{},t)=(1+c)t`$.
Let us search the solution of Eq. (10) in the scaling form
$$\overline{n}(s,s^{},t)=\frac{1}{t}𝒩(\frac{s}{t},\frac{s^{}}{t}).$$
(12)
Then
$`𝒩(\xi ,1)={\displaystyle \frac{\underset{0}{\overset{\xi }{}}𝑑\zeta 𝒩(\zeta ,\xi )+\underset{\xi }{\overset{1}{}}𝑑\zeta ^{}𝒩(\xi ,\zeta ^{})}{\underset{0}{\overset{1}{}}𝑑\xi \left[\underset{0}{\overset{\xi }{}}𝑑\zeta 𝒩(\zeta ,\xi )+\underset{\xi }{\overset{1}{}}𝑑\zeta ^{}𝒩(\xi ,\zeta ^{})\right]}},`$ (13)
$`\left[1\xi {\displaystyle \frac{}{\xi }}\xi ^{}{\displaystyle \frac{}{\xi ^{}}}\right]𝒩(\xi ,\xi ^{})=c{\displaystyle \frac{𝒩(\xi ,\xi ^{})}{\underset{0}{\overset{1}{}}𝑑\xi \underset{\xi }{\overset{1}{}}𝑑\xi ^{}𝒩(\xi ,\xi ^{})}}`$ (14)
and
$$_0^1𝑑\xi 𝒩(\xi ,1)=1.$$
(15)
One sees that
$$_0^1𝑑\xi \left[_0^\xi 𝑑\zeta 𝒩(\zeta ,\xi )+_\xi ^1𝑑\zeta ^{}𝒩(\xi ,\zeta ^{})\right]=2(1+c)$$
(16)
and
$$_0^1𝑑\xi _\xi ^1𝑑\xi ^{}𝒩(\xi ,\xi ^{})=(1+c).$$
(17)
The solution of Eq. (13) may be found in the form:
$$𝒩(\xi ,\xi ^{})=B\xi ^a\xi ^b,$$
(18)
where $`B,a,b`$ are constants. Inserting Eq. (18) into Eq. (13) with account for Eqs. (16) and (17) we obtain the exponents $`a=b=11/[2(1+c)]`$. Eq. (15) gives $`B=1a=1/[2(1+c)]`$. Substitution of Eq. (18) with account for Eq. (12) into Eq. (8) leads to the expression $`\beta =a=11/[2(1+c)]`$ that is exactly the same as in the previous case, see Eq. (5). Note that now it is possible to use the relation between the scaling exponents, Eq. (6), only in the region near $`c=0`$, since only in this region we observed the scaling behaviour of $`P(k)`$. For these values of $`c`$ we again get the old expression Eq. (7) for the exponent $`\gamma `$.
In Fig. 2, we plot the analytically obtained dependences $`\beta `$ and $`\gamma `$ vs. $`c`$ together with the results obtained from the simulation. One may see that the correspondence between the simulation and the theory is quantitative. When $`c`$ changes from $`1`$ to $`0`$, $`\beta `$ increases from $`\mathrm{}`$ to $`1/2`$ passing zero at $`c=1/2`$. Subsequent increase of $`c`$ to $`\mathrm{}`$ leads to growth of $`\beta `$ up to $`1`$ while $`\gamma `$ decreases from $`3`$ to $`2`$. The particular case $`c\mathrm{}`$, $`\gamma =2`$, corresponds to the situation when the network evolves only due to the addition of new links by the defined above rules, that resembles the original small-world networks of Watts and Strogatz .
Our results show that the permanent removing of links leads to a more essential change of the structure of a network than the addition of them. What is the reason for that? One may see that the decaying structure under consideration is, in fact, a changing set of disconnected clusters. Because of finite size effects, we failed to find the position of the percolation threshold that may be defined for networks . Nevertheless, we see that, at high enough rates of link removal, large clusters are certainly absent, and the appearing structures indeed have to demonstrate quite different properties than the networks with $`c0`$. We failed also to find any peculiarity in the distribution of clusters in the point of the scaling break, $`c=1/2`$.
In summary, we have introduced a new parallel component of the evolution of growing networks. In addition to new links connecting new sites and old ones, links between old sites may appear or break with the relative rate $`c`$. We have demonstrated that addition of this component to a scale-free network does not break the scaling behaviour in a wide range of the rate $`c`$.
The following questions remain open. What is the meaning of the point $`c=1/2`$ in which the exponent $`\beta `$ changes sign? Is there any peculiar value of negative $`c`$ below which the power-law behavior of $`P(k)`$ breaks?
SND thanks PRAXIS XXI (Portugal) for a research grant PRAXIS XXI/BCC/16418/98. JFFM was partially supported by the projects PRAXIS/2/2.1/FIS/299/94. We also thank M.C. Marques for reading the manuscript and A.V. Goltsev, Yu.G. Pogorelov and A.N. Samukhin for many useful discussions.
Electronic address: sdorogov@fc.up.pt
Electronic address: jfmendes@fc.up.pt
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# Grassmannians of secant varieties
## 0. Introduction
Let $`X^r`$ be a smooth, non degenerate $`n`$-dimensional projective variety. Take $`P^r`$ and assume that for all lines $`L`$ through $`P`$ the intersection, as a scheme, has length at most one. Then using the projection with center $`P`$ one obtains an embedding $`X^{r1}`$. Such a point $`P`$ exists if and only if $`S_1(X)^r`$, with $`S_1(X)`$ being the secant variety of $`X`$: it is the closure of lines spanned by pairs of distinct points in $`X`$ (see \[JH\] lecture 15). Clearly $`dim(S_1(X))2n+1`$ and $`\mathrm{min}\{r,2n+1\}`$ can be considered as the ”expected dimension” of $`S_1(X)`$. Hence, using such projections, we can embed $`X`$ in $`^{2n+1}`$; we can go further and project $`X`$ isomorphically in some $`^m`$, $`m<2n+1`$ if and only if $`S_1(X)`$ has dimension smaller than the expected one.
The classification of varieties for which $`S_1(X)`$ has dimension less than the expected value was studied by Severi, Terracini and Scorza (see e.g. \[Sc\]) for objects of small dimension and it has been recently reconsidered by several authors. Severi found that the Veronese surface is the unique smooth surface in $`^r`$, $`r5`$ that can be projected isomorphically to $`^4`$. In \[Z\], lower bounds for $`dim(S_1(X))`$ are proved and a classification for varieties attained this lower bound is presented. These varieties can be projected isomorphically to some projective space of dimension much smaller than $`2n+1`$.
In general, projecting a variety $`X^{2n+1}`$ from some disjoint linear subspace $`\pi `$ of dimension $`k>0`$, as $`k`$ increases one expects that points of higher multiplicity must arise in the image. A way (unfortunately not the unique one) to obtain these multiple points is by considering a linear span $`\pi ^{}`$ of $`k+2`$ distinct points of $`X`$ which contains $`\pi `$: clearly $`\pi ^{}`$ is contracted to a $`(k+2)uple`$ point of the projection of $`X`$. The existence of $`\pi ^{}`$ for a general choice of the center $`\pi `$ of projection can be rephrased as follows: consider the Grassmannian $`G(k,r)`$ of $`k`$-planes in $`^r`$ and for a general set $`P_1,\mathrm{},P_{k+2}`$ of points of $`X`$ consider the subset of $`G(k,r)`$ formed by $`k`$-planes contained in the span of the point $`P_i`$’s. As the $`P_i`$’s move, these subsets describe an algebraic subspace $`G(X)G(k,r)`$ which has an obvious expected dimension (see section 1 for more precise definitions) and the singularity above occurs when the closure of $`G(X)`$ coincides with $`G(k,r)`$.
In this setting, varieties for which some space $`G(X)`$ above has dimension less than the expected one are the analogous of varieties with degenerate secant variety; thus we may expect that they are strongly characterized by this property and study their classification.
Notice that this problem was partly considered by classical geometers as a generalization of the Waring problem for forms. Indeed, for varieties $`X`$ which are image of projective spaces under some Veronese embeddings, it arises naturally when one tries to write a set of forms as sum of powers of the same linear forms (see \[B\] and \[T\] for wider discussion).
In the present paper, we propose a systematic study of the subspaces of the Grassmannians arising as above from secant spaces to a projective variety $`X`$; we call these spaces Grassmannians of secant varieties. In the first section we give precise (and more general) definitions. We prove then some general results on the dimensions of these spaces and we show that an exceptional behaviour gives rise to exceptional behaviour also with respect to $`S_k(X)`$, the natural generalization of $`S_1(X)`$.
In the second section, we obtain the classification of irreducible surfaces $`X^r`$, $`r5`$, for which the Grassmannian of lines contained in 3-secant planes has dimension smaller than the expected one.
Indeed, for such surfaces one expects that lines contained in some 3-secant plane describe a subvariety of dimension 8 in $`G(1,r)`$; e.g. when $`r=5`$ one expects that a general line lies in some 3-secant plane (or, in other words, the projection of $`X`$ to $`^3`$ should have some triple point). On the other hand it is classically known that the embedding of $`^1\times ^1`$ in $`^5`$ by a divisor of type $`(2,1)`$ does not enjoy this property. We are able to prove that this is the unique smooth surface for which the Grassmannian of lines in 3-secant planes has dimension less than the expected. Also taking singular (but irreducible) surfaces into account, we have no further examples of degenerate Grasmmannian of 3-secant planes, except for cones.
Notice the curious behaviour of smooth surfaces of minimal degree in $`^5`$. There are 2 types of them: scrolls and the Veronese surface. The second ones have a degenerate secant variety $`S_1(X)`$, but all the Grassmannians of secant varieties are as expected. Surfaces of the first type, instead, have a nice behaviour with respect to secant varieties, but their Grassmannian of 3-secant planes degenerates.
The authors are members of the european AGE project. The first author is supported by the italian MURST fund; he would like to thank Ciro Ciliberto, for many fruitful discussions, and the University of Leuwen, for its warm hospitality. The second author is supported by a Research Fellowship at the University of Leuwen; he would like to thank the University of Siena for its hospitality during the preparation of the manuscript.
## 1. General Properties
Notation
We work over the complex field $``$.
Let $`X^r`$ be an integral non-degenerate variety of dimension $`n`$.
For $`kr`$, a general $`(k+1)`$-uple of points in $`X`$ spans a $`k`$-plane. Therefore we have the ”span” rational map:
$$\mathrm{\Phi }:X^{k+1}\mathrm{}G(k,r)$$
to the Grassmannian of $`k`$-planes in $`^r`$.
For all $`h<k`$, consider the ”incidence” diagram:
$$\begin{array}{ccc}I& \stackrel{\alpha _h}{}& G(h,r)\\ \beta _h& & & & \\ G(k,r)& & \end{array}$$
where $`I`$ is the ”incidence relation” of pairs $`(h,H)`$ with $`hH`$.
Call:
$`G_k(X)=`$ closure of the image of $`\mathrm{\Phi }`$;
$`S_k(X)=`$ closure of $`\alpha _0(\beta _0^1(G_k(X)))`$;
and more generally:
$`G_{h,k}(X)=`$ closure of $`\alpha _h(\beta _h^1(G_k(X)))G(h,r)`$.
In fact:
$`G_{h,k}(X)`$ = closure of $`\{h`$-planes contained in some $`k`$plane, $`(k+1)`$secant to $`X\}`$;
$`S_k(X)=G_{0,k}(X)=`$ closure of the union of all $`k`$planes, $`(k+1)`$-secant to $`X`$.
Observe that we do not consider singular points of $`X`$; e.g. a general line through a double point needs not to be secant in our definition.
For the expected dimensions of these objects, we have:
$`\mathrm{expdim}(G_k(X))=\mathrm{min}\{n(k+1),(rk)(k+1)\};`$
$`\mathrm{expdim}(G_{h,k}(X))=\mathrm{min}\{(kh)(h+1)+n(k+1),(rh)(h+1)\}`$
so that, as usual:
$`\mathrm{expdim}(S_k(X))=\mathrm{min}\{n(k+1)+k,r\}.`$
Notice that all these varieties are irreducible, since $`X`$ is.
###### Proposition 1.1
$`dimG_k(X)=\mathrm{expdim}(G_k(X))`$.
###### Demonstration Proof
First assume $`n>rk`$; then any $`k`$-plane $`\pi `$ meets $`X`$ in (at least) a curve; moving generically $`k+1`$ points of this curve, we see that $`\pi G_k(X)`$. It follows that $`G_k(X)`$ is the whole Grassmannian, hence its dimension is $`(k+1)(rk)`$.
Assume now $`n+kr`$, so that $`\mathrm{expdim}G_k(X)=n(k+1)`$. Since $`X^{(k+1)}`$ has dimension $`n(k+1)`$, the actual dimension of $`G_k(X)`$ is always less or equal than the expected one, and equality means that $`\mathrm{\Phi }`$ has finite general fibers.
Assume $`dimG_k(X)<\mathrm{expdim}(G_k(X))`$ and take $`k`$ minimal, with this property. Since $`k`$ is minimal, the span of $`k`$ general points of $`X`$ meets $`X`$ in a finite set; hence the projection $`X^{}^{rk+1}`$ of X from $`k1`$ general points still has dimension $`n`$, furthermore a general projection from some point of $`X^{}`$ contracts curves on $`X^{}`$; this is possible only if $`X^{}`$ is linear: a contradiction. ∎
###### Corollary 1.2
$`dimG_{h,k}(X)\mathrm{expdim}G_{h,k}(X)`$.
$`dimG_{h,k}(X)=(kh)(h+1)+n(k+1)`$ if and only if a general $`h`$plane in $`G_{h,k}(X)`$ lies only in a finite set of $`k`$-planes in $`G_k(X)`$.
###### Demonstration Proof
Immediate from the fact that $`\beta _h`$ has fibers of dimension $`(kh)(h+1)`$. ∎
###### Corollary 1.3
If $`k<rn`$ and $`\pi G_k(X)`$ is general, then $`\pi X`$ is formed by exactly $`k+1`$ points.
###### Demonstration Proof
Induction on $`k`$. If $`k=1`$ then this is the well known 3-secant lemma: not every secant is a 3-secant. For $`k>1`$, take projection from a general point of $`X`$. ∎
We recall some well-known facts.
###### Proposition 1.4
$`G_{h,k}(X)=G_{h,k+1}(X)`$ implies $`S_k(X)=^r`$.
###### Demonstration Proof
The condition implies $`S_{k+1}(X)=S_k(X)`$; then use \[Z\], V.1.3. ∎
###### Terracini's Lemma
For general points $`P_0,\mathrm{},P_kX`$ and $`u`$ general in their span, one has
$$T_{u,S_k(X)}=<T_{P_0,X},\mathrm{},T_{P_k,X}>.$$
In fact, Terracini’s lemma also works when $`X`$ is reducible.
###### Linear Lemma
Any set of $`m`$-planes such that any two of them meet in a $`(m1)`$-plane, either is contained in some fixed $`^{m+1}`$ or has a $`m1`$-plane for base locus.
###### Demonstration Proof
Call $`H`$ the $`(m+1)`$-plane spanned by two elements $`A,B`$ of the family. Assume that some $`C`$ in the set does not lie in $`H`$; then $`C`$ meets $`H`$ in a $`(m1)`$-plane, since it meets $`A,B`$ in a $`(m1)`$-plane, and $`CH=AB`$; any element of the set contained in $`H`$ must then contain $`CH`$; hence $`AB`$ is the base locus. ∎
Using the linear Lemma and Terracini’s Lemma, we can look at the situation of secant varieties for curves.
###### Proposition 1.5
Let $`X^r`$ be a non degenerate, reduced (but possibly reducible) curve, such that $`dimS_k(X)k+1<r`$. Then $`X`$ is a cone.
###### Demonstration Proof
First take $`k=1`$ and assume $`dimS_1(X)<3`$. Then, by Terracini’s Lemma, the span of two general tangent lines to $`X`$ is a plane, so any pair of tangent lines meets. By \[H\] IV.3.8, if $`X`$ is non degenerate and irreducible, this is impossible. So all components of $`X`$ are degenerate.
If $`X_1`$ and $`X_2`$ are two components, not contained in the same plane, consider a plane $`\pi `$ through $`X_1`$; a general tangent line to $`X_2`$ meets $`\pi `$ in a point $`P`$, which must lie in any tangent line to $`X_1`$; it follows that $`X_1`$ is strange, hence a line through $`P`$. Changing $`X_1`$ and $`X_2`$, we see that all the components of $`X`$ are lines, any two of them meeting at some points. The conclusion now follows from the Linear Lemma.
For $`k>1`$, just work by induction: if $`X`$ is not a cone, then $`dimS_{k1}(X)k+1`$ so $`dimS_k(X)=k+1`$ implies, by Terracini’s Lemma, that for a general choice of $`P_1,\mathrm{},P_kX`$, the linear span of the tangent lines at the points $`P_i`$ is a $`^{k+1}`$ containing the tangent line to any other general point. This is impossible, for $`X`$ is non degenerate.∎
Next, we prove some general results on the behaviour of grassmannians of secant varieties. We are going to use them for the case of surfaces and hope they will prove useful also in higher dimensions, for the classification of varieties whose grassmannians have a degenerate behaviour.
###### Proposition 1.6
If $`dimG_{h,k}(X)=(kh)(h+1)+n(k+1)x`$ for some $`x>0`$, then $`dimG_{h1,k}(X)(kh+1)h+n(k+1)x1`$. In particular $`dimS_k(X)n(k+1)+kxh`$.
###### Demonstration Proof
The first inequality implies that $`\alpha _h`$ has $`x`$-dimensional general fiber, i.e. any $`h`$-plane in some $`(k+1)`$-secant $`k`$-plane is in fact contained in a $`x`$-dimensional family of such planes. If $`\pi ^{}`$ is a general $`(h1)`$-plane in some $`HG_k(X)`$, then we have a ($`kh`$)-dimensional family of $`h`$-planes in $`H`$ containing $`\pi ^{}`$; the inverse image $`I^{}`$ of this family in $`\alpha _h`$ has dimension $`kh+x`$; assume that $`dim\beta _h(I^{})=x^{}x`$. Then over $`H^{}\beta _h(I^{})`$ general we have a fiber of dimension at least $`kh+xx^{}`$; this is only possible if $`x=x^{}`$ and all the $`h`$planes in $`H`$, containing $`\pi ^{}`$, are in fact contained in $`H^{}`$; it follows $`H=H^{}`$, which contradicts $`x>0`$. Therefore $`\beta (I^{})`$ has dimension at least $`x+1`$ and we are done. ∎
###### Proposition 1.7
Assume $`S_{k1}(X)^r`$. If $`X`$ is not a cone and $`dimG_{h,k}(X)=(kh)(h+1)+n(k+1)x`$, $`x>0`$ then either $`dimG_{h1,k}(X)(kh+1)h+n(k+1)x2`$ or $`dimS_k(X)n(k+1)+kk(h+x)`$. In particular $`dimS_k(X)n(k+1)+kx2h`$.
###### Demonstration Proof
Fix $`HG_k(X)`$ general and let $`\pi ^{}`$ be a general $`(h1)`$-plane in $`H`$. Call $`G^{}`$ the $`(kh)`$-dimensional family of $`h`$-planes in $`H`$ through $`\pi ^{}`$ and call $`I^{}`$ its inverse image in $`I`$. We know that $`dimI^{}=kh+x`$ and its image in $`G(k,r)`$ is at least $`(x+1)`$-dimensional, with equality when $`dimG_{h1,k}(X)=(kh+1)h+n(k+1)x1`$; we assume that equality holds and prove that $`dimS_k(X)n(k+1)+kh(h+x)`$.
Take $`H^{}`$ general in the image of $`I^{}`$. The fiber of $`I^{}`$ over $`H^{}`$ is a $`(kh1)`$-dimensional family of $`h`$-planes $`\pi `$ satisfying $`\pi ^{}\pi H`$ and $`\pi H^{}`$; this implies $`dim(HH^{})=k1`$; $`H,H^{}`$ are general moreover the intersections of all elements of $`\beta (I^{})`$ cannot be a fixed $`(k1)`$-plane, for there is an element in $`\beta (I^{})`$ through a general $`h`$-planes of $`H`$ containing $`\pi ^{}`$; it follows by the Linear Lemma that the elements of $`\beta (I^{})`$ lie in a fixed $`(k+1)`$-plane $`V`$.
Since $`HX`$ generate $`H`$, then $`H^{}XHX`$ for $`H^{}\beta (I^{})`$ general; so we may write $`HX=T_1T_2`$ with $`T_1H^{}`$ and $`T_2H^{}=\mathrm{}`$ for $`H^{}\beta (I^{})`$ general. Put $`t=\mathrm{deg}(T_2)`$, so $`k+1t=\mathrm{deg}(T_1)`$. The points of $`T_2`$ move when $`H`$ moves in $`\beta (I^{})`$; this implies that $`V`$ cuts $`X`$ in a curve $`C`$ of degree at least $`t`$.
We claim that $`t3`$. Indeed $`t>0`$ and if $`t<3`$, then the spaces $`H^{}`$ in $`\beta (I^{})`$ have a common ($`k2`$)-plane spanned by $`k1`$ points of $`T_1`$; since $`\pi ^{}`$ is general in $`H`$, it is not contained in this ($`k2`$)-plane; since all $`H^{}\beta (I^{})`$ contain $`\pi ^{}`$, this implies that they have a common $`(k1)`$-plane; we yet know that this leads to a contradiction. Observe that if $`C`$ spans a $`(t1)`$-plane, since $`CH`$ then $`CH`$ spans a $`(t2)`$-plane at most and the span of $`CH`$ and $`T_1`$ has dimension $`t2+k+1t=k1`$; but $`T_2CH`$ and $`<T_1T_2>=H`$, a contradiction. Thus $`dim<C>t3`$.
We claim that $`S_{t2}(C)<C>`$. Otherwise the span of $`S_{t2}(C)T_1`$ contains $`<C>T_1`$ hence also $`H`$; but $`<S_{t2}(C)T_1>S_{t2+k+1t}(X)=S_{k1}(X)`$, hence $`\pi ^{}S_{k1}(X)`$ which in turn implies $`S_k(X)S_{k1}(X)`$, i.e. $`S_{k1}(X)=^r`$, a contradiction.
Now assume $`dimS_1(C)3`$; this yields $`S_{dim<C>2}(C)=<C>`$; since $`HC`$ contains at least $`dim<C>`$ points, we may find $`k+1dim<C>`$ points in $`T_1`$ to conclude that $`VS_{k1}(X)`$, hence again $`S_k(X)S_{k1}(X)`$, a contradiction.
So we have $`dimS_1(C)<3`$. By Proposition 1.5, we see that $`C`$ is a cone; we claim that it has degree $`h+1+x`$. Indeed the $`x`$-dimensional family of $`k`$-planes in $`V`$ through a general $`h`$-plane $`\pi \pi ^{}`$ must contain $`k+1\mathrm{deg}(C)`$ fixed points in $`T_1`$, hence these $`k`$-planes contain a fixed linear subspace of $`H`$ of dimension $`h+k+1\mathrm{deg}(C)`$. It follows $`x=k(h+k+1\mathrm{deg}(C))`$.
Call $`T`$ the vertex of the cone $`C`$. Since $`H`$ is general, by Corollary 1.3 $`HX`$ contains exactly $`k+1`$ points, so $`TH`$; moreover, if $`T`$ is fixed when we move generically one point of $`HX`$ and fix the others, then $`X`$ is a cone, absurd. So $`T`$ describes a subvariety of $`𝒯`$ of $`X`$. We claim that $`dim𝒯`$ is a positive multiple of $`(h+1+x)`$; indeed consider the correspondence $`ZX^{h+1+x}\times X`$, $`Z=\{(P_1,\mathrm{},P_{h+1+x},T):<P_i,T>X`$ for all $`i\}`$. Then the dimension of the fiber of $`Z`$ over $`T`$ is a multiple of $`x+1+h`$, for all points of $`XH`$ can be interchanged, but this fiber has also dimension $`n(h+1+x)dim𝒯`$. In particular $`dim𝒯h+1+x`$.
Fix 2 points $`A,BHX`$ and let the other vary: the corresponding points $`T`$ describe a subvariety of dimension $`h1+x`$ in $`X`$; this yields immediately $`dimT_AT_Bh1+x`$ for the tangent spaces of $`X`$ at $`A,B`$, so that $`dimS_1(X)2n+1(h+x)`$. Now for a third point $`CHX`$, we see that $`dim<T_AT_B>T_Ch1+x`$ so that $`dimS_2(X)3n+22(h+x)`$ and so on: the conclusion $`dimS_k(X)n(k+1)+kk(h+x)`$ follows.
For the last inequality, just call $`h^{}`$ the minimal value such that $`dimG_{h^{},k}(X)n(k+1)+(kh^{})(h^{}+1)x2(hh^{})`$. If $`h^{}=0`$ then $`dimS_k(X)n(k+1)+kx2h`$, otherwise we get $`dimG_{h^{}1,k}(X)n(k+1)+(kh^{}+1)h^{}x2(hh^{})1`$ and the previous conclusion tells us that $`dimS_k(X)n(k+1)+kk(h^{}+x+2(hh^{}))n(k+1)+kx2h`$.∎
###### Remark Example
For $`h=1`$, $`k=2`$ the previous Proposition yields: when $`X`$ is not a cone and $`dimS_1(X)=2n+1<r`$, then $`dimG_{1,2}(X)<3n+2`$ implies also $`dimS_2(X)<3n`$. In this case, the statement is in fact also a consequence of Proposition 1.9 below. Compare with the Proposition 1.6, which just says $`dimS_2(X)3n`$.
###### Theorem 1.8
Assume $`S_{k1}(X)^r`$ and $`dimG_{h,k}(X)<(kh)(h+1)+n(k+1)`$. Then $`dimG_{h1,k1}(X)<(kh)h+nk`$. In particular $`dimS_{kh}(X)<n(kh+1)+kh`$.
###### Demonstration Proof
Take $`\pi G_k(X)`$ general and consider a general $`h`$-plane $`L\pi `$. By Corollary 1.3, we know that $`\pi X=\{p_0,\mathrm{},p_k\}`$. Let $`L^{}`$ be the intersection of $`L`$ with the span of $`p_1,\mathrm{},p_k`$: it is a hyperplane in $`L`$ and it is also a general element of $`G_{h1,k1}(X)`$.
Now move $`\pi `$ in the family of $`(k+1)`$-secant $`k`$-planes through $`L`$; the points $`p_i`$ move consequently, so also their spans move. If $`L^{}`$ moves with $`\pi `$, then it gives a family which is dense in $`L`$; thus any point of $`L`$ lies in some $`k`$-secant $`(k1)`$-plane, which implies $`S_k(X)=S_{k1}(X)`$, whence $`S_{k1}(X)=^r`$ by Proposition 1.4, a contradiction. Thus $`L^{}`$ is fixed, hence it belongs to infinitely many elements of $`G_{k1}(X)`$.
The last statement follows observing that $`S_{kh}(X)S_{kh+1}(X)\mathrm{}S_{k1}(X)`$. ∎
We will apply the previous result to the case $`k=2`$, $`h=1`$. It reads: $`r>2n+1`$ and $`dimG_{1,2}(X)<3n+2`$ implies $`dimS_1(X)<2n+1`$.
###### Theorem 1.9
Assume $`dimG_{1,2}(X)<3n+2`$ and $`dimS_1(X)=2n<r`$. Then $`X`$ is a cone.
###### Demonstration Proof
Take a general 2-plane $`\pi G_2(X)`$; we may assume, by Corollary 1.3, $`\pi X=\{A,B,C\}`$. Take a general point $`P`$ in the line $`<A,B>`$ and fix a general line $`L\pi `$, through $`P`$.
Any line in $`\pi `$ is an element of $`G_{1,2}(X)`$, so, by our assumptions, it is contained in an infinite family of 3-secant planes. Move $`\pi `$ in the family of 3-secant planes through $`L`$, we get a new plane $`\pi ^{}`$, which cuts $`X`$ in $`A^{},B^{},C^{}`$; as observed in the previous proof, by continuity, the hypothesis $`S_1(X)^r`$ implies that $`P`$ is still contained in the line $`<A^{},B^{}>`$; so $`L`$ induces in this way a 1-dimensional family of secant lines through $`P`$.
Move now $`L`$ in $`\pi `$ and consider the induced families of secants. Since $`P`$ is general in $`S_1(X)`$, these families must coincide, for $`P`$ cannot be contained in a 2-dimensional family of secants, by our assumptions. It follows that for $`L^{}\pi `$ general through $`P`$, there exists a 3-secant plane containing $`L^{}`$ and $`<A^{},B^{}>`$; but this implies that every plane through $`A^{},B^{}`$, in the 3-space $`M`$ spanned by $`\pi ,A^{},B^{}`$, is 3-secant to $`X`$. In particular, $`X`$ meets $`M`$ in a curve $`\mathrm{\Gamma }^{}`$, passing through $`C`$.
Observe that $`L`$ is a general line of $`\pi `$ and $`M`$ is also the span of $`\pi `$ and the plane $`\pi ^{}`$ of $`L,A^{},B^{}`$; also observe that $`\pi ^{}`$ is 3-secant and the family of 3-secant planes through $`L`$ is 1-dimensional, by Proposition 1.6, since $`dimS_1(X)=2n`$. It follows that $`M`$ is the span of two general planes in the component of the family of 3-secant planes through $`L`$ containing $`\pi `$; this component is unique, since $`\pi `$ is general.
If $`\mathrm{\Gamma }^{}`$ is non degenerate, then by Proposition 1.5 it must be a cone, otherwise $`S_1(\mathrm{\Gamma }^{})`$ fills up the whole of $`M`$ and $`S_2(X)=S_1(X)^r`$, contradiction. Hence $`\mathrm{deg}\mathrm{\Gamma }^{}2`$. If $`\mathrm{\Gamma }^{}`$ is a conic, then $`X`$ contains a conic through two general points. If $`\mathrm{\Gamma }`$ is such a conic through $`A,B`$, then all the lines through $`P`$ in the plane of $`\mathrm{\Gamma }`$ are secant to $`X`$. Since $`P`$ is general in $`S_1(X)`$, then among these lines there is also the span $`<A^{},B^{}>`$ for there is just one component of secants through $`P`$ containing $`<A,B>`$; on the one hand this is the family of the secants through $`P`$ in the plane of $`\mathrm{\Gamma }`$, on the other hand it contains $`<A^{},B^{}>`$. Thus $`M`$ contains also a conic through $`A,B`$ and proposition 1.5 shows that $`S_1(X)=S_2(X)`$, a contradiction.
We conclude that $`\mathrm{\Gamma }^{}`$ is a line, i.e. $`X`$ contains a line through any general point.
Change now $`C`$ and $`B`$ and change $`P`$ with the point $`Q=L<A,C>`$. As above, we get that there exists a line in $`X`$ passing through $`B`$ and contained in the span of $`\pi `$ and $`\pi ^{}`$, i.e. in $`M`$; hence $`M`$ contains three lines of $`X`$, through the points $`A,B,C`$; since these lines are not in $`\pi `$, they form a non-degenerate curve. Since $`M`$ is contained in $`S_2(X)`$, then by Proposition 1.5 the three lines form a cone, since $`S_1(X)`$ cannot contain $`M`$ by assumption and by Proposition 1.4.
Hence we get that for any triple of points of $`X`$ there pass 3 lines meeting in a common point. As in the proof of Proposition 1.7, this is possible only if $`X`$ is a cone. ∎
## 2. Surfaces and their Grassmannians $`G_{12}`$
Through this section, $`S`$ is always an integral surface in some projective space. We want to classify surfaces for which the variety $`G_{1,2}(S)`$ has dimension smaller than the expected value 8; in other words, a general line in $`G_{1,2}(S)`$ is contained in infinitely many 3-secant planes. Since the situation is clear in $`^4`$, for reasons of dimension, we may assume that $`S^r`$, $`r5`$. It turns out that there are few surfaces with this property; namely:
###### Theorem
The integral surfaces in $`^r`$, satisfying $`dimG_{1,2}(S)<8`$ are either cones or rational normal surfaces of (minimal) degree $`4`$ in $`^5`$, but not Veronese surfaces.
###### Remark Remark 2.1
It is well known that a general line in $`^5`$ lies in some 3-secant plane to a Veronese surface $`S`$. Just to give an elementary argument, look at $`S`$ as the locus of conics of rank 1 and observe that any pencil of conics is contained in some net generated by conics of rank 1. It is classically known that a general projection of a Veronese surface to $`^3`$ has 3 double lines, forming a cone.
The proof of the theorem is divided in several steps. We shall use often the following classical Lemma, due to C.Segre:
###### Segre's Lemma
Let $`S^N`$, $`N4`$ be a non degenerate integral surface containing a 2-dimensional family of plane curves. Then the curves have degree $`2`$ and $`S`$ is either a Veronese surface or a projection of a Veronese surface in $`^4`$.
###### Demonstration Proof
The original proof is in \[S\]. See \[CS\] or \[M\] for a modern proof.∎
###### Remark Step 1
We may reduce ourselves to surfaces in $`^5`$ with $`S_1(S)=S_2(S)=^5`$.
Indeed it follows from Theorem 1.8 that surfaces in $`^r`$, $`r>5`$ for which $`dimG_{1,2}(S)<8`$, also satisfy $`dimS_1(S)<5`$. Since classically it is known that all surfaces in $`^6`$ with this last property are cones, they are yet considered in the classification.
Furthermore, by Theorem 1.9 all surfaces with $`dimG_{1,2}(S)<8`$ and $`dimS_1(S)=4`$ still are cones. Thus we may also assume that $`S`$ is not a Veronese surface.∎
Consider now the incidence variety:
$$I^{}G_2(S)\times ^5I^{}=\{(\pi ,Q):Q\pi \}.$$
call $`p,q`$ the projections; by Proposition 1.1, $`dimI^{}=8`$. Since $`S_2(S)=^5`$, then $`p`$ dominates $`^5`$; so for $`P^5`$ general, all components of $`p^1(P)`$ have dimension 3.
Choose $`P`$ general and choose one component $`L_P`$ of $`p^1(P)`$.
###### Remark Step 2
Let $`W_P=p(q^1(q(L_P)))=`$ the union of all planes belonging to $`L_P`$. Then $`W_P`$ is an irreducible variety containing $`S`$.
Indeed the irreducibility of $`W_P`$ follows immediately from the irreducibility of $`L_P`$.
Assume $`SW_P`$. Then the inverse image of $`L_P`$ in $`S^3\mathrm{}G(2,5)`$ dominates only a curve $`\mathrm{\Gamma }`$. We get an irreducible component $`Y`$ of $`\mathrm{\Gamma }^3`$ such that for $`Q_1,Q_2,Q_3`$ general in $`Y`$, their span lies in $`L_P`$, i.e. $`P<Q_1,Q_2,Q_3>`$ and, for dimensional reasons, general elements of $`L_P`$ are obtained in this way.
Write $`Y=\mathrm{\Gamma }_{01}\times \mathrm{\Gamma }_{02}\times \mathrm{\Gamma }_{03}`$, with each $`\mathrm{\Gamma }_{0i}`$ component of $`\mathrm{\Gamma }`$ and take $`(A,B)`$ general in $`\mathrm{\Gamma }_{01}\times \mathrm{\Gamma }_{02}`$. Since $`P`$ is general and $`S_1(S)=^5`$, then $`P`$ lies only in finitely many secant lines to $`S`$, hence we may assume $`P<A,B>`$; it follows that for $`C\mathrm{\Gamma }_{03}`$ general, we have $`<A,B,C>=<P,A,B>`$ fixed. This contradicts the fact that $`Y`$ induces a 3-dimensional family of planes. ∎
###### Remark Step 3
$`dimW_P=4.`$
Assume $`dimW_P=5`$. Then for $`Q^5`$ general there exists $`\pi L_P`$ with $`Q\pi `$; then the line $`<Q,P>`$ belongs to $`\pi `$. Since the two points $`P,Q`$ are general, we see that a general line belongs to a 3-secant plane, contradicting the assumption $`dimG_{1,2}(S)<8`$.
Assume $`dimW_P3`$. Since $`W_P`$ is irreducible and contains both $`S`$ and the general point $`P`$, then $`dimW_P=3`$. For $`QS`$ general there exists a 3-secant plane contained in $`W_P`$ passing through $`Q`$; it follows that the line $`<P,Q>`$ lies in $`W_P`$, hence $`W_P`$ is the cone over $`S`$ with vertex $`P`$. Since $`W_P`$ contains a 3-dimensional family of planes, the projection of $`W_P`$ from $`P`$ is a surface containing a 3-dimensional family of lines: it is classical that such a surface cannot exist.∎
Choose now two points $`A,BS`$ such that $`P`$ belongs to the line $`\mathrm{}=<A,B>`$; since $`P`$ is general, we may assume that $`A,B`$ are general in $`S`$. Define a rational map:
$$\mathrm{\Psi }:S\mathrm{}G(2,5)$$
which sends $`CS\{A,B\}`$ to the plane $`<A,B,C>`$; call $`L_P^{}`$ the closure of the image. Clearly $`L_P^{}`$ is irreducible and by construction it lies in $`q(p^1(P))`$.
From now on, we start with $`A,B`$ general, then we construct $`L_P^{}`$ and choose the component $`L_P`$ of $`q(p^1(P))`$, containing $`L_P^{}`$.
###### Remark Step 4
$`dimL_P^{}=2`$ and $`p(q^1(L_P^{}))=W_P`$.
By Corollary 1.3, for $`CS`$ general we have $`<A,B,C>S=\{A,B,C\}`$, hence $`\mathrm{\Psi }`$ has finite general fibers, i.e. $`dimL_P^{}=2`$.
Now take $`CS`$ general; the fiber of $`q^1(L_P^{})`$ over $`C`$ is non empty (by construction) and finite, for otherwise $`C`$ lies in infinitely many planes of $`L_P^{}`$. Then $`p(q^1(L_P^{}))`$ has dimension 4. Moreover for $`Q`$ general in $`p(q^1(L_P^{}))`$, there exists $`CS`$ with $`Q<A,B,C>L_P`$; this means $`p(q^1(L_P^{}))W_P`$. The claim follows from irreducibility of $`W_P`$ and Step 3. ∎
For $`\pi L_PL_P^{}`$, write $`\mathrm{\Lambda }_\pi `$ for the span of $`\pi `$ and the line $`\mathrm{}=<A,B>`$.
###### Remark Step 5
$`\mathrm{\Lambda }_\pi `$ is a 3-dimensional space contained in $`W_P`$ and $`W_P`$ is the union of spaces $`\mathrm{\Lambda }_\pi `$, as $`\pi `$ varies. Hence for $`\pi L_P`$ general, the intersection $`\mathrm{\Lambda }_\pi S`$ contains a curve.
Clearly $`dim\mathrm{\Lambda }_\pi =3`$ for $`P\mathrm{}\pi `$. Also for $`Q\pi `$ general, by step 4, there is $`CX`$ with $`Q<A,B,C>`$, whence $`<Q,A,B>W_P`$, so $`\mathrm{\Lambda }_\pi W_P`$. When $`\pi `$ varies, the corresponding $`\mathrm{\Lambda }_\pi `$ define a variety of dimension at least 4, it coincides with $`W_P`$. Since $`SW_P`$, for $`\pi `$ general it follows that $`\mathrm{\Lambda }_\pi S`$ has dimension 1. ∎
Observe that $`\mathrm{\Lambda }_\pi S`$ might have some isolated point, for the ambient fourfold $`W_P`$ might be singular somewhere. Thus although $`A,B`$ belong to $`\mathrm{\Lambda }_\pi S`$, unfortunately we are not allowed to conclude, a priori, that there exists a curve contained in $`\mathrm{\Lambda }_\pi S`$ and passing through $`A,B`$. This makes the argument more involved.
Next step is crucial to override this difficulty and conclude the proof of the Theorem.
###### Remark Step 6
Assume $`S`$ is not a cone. Then for $`\pi ,\pi ^{}L_P`$ general, we have $`\mathrm{\Lambda }_\pi \mathrm{\Lambda }_\pi ^{}=\mathrm{}`$.
Assume on the contrary that $`\mathrm{\Lambda }_\pi \mathrm{\Lambda }_\pi ^{}`$ is a plane $`V`$. First we show that this $`V`$ is fixed as $`\pi ,\pi ^{}`$ vary; indeed otherwise, by the Linear Lemma, the union of the spaces $`\mathrm{\Lambda }_\pi `$ is a $`^4`$ which contains $`S`$, contradicting the assumptions.
The projection $`S\mathrm{}^2`$ from $`V`$ is not dominant, for it contracts the curves $`\mathrm{\Lambda }_\pi S`$, which cover $`S`$. It follows that the projection of $`S`$ from the secant line $`\mathrm{}V`$ is a cone in $`^3`$. The conclusion now follows from:
###### Lemma
If the projection of an integral surface $`S^N`$ $`N4`$, from a general point $`AS`$, is a cone, then $`S`$ itself is a cone.
###### Demonstration Proof
Call $`\tau _A`$ the projection; if $`F`$ is a ruling of the cone $`\tau _A(S)`$, then $`\tau _A^1(F)`$ is a plane curve; thus $`\tau _A`$ determines a 1-dimensional family of plane curves.
Assume these curves are not lines; then they determine their planes, so moving $`A`$ on $`S`$ also the family moves; it follows that $`S`$ is covered by a 2-dimensional family of plane curves; by Segre’s Lemma, $`S`$ is a projection of a Veronese surface. On the other hand it is well known that projecting a Veronese surface or a cubic surface in $`^4`$, which is not a cone, from a general point of the surface, we cannot get a cone.
Assume that all curves $`\tau _A^1(F)`$ are lines, so $`S`$ is ruled; since it is not a cone, its lines meet the line $`R`$ which joins $`A`$ to the vertex of $`\tau _A(S)`$, in a moving point; moving $`A`$ generically, we find a new line $`R^{}`$ which is intersected by all the lines of $`S`$; then $`S`$ lies in the span of $`R,R^{}`$, a contradiction.∎
Since cones are included in our classification, we may assume from now on that $`S`$ is not a cone.
###### Remark Step 7
$`P`$ is not contained in any tangent line to smooth points of $`S`$. Moreover $`\mathrm{}`$ is the unique secant line to $`S`$ contained in $`\mathrm{\Lambda }_\pi `$ and passing through $`P`$.
The first claim follows easily from the observation that the union of all tangent planes to regular points of $`S`$ is a 4-dimensional variety.
Assume that for $`\pi L_P`$ general there exists another secant line $`\mathrm{}^{}`$ with $`P\mathrm{}^{}\mathrm{\Lambda }_\pi `$. $`\mathrm{}^{}`$ must be fixed as $`\pi `$ varies, since $`P`$ is a general point of $`^5`$ and $`S_1(S)=^5`$, so $`P`$ is contained in finitely many secant lines. But this contradicts the previous step, which tells that the intersection of two general spaces $`\mathrm{\Lambda }_\pi \mathrm{\Lambda }_\pi ^{}`$ is $`\mathrm{}`$. ∎
Call now $`\mathrm{\Gamma }_\pi `$ the union of all 1-dimensional components of $`\mathrm{\Lambda }_\pi S`$. As observed in step 5, $`\mathrm{\Gamma }_\pi `$ is non empty, but it may be different from the intersection $`\mathrm{\Lambda }_\pi S`$.
###### Remark Step 8
For $`\pi `$ general in $`L_P`$, there are no components of $`\mathrm{\Gamma }_\pi `$ which are plane curves and contain $`A`$ or $`B`$.
Assume there exists such a component $`\gamma `$ through $`A`$. Since $`\mathrm{}`$ is a general secant line, then it is not tangent to $`S`$, for $`S_1(S)=^5`$. Thus $`\mathrm{}`$ cannot coincide with the embedded tangent space $`T`$ to $`\gamma `$ at $`AS`$. It follows that $`\gamma `$ moves, when $`\pi `$ varies in $`L_P`$, for otherwise a general intersection $`\mathrm{\Lambda }_\pi \mathrm{\Lambda }_\pi ^{}`$ contains $`T`$, contradicting step 6.
Since $`\gamma `$ moves, then $`A`$ is contained in a 1-dimensional family of plane curves lying on $`S`$; since $`A`$ is general, then $`S`$ contains a 2-dimensional family of plane curves; this is impossible by Segre’s Lemma, since $`S_1(S)=^5`$.∎
###### Remark Step 9
The variety $`W_P`$ cannot contain a 2-dimensional family of spaces $`\mathrm{\Lambda }_P`$.
Indeed $`S`$ is non degenerate and we have the following general:
###### Lemma
Let $`Y`$ be an irreducible variety of dimension $`m2`$, containing a 2-dimensional family $``$ of linear spaces of dimension $`m1`$. Then $`Y`$ itself is a linear space.
###### Demonstration Proof
The proof goes by induction on $`m`$, the case $`m=2`$ being classical (and an easy consequence of the Linear Lemma). For $`m>2`$, take a general hyperplane $`H`$ and consider the family $`\{R_iH:R_i\}`$ of subvarieties of $`YH`$. If this family is 2-dimensional, then we conclude by induction. Otherwise, for $`R`$ general, the intersection $`R_H=HR`$ (which is a general hyperplane in $`R`$) must be contained in infinitely many spaces of $``$; but in this case, varying $`H`$, the families $`\{R^{}:R_HR^{}\}`$ together dominate $``$, so $`R`$ meets another general element $`R^{}`$ in a non-fixed subspace of codimension 1. It follows by the Linear Lemma that the total space of $``$ is a linear space of dimension $`m`$, contained in $`Y`$.∎
###### Remark Step 10
We have $`\mathrm{deg}\mathrm{\Gamma }_\pi =3`$ and $`\mathrm{\Gamma }_\pi `$ is a rational normal cubic passing through $`A,B`$.
Since $`L_P`$ is 3-dimensional and all planes of $`L_P`$ belong to some $`\mathrm{\Lambda }_\pi W_P`$, it follows by a dimensional count that the general plane of $`\mathrm{\Lambda }_\pi `$ passing through $`P`$ is a 3-secant plane and by Corollary 1.3 it meets $`S`$ exactly in 3 points; this is clearly impossible unless $`\mathrm{\Lambda }_\pi S`$ is a curve of degree 3.
Take now a general plane $`\pi ^{}\mathrm{\Lambda }_\pi `$ passing through $`\mathrm{}`$; it is a general 3-secant plane, so it meets $`S`$ in exactly 3 points, by Corollary 1.3. But $`\mathrm{deg}\mathrm{\Gamma }_\pi =3`$ and $`\mathrm{\Gamma }_\pi \mathrm{}S\mathrm{}=\{A,B\}`$; also, by step 8, no components of $`\mathrm{\Gamma }_\pi `$ through $`A`$ are plane curves. If $`\mathrm{\Gamma }_\pi `$ does not contain (say) $`A`$, then a general plane through $`\mathrm{}`$ lying in $`\mathrm{\Lambda }_\pi `$ meets $`S`$ in more than 3 points. This is impossible by Corollary 1.3, for this plane is a general 3-secant plane to $`S`$. Thus $`\mathrm{\Gamma }_\pi `$ is an irreducible smooth cubic passing through $`A`$ and $`B`$. ∎
###### Remark Step 11
The linear system $`|L|`$ above determines a birational map from $`S`$ to a quadric surface in $`^3`$.
Take $`A,BS`$ general; for all spaces $`\mathrm{\Lambda }_\pi `$ as above, we find a rational normal cubic $`\gamma S`$, passing through $`A,B`$ and contained in $`\mathrm{\Lambda }_\pi `$; by step 6, moving $`\pi `$ we see that $`S`$ contains a 1-dimensional family of rational normal cubics which intersects transversally in $`A,B`$; moving $`A,B`$, we get a 3-dimensional family of rational normal cubics on $`S`$. Since these curves are contracted by the Albanese map, then $`h^1𝒪_S=0`$ and they belong to a linear system $`|L|`$ of dimension 3; by Noether’s theorem, $`S`$ is rational. Since by step 6 the intersection of two curves of $`|L|`$ through the general points $`A,B`$ is exactly $`\{A,B\}`$, then $`L^2=2`$ and the sequence:
$$0𝒪_S𝒪_S(L)𝒪_\gamma (2)0$$
implies $`h^0𝒪_S(L)=4`$. We get that $`|L|`$ has no base points and it separates general points; the associated map thus sends $`S`$ birationally to a quadric in $`^3`$.∎
Call $`g:S\mathrm{}^3`$ the map associated to $`|L|`$.
###### Remark Step 12: end of the classification
Assume $`g(S)`$ is a smooth quadric. The linear system of hyperplanes of $`S^5`$ corresponds to some divisor class $`(a,b)`$ on $`g(S)`$; since $`g`$ is biratonal on its image, $`a,b>0`$, moreover $`a+b=3`$: it follows that $`S`$ corresponds to the embedding of a smooth quadric surface via the linear system of type $`(2,1)`$ (or $`(1,2)`$).
In fact, embedding a smooth quadric $`Q`$ with the linear system $`|(2,1)|`$ we find a surface $`S^5`$ such that for a general points $`A,B,CS`$ there is a unique rational normal cubic $`\gamma S`$ through them: $`\gamma `$ is the image of the unique conic through the corresponding points on $`Q`$. This curve $`\gamma `$ spans a $`^3`$; call it $`V`$. If $`R`$ is any line contained in the plane spanned by $`A,B,C`$, then a general plane passing through $`R`$ and contained in $`V`$ is 3-secant to $`S`$. Hence $`S`$ is an example of a surface with few lines on 3-secant planes.
Assume $`g(S)`$ is a quadric cone. Call $`X`$ the blow up of $`g(S)`$ at the vertex and let $`T`$ be the class of the proper transform of a ruling; let $`E`$ be the exceptional divisor; the map $`X\mathrm{}^3`$ corresponds to the class $`2T+E`$ on $`X`$. Call $`aE+bT`$ the class corresponding to the map $`X\mathrm{}^5`$; then as above $`(aE+bT)(2T+E)=3`$ which implies $`b=3`$. Since $`g`$ is birational and $`dim(|3T|)=3`$, we se that $`a>0`$; since $`aT+bE`$ has no fixed components and $`E(aE+3T)=2a+3`$, we get $`a=1`$.
Observe that $`E+3T`$ corresponds to rational normal curves passing through the vertex of $`g(S)`$; $`|E+3T|`$ is very ample, in the blow up of $`g(S)`$ and, as above, this system determines a surface $`S^5`$ such that 3 general points of $`S`$ are contained in a rational normal cubic on it; it follows that this embedding of quadric cones also provide example of surfaces in $`^5`$ with few lines on 3-secant planes. ∎
###### Remark Remark 2.2
Rational surfaces of degree 4 in $`^5`$, except for the $`2`$-Veronese embedding of $`^2`$, are rational scrolls of degree 4. There are three types of such scrolls, corresponding to the rank two bundles $`𝒪(2)𝒪(2)`$, $`𝒪(1)𝒪(3)`$ and $`𝒪𝒪(4)`$ over the projective line. Scrolls of the first type can also be seen as a quadric of $`^3`$ embedded in $`^5`$ by means of the linear system $`(2,1)`$. The second type corresponds to the blowing up of a quadric cone at the vertex, embedded by the (pull-back of the) linear system of cubic curves through the vertex. All these scrolls $`X`$ are smooth and have $`dimG_{1,2}(X)<8`$, by step 12 of the theorem.
Scrolls of the third type are cones; also for such $`X`$ one has $`dimG_{1,2}(X)<8`$, for three general points $`A,B,C`$ span, together with the vertex, a 3-plane $`L`$ which meets $`X`$ in three lines, so every line $`r`$ in the 3-secant plane $`<A,B,C>`$ lies in infinitely many 3-secant planes: the planes in $`L`$ through $`r`$.
Observe that rational normal scrolls $`X`$ of degree 4 in $`^5`$ are in fact classically known to project generically to $`^3`$ as surfaces with no triple points; it follows that the general projection of $`X`$ to $`^4`$ has no 3-secant lines through a general point, whence a general line of $`^5`$ does not lie on any 3-secant plane to $`X`$.
Luca Chiantini - Dipartimento di Matematica - Via del Capitano, 15 - 53100 SIENA (Italy) - email: chiantiniunisi.it
Marc Coppens - Department Industrieel Ingenieur en Biotechniek - Katholieke Hogeschool Kempen - Kleinhoefstraat 4 - B 2440 GEEL (Belgium) - email: marc.coppenskhk.be
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# Decoding Information by Following Parameter Modulation With Parameter Adaptive Control
## 1 Introduction
Chaotic dynamics, which have noise-like broadband power spectra, are interesting candidates for encoding and masking information signals in communication. Most approaches proposed to realize this basic idea of using a chaotic signal as broadband carrier are based on the synchronization of coupled chaotic systems . Several schemes have been proposed so for: (I) chaotic masking where the message is added directly to the chaotic carrier with an amplitude much lower than the chaotic carrier; (II) chaotic modulation where the message is injected into the chaotic system to modulate the chaotic carrier, and (III) chaotic switching where a binary message is transmitted by switching between two chaotic attractors associated with two sets of parameters of the system.
Intuitively, the communication is expected to be secure based on two considerations: (1) it is difficult to read out the hidden message by any spectral analysis due to the broadband nature of the chaotic carrier, and (2) exact knowledge of the parameters of the system in the receiver is necessary to recover the message. Thus, a set of the system parameters which serve as the encryption key, is not accessible by any third party.
However, recently some researchers have shown that the security may be spoiled, not by access to the secure set of the parameters, but by some other approaches. For the communication schemes (I) and (II), it has been shown that the hidden message may be unmasked with some nonlinear dynamical forecasting methods . It is believed that this weakness in security is due to low dimension and single positive Lyapunov exponent of the chaotic carrier, and the suggestion is to employ hyperchaotic systems, such as coupled chaotic arrays or time-delay systems in communication. This however may not produce drastic improvement in the security, as shown in a recent report that messages masked by hyperchaos of a six-dimensional system can be unmasked only with a three-dimensional reconstruction, and in our recent work demonstrating that messages masked by chaos of time-delay systems with very hige dimension and many positive Lyapunov exponents can also be extracted successfully.
Another work has also shown that hidden messages can be extracted from chaotic carrier without reconstructing the full dynamics, but using some suitable return maps, which is successfully applied to the Lorenz system for communication schemes I and III.
The idea of encoding by parameter modulation is to use two chaotic attractors $`𝒜`$ and $``$ to represent the two symbols of the digital signals . Since $`𝒜`$, it is possible to construct some suitable return maps which can distinguish the differences between the two attractors, thus reading out the message, just as shown in Ref. . However, if the two attractors are rather complex or the differences between them are very subtle, it may be very difficult to find such distinguishable return maps.
It is natural to ask if it is possible for a motivated intruder to follow the parameter modulation in the transmitter using some parameter adaptive control by the transmitted signal. This paper carrys out security analysis of communication systems using the encoding scheme III. Our analysis is based on the following assumptions:
(a) The intruder does not have access to the precise value of any system parameters in the secure set.
(b) The intruder does know the functional form of the chaotic system in the transmitter.
Our results will show that if a synchronizing system can be constructed using parameter adaptive control by the transmitted signal and the synchronization is robust to parameter mismatches, the messages may be decoded without resorting to the exact parameter values. Since it is practically possible for a motivated intruder to locate a region close to the exact parameters based on the knowledge of the system and the character of the transmitted signal, the security may be spoiled. Robustness of the synchronization to parameter mismatches is an advantage for implementation of the communication scheme but may be a weakness to the security.
## 2 Decoding by parameter adaptive control
Let us consider the following transmitter system
$$\frac{d𝐱}{dt}=𝐅(𝐩,𝐪,𝐱),$$
(1)
where $`𝐩`$ and $`𝐪`$ are system parameters. Binary messages are encoded by switching $`𝐪`$ between $`𝐪_1`$ and $`𝐪_2`$. Based on our assumption, a motivated intruder has known the equations of the system and that $`𝐪`$ are modulated for encoding. He tries to construct a decoding system using parameter adaptive control by the transmitted signal $`s=h(𝐱)`$, as
$`{\displaystyle \frac{d𝐲}{dt}}`$ $`=`$ $`𝐇(𝐩_y,𝐪_y,𝐲,s),`$ (2)
$`{\displaystyle \frac{d𝐪_y}{dt}}`$ $`=`$ $`𝐆(𝐲,𝐪_y,s)(sh(𝐲)),`$ (3)
Suppose that this system (called as the intruder system from now on) is synchronizable with some suitable coupling function $`𝐆`$ if $`𝐩_y=𝐩`$.
In general, it is not always possible that one can find a synchronizable intruder system for any transmitted signal $`s`$ and any subset $`𝐪`$ of the system parameters. However, if with this transmitted signal $`s`$, a synchronizable system $`𝐇(𝐩,𝐪,𝐲,s)`$ can be found by some synchronization schemes, such as Pecora-Carroll decomposition , active-passive decomposition or feedback control without parameter adaptive control, then we can expect that additional parameter adaptive control loops for parameters $`𝐪_y`$ exist if the synchronization is robust to parameter mismatches to some extent, because the system $`𝐇`$ driven by $`s`$ is still stable (the largest conditional Lyapunov exponent is negative) for parameters $`𝐪_y`$ in the vicinity of the point $`𝐪_y=𝐪`$ although exact synchronization is spoiled by parameter mismatches, and the exact synchronization can be restored by bringing the parameters back to the point $`𝐪_y=𝐪`$ using some appropriate control methods. For some systems, such parameter control rules can be found by an analysis based on a global Lyapunov function. In general however, such an analytical treatment may not be possible. In this case, we employ the idea of designing control rule using the information of a control surface constructed by perturbing the parameters $`𝐪_y`$. The essence of the idea is that since the synchronization between the systems $`𝐇`$ and $`𝐅`$ before incorporating parameter adaptive control is robust to parameter mismatches, the synchronization behavior changes smoothly when $`𝐪_y`$ deviate slightly from $`𝐪`$. And there exists a local Lyapunove function with respect to the parameters $`𝐪_y`$ near $`𝐪`$, with the form
$$E(𝐪_y)=𝐔^T𝐔$$
(4)
where $`𝐔=(U_1,U_2,\mathrm{},U_k)^T`$ are time average of some functions $`𝐮=(u_1,u_2,\mathrm{},u_k)^T`$ ($`k`$ is the number of the components in $`𝐪_y`$), i.e.,
$$U_i=\underset{T\mathrm{}}{lim}\frac{1}{T}_0^Tu_i𝑑t(i=1,2,\mathrm{},k).$$
(5)
The function $`u_i`$ has the form $`u_i=\widehat{u}_i(s,𝐲,𝐪_y)(sh(𝐲))`$ so that $`𝐔=\mathrm{𝟎}`$ if $`𝐪_y=𝐪`$. In order that $`E(𝐪_y)`$ is a local Lyapunov function, it is required that $`𝐔`$ is smooth with respect to $`𝐪_y`$ near $`𝐪`$. With this local Lyapunov function, the following evolution system
$$\frac{d𝐔}{dt}=\alpha 𝐔,$$
(6)
is stable at $`𝐔=\mathrm{𝟎}`$. $`\alpha `$ is a convergence parameter. Noting that the convergence of $`𝐔`$ is induced by control of the parameters, we have
$$\frac{d𝐪_y}{dt}=\frac{𝐪_y}{𝐔}\frac{d𝐔}{dt}=\alpha \frac{𝐪_y}{𝐔}𝐔.$$
(7)
In general, it is impossible to obtain an explicit form of $`𝐔`$ for the control rule in Eq. (7). To solve this problem, we use a control surface obtained in simulation or experiment. First we record a time series $`s`$ from system $`𝐅(𝐩,𝐪,𝐱)`$ with a known set of parameter $`(𝐩,𝐪)`$ in the chaotic regime. Then we perturb the parameter $`𝐪_y`$ in the driven system $`𝐇(𝐩,𝐪_y,𝐲,s)`$ slightly from $`𝐪`$ to some values in its vicinity, and compute $`𝐔(𝐪_y)`$ as a function of $`𝐪_y`$. For appropriately chosen function $`𝐮`$, $`𝐔(𝐪_y)`$ are smooth with respect to $`𝐪_y`$ close to the point $`𝐪`$, and in the vivinity, it can be approximated by
$$𝐔(𝐪_y)=M(𝐪_y𝐪),$$
(8)
where the constant $`k\times k`$ matrix $`M`$ is obtained by a local linear fitting of the numerically or experimentally obtained control surface $`𝐔(𝐪_y)`$. Now if the initial value of $`𝐪_y`$ is close to $`𝐪`$, we can replace the Jacobian matrix $`𝐪_y/𝐔`$ in Eq. (7) by the matrix $`M^1`$, i.e.
$$\frac{d𝐪_y}{dt}=\alpha M^1𝐔.$$
(9)
In practice, one can implement the above control rule by replacing $`𝐔`$ with a time average over a period of time $`\tau `$, e.g.
$$U_i(t)=\frac{1}{\tau }_{t\tau }^tu_i𝑑t.$$
(10)
Often, the parameter modulation in the transmitter is much slower than the time scale of the chaotic system, and one can simplify the control rule by replacing the time average $`U_i`$ with $`u_i`$, and finally obtains
$$\frac{d𝐪_y}{dt}=\alpha M^1𝐮.$$
(11)
We can expect that with the above additional parameter adaptive control, the synchronization between the systems is maintained with small enough coupling strength $`\alpha `$ for $`𝐪_y`$ initially close to $`𝐪`$, so that the parameters $`𝐪`$ can be recoved. We can also expect that the synchronization is also robust to mismatches of the rest parameters $`𝐩`$. The function $`u_i`$ can be chosen somewhat arbitrarily, as long as $`U_i`$ (and thus $`E(𝐪_y)`$) is a smooth function of $`𝐪_y`$ in the neighborhood of $`𝐪`$. This scheme provides a general and practical yet simple way to build additional parameter adaptive control loops for originally coupled and synchronized systems, even though a proper choice of the functions $`u_i`$ may still be nontrivial. In this way, the intruder can design systematically an attacking system for the communication scheme based on the knowledge of the system, although so designed control scheme using local information may not be successful when applied to the signal from the transmitter whose parameters the intruder does not know, especially when the transmitter is operating in a parameter region far way from that the intruder uses to build the control rule. However, it is possible for the intruder to get into a neighborhood of the transmitter parameters using some system identification methods based on the knowledge of the system.
Since the intruder system is quite robust to parameter mismatches, the parameter modulation in the transmitter may be revealed and the message decoded without resorting to the exact values of the transmitter parameters $`𝐩`$, but within some tolerable neighborhood. In certain cases, it might also be possible to recover all the transmitter parameters $`(𝐩,𝐪)`$ by designing adaptive control loops for all of them with the above scheme.
In the next sections, we present examples of message decoding based on the above idea. In the first example of Chua’s circuit, the control rule is obtained with a global function analysis and in the second example of Lorenz system, the control rule is designed with the local Lyapunov function scheme.
## 3 An example of Chua’s circuit
As an illustration, we carry out analysis on a specific communication system proposed in Ref. . We first give a brief description of the encoding scheme, and then construct a robust intruder system.
### 3.1 The transmitter
The communication system employs the well-known Chua’s circuit as the chaotic system. The evolution equations for the Chua’s circuit are given by
$`C_1{\displaystyle \frac{dx_1}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{R}}(x_2x_1)h(x_1),`$ (12)
$`C_2{\displaystyle \frac{dx_2}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{R}}(x_1x_2)+x_3,`$ (13)
$`L{\displaystyle \frac{dx_3}{dt}}`$ $`=`$ $`x_2.`$ (14)
The nonlinear characteristic of the Chua’s diode $`h(x_1)`$ is given by
$$h(x_1)=G_1x_1+\frac{1}{2}(G_0G_1)[|x_1+B_p||x_1B_p|],$$
(15)
which is a three-segment piecewise-linear function.
A binary message stream $`I_{in}`$ is encoded by switching between parameters $`G_0`$ , $`G_1`$ and $`G_0^{}=G_0+\frac{1}{r}`$ , $`G_1^{}=G_1+\frac{1}{r}`$ when the stream switches between $`+1`$ and $`1`$, where $`r`$ is a resistor in parallel with Chua’s diode. The parameters used are shown in Table I. Since $`1/r`$ is small (about $`1\%`$ with respect to $`G_0`$ and $`G_1`$), the two chaotic attractors are very similar, as shown in Fig. 1. To examine the similarity, we construct return maps using the consecutive maxima $`x_{max}(n)`$ and $`x_{max}(n+1)`$ from the transmitter signal $`x_1`$, as done in Ref. . The results are shown in Fig. 2, with circles for $`I_{in}=1`$ and crosses for $`I_{in}=1`$ respectively. It is seen that the maps are quite complicated, and most of the points of the two maps coincide and entangle with each other. Distinguishable difference between the two maps is only seen for $`x_{max}(n)`$ around $`0.5`$, which consists of only a small fraction of the maxima. Extracting the message, although is not totally impossible, can be done only for a small portion of the message bits.
### 3.2 The intruder system
Based on our assumption, the intruder does know that the chaotic system of the transmitter is the Chua’s circuit and that the message is encoded by the modulation of $`G_0`$ and $`G_1`$, but does not know the value of any of the parameters $`C_1,C_2,R,L,B_p,G_0,G_1,r`$. Based on the available information, the intruder constructs a receiver system based on the idea of parameter adaptive control as follows:
$`C_1{\displaystyle \frac{dy_1}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{R}}(y_2y_1)h(x_1),`$ (16)
$`C_2{\displaystyle \frac{dy_2}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{R}}(y_1y_2)+y_3,`$ (17)
$`L{\displaystyle \frac{dy_3}{dt}}`$ $`=`$ $`y_2,`$ (18)
$`{\displaystyle \frac{dQ_0}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}x_1[y_1x_1][1\text{sgn}(|x_1|B_p)],`$ (19)
$`{\displaystyle \frac{dQ_1}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}x_1[y_1x_1][1+\text{sgn}(|x_1|B_p)].`$ (20)
where $`Q_0`$ and $`Q_1`$ are controllable parameters of $`h(x_1)`$ which is now
$$h(x_1)=Q_1x_1+\frac{1}{2}(Q_0Q_1)[|x_1+B_p||x_1B_p|].$$
(21)
Eqs. (19-20) mean that $`Q_0`$ is modified when $`|x_1|B_p`$ and $`Q_1`$ is modified when $`|x_1|>B_p`$.
The intruder system of Eqs. (16-21) will synchronize with the transmitter Eqs. (12-15) if the parameters $`C_1,C_2,R,L,B_p`$ are identical for the two systems. To prove it, let us examine the dynamics of the difference $`e_i=y_ix_i(i=1,2,3)`$, $`e_4=Q_0G_0`$ and $`e_5=Q_1G_1`$ given by
$`C_1{\displaystyle \frac{de_1}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{R}}(e_2e_1){\displaystyle \frac{1}{2}}x_1e_4[1\text{sgn}(|x_1|B_p)]{\displaystyle \frac{1}{2}}x_1e_5[1+\text{sgn}(|x_1|B_p)],`$ (22)
$`C_2{\displaystyle \frac{de_2}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{R}}(e_1e_2)+e_3,`$ (23)
$`L{\displaystyle \frac{de_3}{dt}}`$ $`=`$ $`e_2,`$ (24)
$`{\displaystyle \frac{de_4}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}x_1e_1[1\text{sgn}(|x_1|B_p)],`$ (25)
$`{\displaystyle \frac{de_5}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}x_1e_1[1+\text{sgn}(|x_1|B_p)].`$ (26)
The global Lyapunov function
$$E=C_1e_1^2+C_2e_2^2+Le_3^2+e_4^2+e_5^2,$$
(27)
$$\frac{dE}{dt}=\frac{2}{R}(e_1^2+e_2^2)0,$$
(28)
suggests that the state and parameters of the intruder system will converge to those of the transmitter. Fig. 3 illustrates the synchronization process of the system with the attractor in Fig. 1(a) ($`I_{in}=1)`$. The synchronization error decreases exponentially with time, with fluctuations only within small time scales, and we can expect that the synchronization is robust to parameter mismatches. Note that the stable values $`Q_0=0.753`$ and $`Q_1=0.396`$ are just the values of $`G_0`$ and $`G_1`$ in the transmitter.
When the information stream enters the transmitter, lasting a time interval $`T`$ for each bit, the transmitted signal is a sequence switching between the two chaotic attractors in Fig. 1. We take $`T=4.65`$ ms as in Ref. . With the transmitted signal $`s=x_1`$ (Fig. 4(a)) carrying a random message stream, the intruder observed the change of parameter $`Q_0`$ and $`Q_1`$ in Fig. 4(b) and (c) respectively. Switching between the two chaotic attractors results in only small fluctuations of $`Q_0`$, but sudden jumps of $`Q_1`$, because $`|x_1|>B_p`$ most of the time so that $`Q_1`$ is modified more frequently. After a transient of about $`50`$ ms, $`Q_1`$ comes to oscillate slightly about $`0.395`$ for bit $`+1`$ and $`0.385`$ for bit $`1`$. A comparison between the evolution of $`Q_1`$ and the parameter modulation in the transmitter shows clearly that the message can be decoded correctly except for a few bits during the synchronization transient. An interesting thing is that, since $`T`$ is much smaller than the relaxation time of synchronization (about 50 ms, see Fig. 3), the intruder operates in a regime of synchronization transient after the message stream switches from one value to the other. As a result, the oscillation amplitude of $`Q_1`$ (about 10 $`\mu `$s) is larger than the parameter modulation $`\frac{1}{r}=6\mu `$s in the transmitter, which can be an advantage for message decoding.
So far, we use the exact values of the transmitter parameters $`C_1,C_2,R,L,B_p`$ in the intruder system to demonstrate that it is able to follow the parameter modulation in the transmitter by adaptive control. By assumption, the intruder does not have access to these values. However, it is possible to locate an approximate region in the parameter space using some characteristic quantities for system identification based on the knowledge of the chaotic system and at the same time monitor the synchronization error during the scanning of the parameter space. In the following simulation, we suppose that the intruder is able to locate a region within $`20\%`$ deviation from the precise values of the parameters. We choose 5 random values in $`[0.2,0.2]`$ as the relative difference of the above parameters between the intruder and the transmitter, as displayed in table II as an example. For the same information stream in Fig. 4(c), the evolution of $`Q_1`$ is now presented in Fig. 5(a) which is shifted to oscillate quite noisely around $`0.36`$ due to the parameter mismatches. After smoothing the fluctuations with a moving average filter with length of 4 ms, the oscillation of $`Q_1`$ reveals most of the parameter modulation in the transmitter correctly, as seen in Fig. 5(b). We use a simple threshold testing to recover the message, as shown in Fig. 5(c) with a threshold $`Q_{th}=0.365`$. A comparison between the recovered and the original messages has clearly shown that the security is compromised. The results are also robust to external noise, as seen in Fig. 6 for the same parameters as in Fig. 5, but with $`x_1`$ containing noise between $`[0.2,0.2]`$.
## 4 An example of Lorenz system
In the above example, we are able to write down the parameter adaptive control rules based on an analytical treatment by a global Lyapunov function. In the following example, we revisit the communication system in Refs. to illustrate the idea of designing an intruder system using the above local Lyapunov function method, although it has been shown that the message can be extracted using some suitable return maps .
The communication system using Lorenz system is
$`{\displaystyle \frac{dx_1}{d\tau }}`$ $`=`$ $`\sigma (x_2x_1),`$ (29)
$`{\displaystyle \frac{dx_2}{d\tau }}`$ $`=`$ $`rx_1x_2x_1x_3,`$ (30)
$`{\displaystyle \frac{dx_3}{d\tau }}`$ $`=`$ $`x_1x_2bx_3,`$ (31)
where $`\sigma =16.0,r=45.6`$ and $`b`$ is modulated between $`b=4.0`$ and $`b=4.4`$. $`s=x_1`$ is the transmitted signal.
We can design an attacking system with parameter adaptive control for parameter $`b`$ based on the following system coupled by feedback ,
$`{\displaystyle \frac{dy_1}{d\tau }}`$ $`=`$ $`\sigma (y_2y_1)+ϵ(sy_1),`$ (32)
$`{\displaystyle \frac{dy_2}{d\tau }}`$ $`=`$ $`ry_1y_2y_1y_3,`$ (33)
$`{\displaystyle \frac{dy_3}{d\tau }}`$ $`=`$ $`y_1y_2by_3,`$ (34)
which will be synchronized with the system $`x`$ for large enough coupling strength $`ϵ`$. The synchronization is also quite robust to parameter mismatches for large $`ϵ`$. Since by assumption, we do not know the parameter values in the transmitter, we uses $`(\sigma ,r,b)=(10,28,8/3)`$ from a chaotic region in experiment or simulation. With $`ϵ=40`$, for example, the two system is synchronized. Now let systems $`x`$ and $`y`$ have the same $`\sigma `$ and $`r`$, but perturb the parameter $`b`$ in the system $`y`$ around $`b=8/3`$, e.g. $`b_y=b(1+\mathrm{\Delta })`$. We calculate $`U(\mathrm{\Delta })`$ as a function of $`\mathrm{\Delta }`$ by trying the following three simplest functions $`u_1=(sy_1)y_1`$, $`u_2=(sy_1)y_2`$ and $`u_3=(sy_1)y_3`$ . The results of $`U`$ are shown in Fig. 7. It is seen that $`U`$ is a smooth function for $`\mathrm{\Delta }`$ close to zero for the functions $`u_1`$ and $`u_2`$, but not for $`u_3`$. And it is obvious that $`U`$ is also a smooth function for any linear combination of $`u_1`$ and $`u_2`$. Some other choices of $`u`$ is possible, for example $`u=(sy_1)y_1y_3`$. Now we can introduce the additional parameter control loop
$$\frac{db}{d\tau }=\alpha u,$$
(35)
where $`u`$ can be $`u_1`$, $`u_2`$ or their any linear combinations or many other possible choices. The sign of $`\alpha `$ is determined by the sign of $`dU/d\mathrm{\Delta }`$. Simulations have demonstrated that so designed control rules maintain the synchronization for $`\alpha `$ small enough. $`\alpha `$ is allowed to be larger for larger $`ϵ`$. The control is still stable if the system parameters are shifted from $`(\sigma ,r,b)=(10,28,8/3)`$ to those of the transmitter, and the initial values of the parameter $`b`$ in the system $`y`$ does not need to be close to that of the system $`x`$. Since our purpose is to illustrate the designing idea, we do not go into great details on the synchronization behavior in the parameter space $`(ϵ,\alpha )`$. The fact is, in a large region of this parameter space, the synchronization is very robust to mismatches of the rest parameters $`\sigma `$ and $`r`$. An example of the synchronization without and with the additional parameter control loop is shown in Fig. 8 for $`u=u_1=(sy_1)y_1`$, $`(\sigma ,r,b)=(16.0,45.6,4.0)`$ and $`(ϵ,\alpha )=(40,0.1)`$. The synchronization is a little slower when parameter adaptive control is introduced, and it is robust to parameter mismatches because the synchronization error decreases exponentially, with fluctuations only within very small time scales. The parameter $`b`$ comes very close to $`b=4.0`$ in the transmitter from a large initial value $`b=17.0`$ within only a few ms (in the new time scale below).
Now let us use the system to attack the secure communication. In order that the time scale agrees with the system in , we introduce a new time scale $`t=\tau /K`$ where $`K=2505`$ is a scale factor. In the transmitter, the bit duration is $`4`$ ms. As in the above section, we first demonstrate the parameter recovery for identical parameters $`\sigma `$ and $`r`$ in the transmitter and intruder systems, as seen in Fig. 9. Then, we examine the robustness to the parameter mismatches of $`\sigma `$ and $`r`$. The message can be recoved quite reliablely if $`(\sigma _y,r_y)`$ in the intruder system is within a relatively close neighborhood of the transmitter, say, within a $`20\%`$ deviation. Message decoding is generally extremely robust for $`\sigma _y<\sigma `$ and $`r_y>r`$. An example for $`\sigma _y=0.37\sigma `$ and $`r_y=1.72r`$ is shown in Fig. 10. It is seen that $`b`$ has been made to oscillate around $`b=2.2`$ rather than $`b=4.2`$ in the transmitter due to very large parameter mismatches, however, the message is recoved correctly with a moving average filter with length of 2 ms and a simple threshold test.
In the following, we go further to design adaptive control loops for all the three parameters $`(\sigma ,r,b)`$ in the Lorenz system. We find that $`𝐔`$ changes smoothly when the parameters in the originally synchronized systems change slightly if we choose the following three functions
$$u_1=(sy_1)y_1y_3,u_2=(sy_1)y_2,u_3=(sy_1)(y_1+y_2)$$
(36)
The control surface is obtained by perturbing the parameters in the driven system within a $`2\%`$ vicinity of $`(\sigma ,r,b)=(10,28,8/3)`$. $`U_i(i=1,2,3)`$ is the time average of $`u_i`$ in a period of 0.1 second in the time scale $`t`$. After evaluating the matrix $`M^1`$, we obtained the following attacking system with parameter adaptive control loops:
$`{\displaystyle \frac{dy_1}{d\tau }}`$ $`=`$ $`\sigma (y_2y_1)+ϵ(sy_1),`$ (37)
$`{\displaystyle \frac{dy_2}{d\tau }}`$ $`=`$ $`ry_1y_2y_1y_3,`$ (38)
$`{\displaystyle \frac{dy_3}{d\tau }}`$ $`=`$ $`y_1y_2by_3,`$ (39)
$`{\displaystyle \frac{d\sigma }{d\tau }}`$ $`=`$ $`\alpha (0.293u_118.5u_2+15.7u_3),`$ (40)
$`{\displaystyle \frac{dr}{d\tau }}`$ $`=`$ $`\alpha (1.18u_1+95.4u_275.4u_3),`$ (41)
$`{\displaystyle \frac{db}{d\tau }}`$ $`=`$ $`\alpha (0.123u_110.2u_2+8.10u_3).`$ (42)
This system is synchronized for small enough $`\alpha `$ even if the parameters has been shifted to those in the transmitter, i.e. $`(\sigma ,r,b)=(16.0,45.6,4.0)`$. An example of the synchronization and parameter recovery process is shown in Fig. 11 for $`(ϵ,\alpha )=(100,0.2)`$. The convergence rate with the additional parameter adaptive control is slower than that without these control loops. Now if the bit duration in the transmitter is longer than the synchronization transient, the attacking system can follow the parameter modulation in the transmitter and thus decode the message. An example is shown in Fig. 12, where the bit duration is 16 ms. While the parameter $`b`$ clearly follows the modulation, the other two parameters also reflect the switch of the message from one value to the other. It is seen again in Fig. 11 that the synchronization error decreases exponentially, with fluctuations only within very small time scale, so that the message decoding is also very robust to channel noise. Fig. 13 shows the recoved parameters when the transmitted signal $`s=x_1`$ contains an additive noise in $`[1,1]`$.
## 5 Discussion
Based on the assumption that the chaotic system structure is in the public domain and the system parameters are kept in secret as the encryption key in a secure communication system encoding digital message by parameter modulation, we have shown how an intruder might decode the message using an appropriate attacking system with parameter adaptive control by the transmitted signal, even though the intruder does not have access to the exact parameter values in the transmitter. A requirement for the success of this attacking approach is that the intruder can design a synchronizing parameter adaptive control system which is quite robust to mismatches between the parameters of the two systems, so that the message can be recoved without resorting to the exact parameters in the transmitter, but within some neighborhood. Based on the knowledge of the system, it is practically possible for the intruder to get into such a neighborhood using some system identification methods.
For some systems, such a robust synchronizing intruder system with parameter adaptive control can be constructed based on an analysis of a global Laypunov function. Generally, such an analytical treatment is impossible, and we proposed a quite general and practical local Lyapunov function method to design parameter adaptive control rules based on a system which has been synchronized by the transmitted signal. Such a synchronizing system is often obtainable using many possible approaches for constructing synchronization chaotic systems, such as Pecora-Carroll decomposition, active-passive decomposition or feed-back control. In many systems, the synchronization is robust to parameter mismatches if the coupling is not close to the synchronization threshold. The parameter control rules are designed by seeking appropriate functions of the synchronization error whose time average change smoothly when the parameters in the originally synchronized systems deviate slightly from each other. The smooth control surface is obtained in simulation or experiment by perturbing the parameters that will be involved in adaptive control. Although this scheme is quite general, in practice, finding a set of appropriate functions may not be a trivial task when many parameters are involved in modulation. In some cases, the originally synchronized system may be very sensitive to parameter mismatches due to unstable invariant sets embedded in the synchronization manifold , and the proposed local Lyapunov function may not be successful in designing additional parameter adaptive control loops for such systems. However, such systems will not be used in the communication systems because the authorized receiver cannot decode the message correctly in practical environment with unavoidable perturbations.
Employing some system identification methods, the intruder may be able to identify a region near the transmitter parameters in the parameter space in order to design a stable intruder system. Furthermore, the intruder may be able to get close enough to the transmitter parameters by monitoring the synchronization error while scanning the parameter space, so that the message can be decoded with very low rate of errors. During the decoding process, the intruder can improve the decoding by comparison of the results using different parameters in the identified region. In some systems, it is also possible to design adaptive control loops for all the system parameters, so that the message can be decoded even more reliablely.
Based on our investigation, we would like to point out an interesting paradox between robustness and security in chaotic communications. Since in practice, parameter mismatches and external noise is unavoidable, we would require the synchronization systems to be robust to these perturbations, so that high-quality synchronization can be established between the transmitter and the authorized receiver to recover the message correctly in practical implementation. On the other hand, this robustness may be employed by a third party to compromise the security. How to improve the security while maintaining the robustness is an interesting and meaningful research topic for chaotic communications.
Acknowledgements:
This work was supported in part by research grant No. RP960689 at the National University of Singapore. Zhou is supported by NSTB.
Table I. Values of the parameters of the Chua’s circuit in the transmitter.
| $`C_1`$ (nF) | $`C_2`$ (nF) | $`R`$ ($`\mathrm{\Omega }`$) | $`L`$ (mH) | $`G_0`$ (ms) | $`G_1`$ (ms) | $`1/r`$ ($`\mu `$s) | $`B_p`$ (V) |
| --- | --- | --- | --- | --- | --- | --- | --- |
| 10 | 100 | 1680 | 18 | $``$0.753 | $``$0.396 | 6 | 1 |
Table II. Parameters of the transmitter and the intruder and their relative differences.
| transmitter | $`C_1`$ | $`C_2`$ | $`R`$ | $`L`$ | $`B_p`$ |
| --- | --- | --- | --- | --- | --- |
| intruder | $`0.817C_1`$ | $`1.163C_2`$ | $`1.072R`$ | $`0.897L`$ | $`0.849B_p`$ |
| differences ($`\%`$) | $``$18.3 | 16.3 | 7.2 | $``$10.3 | 15.1 |
Figure Captions
Fig. 1. Two attractors used to encode binary information, with (a) corresponding to bit $`+1`$ and (b) to $`1`$.
Fig. 2. Return maps of the consecutive maxima of the transmitted signal $`x_1`$. Most of the points for the attractor in Fig. 2(a) (circles) and those for the attractor in Fig. 2(b) (crosses) do not have distinguishable separation.
Fig. 3. Synchronization process of the intruder to the attractor in Fig. 2(a)
(a) Synchronization error $`\sqrt{e_i^2}`$. (b) Convergence of parameters $`Q_0`$ and $`Q_1`$ to those of the transmitter $`G_0=0.753`$ and $`G_1=0.396`$, respectively. In all the figures in this paper, the unit of time is ms.
Fig. 4. The process of following the parameter modulation in the transmitter.
(a) The transmitted signal $`s=x_1`$. (b) Change of $`Q_0`$. (c) Change of $`Q_1`$. The dotted line shows the parameter modulation in the transmitter.
Fig. 5. An example of decoding process.
(a) Evolution of $`Q_1`$. It oscillates noisely due to quite large parameter mismatches. (b) Smoothed $`Q_1`$ by moving average with length of 4 ms. The dotted line shows the parameter modulation in the transmitter. (c) The decoded message by threshold testing with a threshold $`Q_{th}=0.365`$. (d) The original message stream.
Fig. 6. Analogous to Fig. 6, but with $`s`$ containing noise between $`[0.2,0.2]`$.
Fig. 7. $`U`$ as a function of the relative deviation of the parameter $`b`$ for different choice of function $`u`$. The smooth functions can be used to design parameter adaptive control loop.
Fig. 8. Synchronization and parameter recovery with the additional parameter adaptive control loop of Eq. (35). (a) and (b) are synchronization errors without and with this control loop respectively. (c) is the convergence of the parameter $`b`$.
Fig. 9. The process of following the parameter modulation in the case that the rest of the parameters $`\sigma `$ and $`r`$ are identical.
Fig. 10. Illustration of message decoding when the parameters $`\sigma `$ and $`r`$ have large mismatches between the systems. (a) Evolution of $`b`$. (b) smoothed $`b`$ by moving average filter with length of 2ms. (c) The decoded message by threshold testing with a threshold $`b_{th}=2.2`$. (d) The original message stream.
Fig. 11. Synchronization and parameter recovery with the additional parameter adaptive control loops for all the three parameters. (a) Plots 1 and 2 are synchronization errors without and with these control loops respectively. (b), (c) and (d) are the evolutions of the parameters $`\sigma `$, $`r`$ and $`b`$, respectively.
Fig. 12. The process of following the parameter modulation in the transmitter. (a) is an input message stream. (b), (c) and (d) are the evolutions of the parameters $`\sigma `$, $`r`$ and $`b`$, respectively. The bit duration is 16 ms.
Fig. 13. Robustness of the message decoding in the presence of channel noise within $`[1,1]`$.
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# Expansion of real valued meromorphic functions into Fourier trigonometric series
## 1. Introduction
From the author’s viewpoint, proper attention should be paid to a class of mathematical expressions reducing, in the boundary case, to the difference of infinities $`\mathrm{}\mathrm{}`$. Namely, in general case, the difference of infinities $`\mathrm{}\mathrm{}`$ is an indefinite expression taking any value from the extended numeric straight-line (real axis) $`R^{}`$, that is, from the segment $`[\mathrm{},+\mathrm{}]`$ . Causality related to the value itself of an indefinite expression is behavior of the mathematical expression reducing to it in the boundary case. Therefore, a case of so-called alternative numerical series is indicative one, . In fact, since each alternative series can be expressed, according to its definition - Definition 1, Section 2.6, Chapter 2, p. 28, \- by difference of two positive numerical series, then in the case when each of them definitely diverges the alternative series reduces to an indefinite expression $`\mathrm{}\mathrm{}`$. The typical representative of such a class of series, is the alternative numerical series $`\underset{k=1}{\overset{+\mathrm{}}{}}\left(1\right)^k`$. By means of Caushy’s definition of both the sequence convergence and the numerical series convergence , it can be immediately shown that the alternative numerical series $`\underset{k=1}{\overset{+\mathrm{}}{}}\left(1\right)^k`$, indefinitely diverges. In other words, since in this emphasized case, the limiting value of partial sums $`\underset{k=1}{\overset{𝑛}{}}\left(1\right)^k`$ of the numerical series $`\underset{k=1}{\overset{+\mathrm{}}{}}\left(1\right)^k`$ does not exist when $`n+\mathrm{}`$, then the sum of the observed numerical series, in the Caushy’s sense, does not exist too. Accordingly, it is natural to ask the following questions: How much is the sum value of this numerical series, more exactly, is this numerical series summable? Closely related to these questions is another: How much is, in this emphasized instance, the numerical value of an indefinite expression of difference of infinities $`\mathrm{}\mathrm{}`$ ? Clearly, conceptually distinction should be made between summation of series in the Caushy’s sense and its summability. In the modern mathematical analysis, more exactly, in the series theory, and from the point of view of the general convergence (summability) of numerical series, an answer to the former questions was given by Frobenius, Holder and Cesaro . In this paper it is presented slightly different and more indirect answer more related to the problem of exact determining rather than to the problem of redefining the sum itself of indefinitely divergent series. It is more indirect because the fundamental conclusions will be based on the results of the complex analysis theory, more exactly, on contour integration of the functions of a complex variable. In this case also, a proper attention will be paid to the mathematical expressions reducing, in the boundary case, to an indefinite expression of difference of infinities $`\mathrm{}\mathrm{}`$.
## 2. The main results
### 2.1. Contour integration and improper integral
The concept of an improper integral absolute existence of a meromorphic function $`f\left(z\right)`$: $`C^1C^1`$; $`C`$\- is the set of complex numbers, clearly in the case when the singularities of the function lie onto the integration contour, is based on concept of total value ($`v.t.`$) of an improper integral of a meromorphic function which is defined to be the sum of Caushy’s principal value ($`v.p.`$) and Jordan’s singular value ($`v.s.`$) of an improper integral of a meromorphic function $`f\left(z\right)`$. Jordan’s singular value ($`v.s.`$) of an improper integral of a meromorphic function is defined to be a limiting value, as $`\epsilon 0^+`$, of an integral of the function $`f\left(z\right)`$ over a certain part $`\stackrel{}{PQ}`$ of a circular path of integration $`G_\epsilon `$: $`G_\epsilon =\{zz\left(\theta \right)=\epsilon e^{i\theta }\text{;}\text{ }\theta [0,2\pi ]\}`$, bypassing a singularity of the function, where the points of the complex plane: $`P`$ and $`Q`$, are intersection points of the circular contour $`G_\epsilon `$ and an integration contour $`G`$, $`i`$\- denotes an imaginary unit and $`e`$\- is a base of natural logarithm. In fact, in other words a concept of improper integrals absolute existence of meromorphic functions generalizes the fundamental concept of improper integrals convergence (existence).
By results of both so-called Jordan’s lemma - Theorem 1, Subsection 3.1.4, Section 3.1, Chapter 3, p. 52, \- and the fundamental Caushy’s theorem on residues - Theorem 1, Subsection 3.6.2, Section 3.6, Chapter 3, p. 226, \- the sum of Caushy’s principal value ($`v.p.`$) and Jordan’s singular value ($`v.s.`$) of an improper integral of a meromorphic function $`f\left(z\right)`$ whose only singularity (simple pole) lies onto closed integration contour $`G`$, can be proved to be<sup>1</sup><sup>1</sup>1Symbol $`\underset{𝐺}{\overset{}{}}`$ denotes an integration over the closed contour of integration G, in this case in the positive mathematical direction
$$v.t.\underset{𝐺}{\overset{}{}}f\left(z\right)dz=v.p.\underset{𝐺}{\overset{}{}}f\left(z\right)dz+v.s.\underset{G_\epsilon }{\overset{}{}}f\left(z\right)dz=$$
(2.1)
$$=v.p.\underset{𝐺}{\overset{}{}}f\left(z\right)dz+\{\begin{array}{c}i\alpha A\hfill \\ i\left(2\pi \alpha \right)A\hfill \end{array}=\{\begin{array}{c}0\hfill \\ i2\pi A\hfill \end{array},$$
where $`\alpha `$ is an absolute value of an angular difference of arguments of intersection points: $`Q`$ and $`P`$, with respect to the point $`z_0`$, respectively, in the limit as $`\epsilon 0^+`$, and $`A`$ is a residue of the function $`f\left(z\right)`$ at the point $`z_0`$, that is: $`A=\underset{zz_0}{lim}\left(zz_0\right)f\left(z\right)`$, on condition that such limiting value exists, and . In this emphasized case, as distinguished from Caushy’s principal value ($`v.p.`$), Jordan’s singular value ($`v.s.`$), just as well as the total value ($`v.t.`$), of an improper integral of a meromorphic function $`f\left(z\right)`$, are not uniquely defined ones, already they depend upon the choice of the part of a circular path bypassing the singularity of the function.
In the case when the singularity of the meromorphic function $`f\left(z\right)`$ at the point $`z_0`$ is a pole of a higher order, the conclusion essentially differs from the preceding one. Namely, in that case, both the Caushy’s principal value ($`v.p.`$) and Jordan’s singular value ($`v.s.`$) of an improper integral do not exist in the limit as $`\epsilon 0^+`$. The improper integral reduces to an indefinite expression of the difference of infinities $`\mathrm{}\mathrm{}`$. Thus, as for the meromorphic function $`z\frac{1}{\left(zz_0\right)^k}`$ , where $`k2`$ ($`kN`$) and $`N`$ is a set of natural numbers, the improper integral along an any closed integration path passing through the point $`z_0`$ absolutely exists and its unique total value ($`v.t.`$) is identically zero. This is in agreement with both the general Caushy-Goursat’s integral theorem - Theorem 1, Subsection 3.5.2, Section 3.5, Chapter 3, p. 203, \- and the results of contour integration of rational functions - Subsubsection 3.1.2.3, Subsection 3.1.2. Section 3.1, Chapter 3, p. 46, . Now then, in this emphasized case, a sum of values of integrals of meromorphic function $`f\left(z\right)`$ over a part of any integration path $`G`$ between intersection points as well as over the part $`\stackrel{}{PQ}`$ of circular path $`G_\epsilon `$, is identically zero for each $`\epsilon `$. Since choice of an arc path: $`\stackrel{}{PQ}`$ or $`\stackrel{}{PQ}`$, bypassing a singularity of the function is arbitrary, the above-mentioned unique sum remains zero in the limit as $`\epsilon 0^+`$.
#### 2.1.1. Example
Let the part of an integration path between points: $`P`$ and $`Q`$, be a part of circumference of a circle centred at the orgin and of radius $`a`$: $`aR_+^1`$ ($`R_+^1`$\- is a set of positive real numbers) and $`a`$ be also a singularity of the function $`f\left(z\right)`$: $`f\left(z\right)=\frac{1}{\left(za\right)^k^{}}`$, $`k^{}2`$ and $`k^{}N`$. Then, an integral of the function $`f\left(z\right)`$, over the part of a circular integration path $`G`$ from the point $`P`$ to the point $`Q`$ which dose not contain the singularity $`a`$, reduces to the integral
$$a\underset{\alpha ^{}}{\overset{2\pi \alpha ^{}}{}}\frac{e^{i\theta ^{}}}{\left(ae^{i\theta ^{}}a\right)^k^{}}d\theta ^{}=$$
(2.2)
$$=\frac{\left(i\right)^k}{\left(2a\right)^{k+1}}\left[\underset{𝛼}{\overset{\pi \alpha }{}}\frac{\mathrm{cos}\left(k\theta \right)}{\left(\mathrm{sin}\theta \right)^{k+2}}d\theta i\underset{𝛼}{\overset{\pi \alpha }{}}\frac{\mathrm{sin}\left(k\theta \right)}{\left(\mathrm{sin}\theta \right)^{k+2}}d\theta \right]\text{,}$$
where $`\alpha ^{}`$ is an argument of the point $`P`$ with respect to the origin ($`\alpha =\frac{\alpha ^{}}{2}`$), $`\theta =\frac{\theta ^{}}{2}`$ and $`k=k^{}2`$ ($`kN_0`$; $`N_0=\{0,1,2,\mathrm{}\}`$).
As results of partial integration the following integral dependencies are obtained:
$$i(1+k)\underset{𝛼}{\overset{\pi \alpha }{}}\frac{2ae^{2i\theta }}{\left(ae^{2i\theta }a\right)^{k+2}}d\theta =\left(\frac{i}{2a}\right)^{k+1}[\frac{\mathrm{cos}\left(k\theta \right)}{\left(\mathrm{sin}\theta \right)^k}\mathrm{cot}\theta |{}_{\alpha }{}^{\pi \alpha }$$
(2.3)
$$k\underset{𝛼}{\overset{\pi \alpha }{}}\frac{\mathrm{sin}\left[\left(k1\right)\theta \right]}{\left(\mathrm{sin}\theta \right)^{k+1}}d\theta i(k+1)\underset{𝛼}{\overset{\pi \alpha }{}}\frac{\mathrm{sin}\left(k\theta \right)}{\left(\mathrm{sin}\theta \right)^{k+2}}d\theta ]$$
and
$$(k+2)\underset{𝛼}{\overset{\pi \alpha }{}}\frac{\mathrm{sin}\left[\left(k+1\right)\theta \right]}{\left(\mathrm{sin}\theta \right)^{k+3}}d\theta =2\frac{\mathrm{cos}\left(k\theta \right)}{\left(\mathrm{sin}\theta \right)^k}\mathrm{cot}\theta |{}_{\alpha }{}^{\pi \alpha }$$
(2.4)
$$2k\underset{𝛼}{\overset{\pi \alpha }{}}\frac{\mathrm{sin}\left[\left(k1\right)\theta \right]}{\left(\mathrm{sin}\theta \right)^{k+1}}d\theta \frac{\mathrm{sin}\left(k\theta \right)}{\left(\mathrm{sin}\theta \right)^{k+1}}\frac{1}{\mathrm{sin}\theta }|{}_{\alpha }{}^{\pi \alpha }\text{.}$$
As, in this emphasized case, $`2a\mathrm{sin}\alpha =\epsilon `$ a by-pass integral value is:
$$\epsilon ^{\left(k+1\right)}\underset{\left({\displaystyle \frac{\pi }{2}}+\alpha \right)}{\overset{\frac{\pi }{2}+\alpha }{}}e^{i\left(k+1\right)\theta }d\theta =\{\begin{array}{c}\frac{\left(1\right)^n}{n}\frac{\mathrm{sin}\left(2n\alpha \right)}{\left(2a\mathrm{sin}\alpha \right)^{2n}}\text{;}\text{ }k=2n1\hfill \\ \frac{2\left(1\right)^n}{2n+1}\frac{\mathrm{cos}\left[\left(2n+1\right)\alpha \right]}{\left(2a\mathrm{sin}\alpha \right)^{2n+1}}\text{;}\text{ }k=2n\hfill \end{array}\text{}nN\text{.}$$
(2.5)
By the comparative analysis of preceding equalities it can be easily shown that if the sum of the integral value of the function $`f\left(z\right)`$ over the circular integration contour $`G`$ from the point $`P`$ to the point $`Q`$ and the by-pass integral value, for arbitrarily chosen $`\alpha `$ and $`k=2n1`$ respectively, is identically zero, just as well as the value of the integral $`\underset{𝛼}{\overset{\pi \alpha }{}}\frac{\mathrm{sin}\left[\left(k1\right)\theta \right]}{\left(\mathrm{sin}\theta \right)^{k+1}}d\theta `$ for $`k=2n1`$, see the equation (2.4), then for $`k=2n1`$ and $`nN`$ it holds
$$\underset{𝛼}{\overset{\pi \alpha }{}}\frac{\mathrm{sin}\left(k\theta \right)}{\left(\mathrm{sin}\theta \right)^{k+2}}d\theta =\frac{1}{n}\frac{\mathrm{sin}\left(2n\alpha \right)}{\left(\mathrm{sin}\alpha \right)^{2n}}\text{.}$$
(2.6)
Since the second integral on the right-hand side of the equation (2.3) is identically zero for $`k=2n`$, then in view of the preceding result (2.6) it follows that
$$\underset{𝛼}{\overset{\pi \alpha }{}}\frac{2ae^{2i\theta }}{\left(ae^{2i\theta }a\right)^{2\left(n+1\right)}}d\theta =\frac{2\left(1\right)^n}{2n+1}\left\{\frac{\mathrm{cos}\left[\left(2n+1\right)\alpha \right]}{\left(2a\mathrm{sin}\alpha \right)^{2n+1}}\right\}\text{}nN\text{.}$$
(2.7)
If one considers the fact that, for a corresponding natural number $`n`$ ($`nN`$), the sum of functional expressions on right-hand sides of relations: (2.5) (for $`k=2n`$) and (2.7), is identically zero for arbitrarily chosen $`\alpha `$, then the sum of integrals on the left-hand sides of these equations is identically zero too. Thus, the improper integral of the function $`f\left(z\right)`$ along the circular contour of integration $`G`$ absolutely exists and its total value, as limiting value of a sum of integrals (2.2) and (2.5) as $`\alpha 0^+`$, is identically zero for each $`k^{}2`$ and $`k^{}N`$, and what has been just proved by method of a mathematical induction<sup>2</sup><sup>2</sup>2For $`k=0`$ and $`k=1`$ i.e. $`k^{}=2`$ and $`k^{}=3`$, on the basis of the relation (2.3) it holds $`\underset{𝛼}{\overset{\pi \alpha }{}}\frac{2ae^{2i\theta }}{\left(ae^{2i\theta }a\right)^{k+2}}d\theta =\frac{\mathrm{cot}\alpha }{a}`$ and $`\underset{𝛼}{\overset{\pi \alpha }{}}\frac{2ae^{2i\theta }}{\left(ae^{2i\theta }a\right)^{k+2}}d\theta =\frac{\mathrm{cot}\alpha }{2a^2}`$, respectively. From the relation (2.5), the values of by-pass integrals, in these emphasized cases, are equal to: $`\frac{1}{\epsilon }\underset{\left({\scriptscriptstyle \frac{\pi }{2}}+\alpha \right)}{\overset{\frac{\pi }{2}+\alpha }{}}e^{i\theta }d\theta =\frac{\mathrm{cot}\alpha }{a}`$ and $`\frac{1}{\epsilon ^2}\underset{\left({\scriptscriptstyle \frac{\pi }{2}}+\alpha \right)}{\overset{\frac{\pi }{2}+\alpha }{}}e^{2i\theta }d\theta =\frac{\mathrm{cot}\alpha }{2a^2}`$, respectively. Accordingly, the total value of improper integral $`i\underset{𝑜}{\overset{2\pi }{}}\frac{ae^{i\theta ^{}}}{\left(ae^{i\theta ^{}}a\right)^k^{}}d\theta ^{}`$, is identically zero, for both $`k^{}=2`$ and $`k^{}=3`$..$`\mathrm{}`$
As it has been just illustrated by the previous example, the concept of an improper integral absolute existence is more general concept with respect to the concept of an improper integral convergence. This is in connection with indefinite expression of the difference of infinities $`\mathrm{}\mathrm{}`$ to which the improper integral value is reduced in a boundary case. Namely, independently from the fact that the improper integral absolutely exists, in some of the concrete cases its principal value does not exist. Hence, by introducing a by-pass integral into the analysis the concept itself of improper integral convergence (existence) is generalized to the concept of improper integral absolute existence.$`\mathrm{}`$
### 2.2. Fourier trigonometric series of real valued meromorphic functions
#### 2.2.1. An analysis of an idea
Without loss of the generality, one may assume that a complex meromorphic function $`g(z,t)`$, where the variable $`t`$ is independent one with respect to the complex variable $`z`$, has infinitely but a count of many simple poles: $`a_1,a_2,\mathrm{}`$ onto the imaginary axis of the complex plane $`C^1`$. In that emphasized case, there exists a sequence of circular contours of integration $`G_r`$, centred at the origin and of radius $`r`$, such that onto theirs boundaries there are no singularities of the function $`g(z,t)`$. Hence, by the fundamental Cauchy’s theorem on residues the sequence of the partial sums can be formed
$$\underset{k=1}{\overset{𝑛}{}}A_k\left(t\right)=\frac{1}{2\pi i}\underset{G_r}{\overset{}{}}g(z,t)dz\text{,}$$
(2.8)
where $`A_k\left(t\right)`$: $`A_k\left(t\right)=\underset{z=a_k}{Res}g(z,t)`$, are residues of the function $`g(z,t)`$ at the points: $`z=a_k`$, $`k=1,2,\mathrm{},n`$.
On the one hand, on the basis of the second Jordan’s lemma - Theorem 2, Subsection 3.1.4, Section 3.1, Chapter 3, p. 52, \- if there exists an unique limiting value: $`\underset{\left|z\right|+\mathrm{}}{lim}\left[zg(z,t)\right]`$; for each $`zC^1`$, then the sequence of the partial sums $`\underset{k=1}{\overset{𝑛}{}}A_k\left(t\right)`$ converges, in other words there exists a sum of the infinite functional series $`\underset{k=1}{\overset{+\mathrm{}}{}}A_k\left(t\right)`$ in the Cauchy’s sense:
$$\underset{k=1}{\overset{+\mathrm{}}{}}A_k\left(t\right)=\underset{\left|z\right|=+\mathrm{}}{Res}g(z,t)\text{,}$$
(2.9)
where $`2\pi i\underset{\left|z\right|=+\mathrm{}}{Res}g(z,t)=\underset{r+\mathrm{}}{lim}\underset{G_r}{\overset{}{}}g(z,t)dz=2\pi i\underset{\left|z\right|+\mathrm{}}{lim}\left[zg(z,t)\right]`$.
However, on the other hand, by the same Jordan’s lemma, if there exists no an above mentioned limiting value: $`\underset{\left|z\right|+\mathrm{}}{lim}\left[zg(z,t)\right]`$; for each $`zC^1`$, already there exist only partial limiting values: $`\underset{\left|z\right|+\mathrm{}}{lim}\left[zg(z,t)\right]`$; $`\mathrm{R}ez>0`$ and $`\underset{\left|z\right|+\mathrm{}}{lim}\left[zg(z,t)\right]`$; $`\mathrm{R}ez<0`$, then there exists a infinite sum of the residues of the function $`g(z,t)`$ that is equal to the limiting sum of integral values:
$$\underset{k=1}{\overset{+\mathrm{}}{}}A_k\left(t\right)=\frac{1}{2\pi i}\underset{r+\mathrm{}}{lim}\left[\underset{G_r^R}{\overset{}{}}g(z,t)dz+\underset{G_r^L}{\overset{}{}}g(z,t)dz\right]\text{,}$$
(2.10)
where the integral paths: $`G_r^R=\{zz\left(\theta \right)=re^{i\theta }\text{;}\text{ }\theta [\frac{\pi }{2}+\delta \left(r\right),\frac{\pi }{2}\delta \left(r\right)]\}`$ and $`G_r^L=\{zz\left(\theta \right)=re^{i\theta }\text{;}\text{ }\theta [\frac{\pi }{2}+\delta \left(r\right),\frac{3\pi }{2}\delta \left(r\right)]\}`$, are arc parts of the circular path of integration $`G_r`$ in the right-hand and left-hand half-plane of the complex plane $`C^1`$, respectively, and an arbitrary angular function $`\delta \left(r\right)`$, which is of sufficiently small real positive values for any positive values of $`r`$, satisfies the condition: $`\underset{r+\mathrm{}}{lim}\delta \left(r\right)=0`$. In other words, although in this emphasized case there exists no a sum of the infinite functional series $`\underset{k=1}{\overset{+\mathrm{}}{}}A_k\left(t\right)`$ in the Cauchy’s sense, this infinite functional series is summable. Note that in this case too: $`\underset{k=1}{\overset{+\mathrm{}}{}}A_k\left(t\right)=\underset{\left|z\right|=+\mathrm{}}{Res}g(z,t)`$, where $`2\pi i\underset{\left|z\right|=+\mathrm{}}{Res}g(z,t)=\underset{r+\mathrm{}}{lim}\left[\underset{G_r^R}{\overset{}{}}g(z,t)dz+\underset{G_r^L}{\overset{}{}}g(z,t)dz\right]`$.
#### 2.2.2. Cauchy’s formula
As it is well-known, during the deriving Caushy’s formula for expansion of real valued functions into an infinite functional series, and (taken over from ) - Formula (9), Subsection 4.6.2, Section 4.6, Chapter 4, p. 94, \- in a real axis interval $`(t_0,t_1)`$
$$f\left(t\right)=\underset{k=1}{\overset{+\mathrm{}}{}}\frac{w\left(a_k\right)}{\frac{dq\left(z\right)}{dz}|_{z=a_k}}\underset{t_0}{\overset{t_1}{}}e^{a_k\left(t\tau \right)}f\left(\tau \right)d\tau \text{;}\text{ }tt_{si}\text{,}$$
(2.11)
where $`t_{si}`$ are break points of the function $`tf\left(t\right)`$ in $`(t_0,t_1)`$, the conditions for existence of finite limiting values of the functional expressions:
$$\underset{\left|z\right|+\mathrm{}}{lim}\left[z\frac{p\left(z\right)}{q\left(z\right)}\underset{t_0}{\overset{𝑡}{}}e^{z\left(t\tau \right)}f\left(\tau \right)d\tau z\frac{p\left(z\right)}{q\left(z\right)}\underset{t_0}{\overset{𝑡}{}}e^{z\left(t\tau \right)}f\left(\tau \right)d\tau \right]\text{;}$$
(2.12)
$$\underset{\left|z\right|+\mathrm{}}{lim}\left[z\frac{w\left(z\right)}{q\left(z\right)}\underset{𝑡}{\overset{t_1}{}}e^{z\left(t\tau \right)}f\left(\tau \right)d\tau z\frac{w\left(z\right)}{q\left(z\right)}\underset{𝑡}{\overset{t_1}{}}e^{z\left(t\tau \right)}f\left(\tau \right)d\tau \right]\text{,}$$
(2.13)
which have to be satisfied by the function $`tf\left(t\right)`$ in $`(t_0,t_1)`$, are of the most importance.
Namely, let $`g_1(z,t)=\frac{p\left(z\right)}{q\left(z\right)}\underset{t_0}{\overset{𝑡}{}}e^{z\left(t\tau \right)}f\left(\tau \right)d\tau `$ and $`g_2(z,t)=\frac{w\left(z\right)}{q\left(z\right)}\underset{𝑡}{\overset{t_1}{}}e^{z\left(t\tau \right)}f\left(\tau \right)d\tau `$, where an analytic function $`q\left(z\right)`$: $`q\left(z\right)=p\left(z\right)+w\left(z\right)`$, has infinitely but a count of many simple poles: $`a_1,a_2,\mathrm{}`$ onto the imaginary axis as singularities. If under an assumption that: $`\underset{\left|z\right|+\mathrm{}}{lim}\frac{p\left(z\right)}{q\left(z\right)}e^{z\left(tt_0\right)}=0`$ and $`\underset{\left|z\right|+\mathrm{}}{lim}\frac{p\left(z\right)}{q\left(z\right)}=1`$ as well as $`\underset{\left|z\right|+\mathrm{}}{lim}\frac{w\left(z\right)}{q\left(z\right)}e^{z\left(t_1t\right)}=0`$ and $`\underset{\left|z\right|+\mathrm{}}{lim}\frac{w\left(z\right)}{q\left(z\right)}=1`$, the following functional expressions: $`\underset{\left|z\right|+\mathrm{}}{lim}z\underset{t_0}{\overset{𝑡}{}}e^{z\left(\tau t_0\right)}f\left(\tau \right)d\tau `$ and$`\underset{\left|z\right|+\mathrm{}}{lim}z\underset{𝑡}{\overset{t_1}{}}e^{z\left(t_1\tau \right)}f\left(\tau \right)d\tau `$, converge for each $`\mathrm{R}ez>0`$ and $`t`$ $`(t_0,t_1)`$, just as well as functional expressions: $`\underset{\left|z\right|+\mathrm{}}{lim}z\underset{t_0}{\overset{𝑡}{}}e^{z\left(t\tau \right)}f\left(\tau \right)d\tau `$ and$`\underset{\left|z\right|+\mathrm{}}{lim}z\underset{𝑡}{\overset{t_1}{}}e^{z\left(\tau t\right)}f\left(\tau \right)d\tau `$, in such a way that: $`\underset{\left|z\right|+\mathrm{}}{lim}z\underset{t_0}{\overset{𝑡}{}}e^{z\left(t\tau \right)}f\left(\tau \right)d\tau =f\left(t\right)`$ and$`\underset{\left|z\right|+\mathrm{}}{lim}z\underset{𝑡}{\overset{t_1}{}}e^{z\left(\tau t\right)}f\left(\tau \right)d\tau =f\left(t\right)`$; $`tt_{si}`$, then on the basis of previous analyzed idea it can be proved that at all points of $`(t_0,t_1)`$, at which a function $`tf\left(t\right)`$ is continuous, there exists a sum of an infinite functional series on the right-hand side of the equation (2.11) which is just equal to the function values at those points, more exactly, in the general case, Cauchy’s infinite functional series of the function $`tf\left(t\right)`$ is summable.
If a function $`tf\left(t\right)`$ satisfies the general well-known Dirichlet’s conditions in $`(t_0,t_1)`$, then partial sums of Cauchy’s infinite functional series of the function $`tf\left(t\right)`$ at all points of $`(t_0,t_1)`$, at which the function $`tf\left(t\right)`$ is continuous, converge to the function values at those points . At the break points $`t_{si}`$ of the function $`tf\left(t\right)`$ in $`(t_0,t_1)`$ the partial sums of Cauchy’s infinite functional series of the function $`tf\left(t\right)`$ converge to the following functional values
$$\frac{1}{2}\left[\underset{\epsilon 0^+}{lim}f\left(t_{si}+\epsilon \right)+\underset{\eta 0^+}{lim}f\left(t_{si}\eta \right)\right]\text{.}$$
(2.14)
At the extreme points of the segment $`[t_0,t_1]`$: $`t_0`$ and $`t_1`$, at which a function $`tf\left(t\right)`$ is continuous on one’s right and left respectively, the sum of Cauchy’s infinite functional series of the function $`tf\left(t\right)`$ is equal to the following functional value
$$\frac{1}{2}\left[\underset{\epsilon 0^+}{lim}f\left(t_0+\epsilon \right)+\underset{\eta 0^+}{lim}f\left(t_1\eta \right)\right]\text{.}$$
(2.15)
#### 2.2.3. Interval of improper integrals convergence of real valued functions
Let $`\nu f\left(\nu \right)`$, be an analytic function of complex variable $`\nu `$ on some neighborhood $`V_0`$ of the point $`\nu =0`$ at which the function $`\nu f\left(\nu \right)`$ has a pole of arbitrary order as a singularity. A function $`\nu f\left(\nu \right)e^{z\left(t\nu \right)}`$, where $`zC^1`$ is a complex parameter and $`tR_+^1`$ ($`\mathrm{R}e\nu =t`$) is fixed point belonging to the neighborhood $`V_0`$, is parametric analytic function on $`V_0`$. Further, a smooth one-one mapping $`\nu \left(\theta \right)`$: $`R^1C^1`$ ($`\nu \left(\theta \right)=\frac{tt_0}{2}e^{i\theta }+\frac{t_0+t}{2}`$) of the real axis segment $`[\pi ,0]`$ ($`\theta [\pi ,0]`$) onto the set of complex points $`\nu `$ of the complex plane $`C^1`$ is defined. A fixed point $`\mathrm{R}e\nu =t_0`$ ($`t_0<0`$) also belongs to the neighborhood $`V_0`$ of the zero point $`\nu =0`$. An arbitrary $`n`$-division $`P_n`$: $`P_n=\left\{\theta _0=\pi ,\theta _1,\mathrm{},\theta _i,\mathrm{},\theta _n=0\right\}`$, where $`nN`$, is one of all possible $`n`$-divisions of the segment $`[\pi ,0]`$.
Accordingly, since a complex function $`f\left[\nu \left(\theta \right)\right]e^{z\left[t\nu \left(\theta \right)\right]}`$ of a real variable $`\theta `$ is a regular that in the segment $`[\pi ,0]`$, and in those circumstances its both a real and an imaginary part satisfies all conditions of Langrange’s mean value theorem of the differential calculus in the segment $`[\pi ,0]`$, then, for each partial segment $`[\theta _{i1},\theta _i]`$ of the segment $`[\pi ,0]`$, it holds
$$\mathrm{R}e\left\{\left\{\frac{d}{d\theta }\left\{f\left[\nu \left(\theta \right)\right]e^{z\left[t\nu \left(\theta \right)\right]}\right\}\right\}\right|{}_{\theta =\theta _i^{}}{}^{}\}=$$
(2.16)
$$=\frac{\mathrm{R}e\left\{f\left[\nu \left(\theta _i\right)\right]e^{z\left[t\nu \left(\theta _i\right)\right]}\right\}\mathrm{R}e\left\{f\left[\nu \left(\theta _{i1}\right)\right]e^{z\left[t\nu \left(\theta _{i1}\right)\right]}\right\}}{\theta _i\theta _{i1}}\text{;}$$
$$\mathrm{I}m\left\{\left\{\frac{d}{d\theta }\left\{f\left[\nu \left(\theta \right)\right]e^{z\left[t\nu \left(\theta \right)\right]}\right\}\right\}\right|{}_{\theta =\theta _i^{}}{}^{}\}=$$
(2.17)
$$=\frac{\mathrm{I}m\left\{f\left[\nu \left(\theta _i\right)\right]e^{z\left[t\nu \left(\theta _i\right)\right]}\right\}\mathrm{I}m\left\{f\left[\nu \left(\theta _{i1}\right)\right]e^{z\left[t\nu \left(\theta _{i1}\right)\right]}\right\}}{\theta _i\theta _{i1}}\text{,}$$
where $`\{\theta _i^{},\theta _i^{}\}[\theta _{i1},\theta _i]`$.
By virtue of (2.16) and (2.17), it is possible to form the integral sums
$$\underset{i=1}{\overset{𝑛}{}}\mathrm{R}e\left\{\left\{\frac{d}{d\theta }\left\{f\left[\nu \left(\theta \right)\right]e^{z\left[t\nu \left(\theta \right)\right]}\right\}\right\}\right|{}_{\theta =\theta _i^{}}{}^{}\}(\theta _i\theta _{i1})=$$
(2.18)
$$=\mathrm{R}e\left[f\left(t\right)\right]\mathrm{R}e\left[f\left(t_0\right)e^{z\left(tt_0\right)}\right]\text{;}$$
$$\underset{i=1}{\overset{𝑛}{}}\mathrm{I}m\left\{\left\{\frac{d}{d\theta }\left\{f\left[\nu \left(\theta \right)\right]e^{z\left[t\nu \left(\theta \right)\right]}\right\}\right\}\right|{}_{\theta =\theta _i^{}}{}^{}\}(\theta _i\theta _{i1})=$$
(2.19)
$$=\mathrm{I}m\left[f\left(t\right)\right]\mathrm{I}m\left[f\left(t_0\right)e^{z\left(tt_0\right)}\right]\text{,}$$
that is, after the performed differentiation
$$\underset{i=1}{\overset{𝑛}{}}\mathrm{R}e\left\{\left\{\frac{d}{d\theta }f\left[\nu \left(\theta \right)\right]e^{z\left[t\nu \left(\theta \right)\right]}\right\}\right|{}_{\theta =\theta _i^{}}{}^{}\}(\theta _i\theta _{i1})+$$
(2.20)
$$+\underset{i=1}{\overset{𝑛}{}}\mathrm{R}e\left\{\left\{zf\left[\nu \left(\theta \right)\right]e^{z\left[t\nu \left(\theta \right)\right]}\right\}\right|{}_{\theta =\theta _i^{}}{}^{}\}(\theta _i\theta _{i1})=$$
$$=\mathrm{R}e\left[f\left(t\right)\right]\mathrm{R}e\left[f\left(t_0\right)e^{z\left(tt_0\right)}\right]\text{;}$$
$$\underset{i=1}{\overset{𝑛}{}}\mathrm{I}m\left\{\left\{\frac{d}{d\theta }f\left[\nu \left(\theta \right)\right]e^{z\left[t\nu \left(\theta \right)\right]}\right\}\right|{}_{\theta =\theta _i^{}}{}^{}\}(\theta _i\theta _{i1})+$$
(2.21)
$$+\underset{i=1}{\overset{𝑛}{}}\mathrm{I}m\left\{\left\{zf\left[\nu \left(\theta \right)\right]e^{z\left[t\nu \left(\theta \right)\right]}\right\}\right|{}_{\theta =\theta _i^{}}{}^{}\}(\theta _i\theta _{i1})=$$
$$=\mathrm{I}m\left[f\left(t\right)\right]\mathrm{I}m\left[f\left(t_0\right)e^{z\left(tt_0\right)}\right]\text{.}$$
Having in view the fact that $`\nu \left(\theta \right)=\frac{tt_0}{2}e^{i\theta }+\frac{t_0+t}{2}`$, it is possible to define an interval (a semi-interval) of a change of the argument $`\phi `$: $`\phi R^1`$, of the complex parameter $`z`$, for which it holds: $`\underset{\left|z\right|+\mathrm{}}{lim}e^{z\left[t\nu \left(\theta \right)\right]}=0`$; $`\theta (\pi ,0)`$. Namely, since for $`\mathrm{R}ez0`$
$$z\left[t\nu \left(\theta \right)\right]=z\left(\frac{tt_0}{2}\right)\left(1e^{i\theta }\right)=$$
(2.22)
$$=\left(\frac{tt_0}{2}\right)\left[\mathrm{R}ez\left(1\mathrm{cos}\theta \right)\right]\left(1+\frac{\mathrm{I}mz}{\mathrm{R}ez}\frac{\mathrm{sin}\theta }{1\mathrm{cos}\theta }\right)+$$
$$+i\left[\mathrm{I}mz\left(1\mathrm{cos}\theta \right)\mathrm{R}ez\mathrm{sin}\theta \right]\text{,}$$
then, if the condition
$$\mathrm{R}ez\left(1\mathrm{cos}\theta \right)\left(1+\frac{\mathrm{I}mz}{\mathrm{R}ez}\frac{\mathrm{sin}\theta }{1\mathrm{cos}\theta }\right)>0$$
(2.23)
is satisfied, it follows that $`\underset{\left|z\right|+\mathrm{}}{lim}e^{z\left[t\nu \left(\theta \right)\right]}=0`$; $`\theta (\pi ,0)`$.
In view of the fact that $`\frac{\mathrm{sin}\theta }{1\mathrm{cos}\theta }=\mathrm{cot}\frac{\theta }{2}=\left(\mathrm{tan}\frac{\theta }{2}\right)^1`$ and $`\frac{\mathrm{I}mz}{\mathrm{R}ez}=\mathrm{tan}\phi `$, the condition (2.23) is satisfied if and only if $`\phi (\frac{\pi }{2},0]`$, in other words for each $`\phi (\frac{\pi }{2},0]`$ it holds $`\underset{\left|z\right|+\mathrm{}}{lim}e^{z\left[t\nu \left(\theta \right)\right]}=0`$; $`\theta (\pi ,0)`$.
Hence, and on the basis of derived relations: (2.20) and (2.21), in the limit as $`n+\mathrm{}`$, more exactly, when a maximum partial segment $`[\theta _{i1},\theta _i]`$ of the segment $`[\pi ,0]`$ vanishes, the condition $`\phi (\frac{\pi }{2},0]`$ becomes a condition of convergence of a parametric contour integral
$$\underset{\left|z\right|+\mathrm{}}{lim}z\underset{𝐺}{}f\left(\nu \right)e^{z\left(t\nu \right)}d\nu =f\left(t\right)\text{,}$$
(2.24)
where $`G=\{\nu \text{ }\nu =\frac{tt_0}{2}e^{i\theta }+\frac{t_0+t}{2}\}`$ and $`\theta [\pi ,0]`$.
In other words, the semi-interval $`(\frac{\pi }{2},0]`$ is a semi-interval of a complex parametric contour integral convergence of the function $`f\left(\nu \right)e^{z\left(t\nu \right)}`$ along the given contour of integration $`G`$, in the limit as $`\left|z\right|+\mathrm{}`$.
In the next step, instead of the above integration path $`G`$, a complex plane curve $`G^{}`$ consisting of parts of real axis defined by segments: $`[t_0,\epsilon ]`$ and $`[\epsilon ,t]`$ ($`\epsilon R_+^1`$) as well as of a part of a circular path defined by a smooth one-one mapping $`\nu \left(\theta \right)`$: $`R^1C^1`$ ($`\nu \left(\theta \right)=\epsilon e^{i\theta }`$) of the segment $`[\pi ,0]`$ of a real axis onto a set of points of the complex plane, is taken for a contour of integration. The integral $`\underset{G^{}}{}f\left(\nu \right)e^{z\left(t\nu \right)}d\nu `$, defined over the integration contour $`G_\epsilon ^{}=\{\nu \text{ }\nu \left(\theta \right)=\epsilon e^{i\theta }\}`$ bypassing a singularity of a function $`f\left(\nu \right)e^{z\left(t\nu \right)}`$ at the point $`\nu =0`$, is a by-pass integral. In this case too, similarly to the previous analysis, it is possible to define an interval (a semi-interval) of a parametric contour integral convergence of the function $`f\left(\nu \right)e^{z\left(t\nu \right)}`$. Namely, since in this case
$$z\left[t\nu \left(\theta \right)\right]=z\left[t\epsilon \left(\mathrm{cos}\theta +i\mathrm{sin}\theta \right)\right]=$$
(2.25)
$$=\mathrm{R}ez\left(t\epsilon \mathrm{cos}\theta \right)+\epsilon \mathrm{I}mz\mathrm{sin}\theta +$$
$$+i\left[\mathrm{I}mz\left(t\epsilon \mathrm{cos}\theta \right)\epsilon \mathrm{R}ez\mathrm{sin}\theta \right]\text{,}$$
then the condition (2.23), for the contour of integration $`G`$, reduces to the condition
$$\mathrm{R}ez\left[\left(t\epsilon \mathrm{cos}\theta \right)+\epsilon \frac{\mathrm{I}mz}{\mathrm{R}ez}\mathrm{sin}\theta \right]>0\text{}\mathrm{R}ez>0\text{,}$$
(2.26)
for the contour of integration $`G_\epsilon ^{}`$.
In view of the fact that $`t>\epsilon >0`$, there exists a positive real number $`k`$: $`kR_+^1`$, such that $`t=\left(1+k\right)\epsilon `$. Hence, the condition (2.26) reduces to the condition
$$\epsilon \mathrm{R}ez\left[k+\left(1\mathrm{cos}\theta \right)+\mathrm{tan}\phi \mathrm{sin}\theta \right]>0\text{}\mathrm{R}ez>0\text{.}$$
(2.27)
As $`\left|\mathrm{sin}\theta \right|1`$, for $`\theta (\pi ,0)`$, the condition (2.27) is satisfied if and only if $`\phi (\frac{\pi }{2},\mathrm{arctan}k]`$, in other words for each $`\phi (\frac{\pi }{2},\mathrm{arctan}k]`$ it holds
$$\underset{\left|z\right|+\mathrm{}}{lim}z\underset{G_\epsilon ^{}}{}f\left(\nu \right)e^{z\left(t\nu \right)}d\nu =0\text{.}$$
(2.28)
On the other hand, since the real meromorphic function $`f\left[\left(\mathrm{R}e\nu \right)\right]`$ satisfies in the semi-segment $`[t_0,0)`$, as well as in the semi-interval $`(0,t]`$, general well-known Dirichlet’s conditions, and , then, for each $`\mathrm{R}ez0`$, it holds:
$$\underset{\left|z\right|+\mathrm{}}{lim}z\underset{t_0}{\overset{\epsilon }{}}f\left(\tau \right)e^{z\left(t\tau \right)}d\tau =0\text{ and}\underset{\left|z\right|+\mathrm{}}{lim}z\underset{𝜀}{\overset{𝑡}{}}f\left(\tau \right)e^{z\left(t\tau \right)}d\tau =f\left(t\right)\text{,}$$
respectively, where $`\tau =\mathrm{R}e\nu `$. Therefore, the semi-interval $`(\frac{\pi }{2},\mathrm{arctan}k]`$ is a semi-interval of convergence of the parametric contour integral
$$\underset{\left|z\right|+\mathrm{}}{lim}z\underset{G^{}}{}f\left(\nu \right)e^{z\left(t\nu \right)}d\nu =f\left(t\right)\text{.}$$
(2.29)
Taking into consideration the fact that in the limit as $`\epsilon 0^+`$ and for each $`t>0`$: $`k+\mathrm{}`$ ( $`t=\left(1+k\right)\epsilon `$), the interval $`(\frac{\pi }{2},\frac{\pi }{2})`$ ($`\phi (\frac{\pi }{2},\frac{\pi }{2})`$) becomes an interval of convergence of an improper integral
$$\underset{\left|z\right|+\mathrm{}}{lim}z[v.t.\underset{t_0}{\overset{𝑡}{}}f\left(\tau \right)e^{z\left(t\tau \right)}d\tau ]=f\left(t\right)$$
(2.30)
absolutely existing in the segment $`[t_0,t]`$; $`t>0`$. In other words, for each $`\mathrm{R}ez>0`$, it holds (2.30).
In the similar manner it can be proved to be
$$\underset{\left|z\right|+\mathrm{}}{lim}z[v.t.\underset{t_0}{\overset{𝑡}{}}f\left(\tau \right)e^{z\left(\tau t_0\right)}d\tau ]=f\left(t_0\right)\text{,}$$
(2.31)
for each $`\mathrm{R}ez>0`$ and $`t>0`$ ($`t(t_0,t_1)`$.
It should be emphasized that as distinguished from limiting values of contour integrals:
$$\underset{\left|z\right|+\mathrm{}}{lim}z\underset{𝐺}{}f\left(\nu \right)e^{z\left(t\nu \right)}d\nu \text{ and}\underset{\left|z\right|+\mathrm{}}{lim}z\underset{G^{}}{}f\left(\nu \right)e^{z\left(t\nu \right)}d\nu \text{,}$$
which are equal, their intervals of convergence are different.$`\mathrm{}`$
#### 2.2.4. Fourier formula
Based on the results obtained by previous analysis one may say that in addition to a class of real valued functions satisfying so-called general Dirichlet’s conditions in the real axis segment $`[t_0,t_1]`$, for which the functional expressions: (2.12) and (2.13), converge for each $`\mathrm{R}ez0`$, , there exists an one other class of real valued functions for which Cauchy’s formula is still in effect, a class of the real valued meromorphic functions whose finitely many isolated singularities lie onto the segment $`[t_0,t_1]`$.
From Cauchy’s formula, and for $`p\left(z\right)=1`$ and $`w\left(z\right)=e^{az}`$ ($`aR_+^1`$) i.e. $`q\left(z\right)=e^{az}1`$, it is immediately obtained that for $`t_1a<t<t_0+a`$ and $`tt_{si}`$
$$f\left(t\right)=\frac{1}{a}[v.t.\underset{t_0}{\overset{t_1}{}}f\left(\tau \right)d\tau ]+\frac{2}{a}\underset{k=1}{\overset{+\mathrm{}}{}}[v.t.\underset{t_0}{\overset{t_1}{}}f\left(\tau \right)\mathrm{cos}\frac{2k\pi \left(t\tau \right)}{a}d\tau ]\text{.}$$
(2.32)
If $`a=2\pi `$, $`t_0=\pi `$ and $`t_1=\pi `$, then the equation (2.32) represents an expansion of a real valued meromorphic function $`tf\left(t\right)`$ into a Fourier trigonometric series in the interval $`(\pi ,\pi )`$, more exactly, for each $`\pi <t<\pi `$ and $`tt_{si}`$, where $`t_{si}`$ are break points of the function $`tf\left(t\right)`$ in the interval $`(\pi ,\pi )`$, it holds
$$f\left(t\right)=\frac{1}{2}A_0+\underset{k=1}{\overset{+\mathrm{}}{}}\left[A_k\mathrm{cos}\left(kt\right)+B_k\mathrm{sin}\left(kt\right)\right]\text{,}$$
(2.33)
where
$$A_k=\frac{1}{\pi }[v.t.\underset{\pi }{\overset{𝜋}{}}f\left(\tau \right)\mathrm{cos}\left(k\tau \right)d\tau ]\text{;}\text{ }kN_0$$
(2.34)
and
$$B_k=\frac{1}{\pi }[v.t.\underset{\pi }{\overset{𝜋}{}}f\left(\tau \right)\mathrm{sin}\left(k\tau \right)d\tau ]\text{;}\text{ }kN\text{.}$$
(2.35)
According to the well-known result of Dirichlet’s theorem, see - Theorem 1, Section 5.4, Chapter 5, p. 65, \- in the general case of a class of real valued functions $`tf_d\left(t\right)`$ satisfying the general Dirichlet’s conditions in the segment $`[\pi ,\pi ]`$, the Fourier trigonometric series on the right-hand side of the equation (2.33) can be said to converge to a function $`F_d\left(t\right)`$. Clearly, at all points of the interval $`(\pi ,\pi )`$ at which the function $`tf_d\left(t\right)`$ is continuous, a convergent value of Fourier series is equal to the function value: $`F_d\left(t\right)=f_d\left(t\right)`$. Considering the consequence of whether Bessel’s inequality or Riemann-Lebesque’s theorem - Theorem 2, Section 6.2, Chapter 5, p. 96, \- Fourier’s coefficients of the function $`tf_d\left(t\right)`$: $`A_k`$ and $`B_k`$, tend to zero as $`k+\mathrm{}`$. This is important from the viewpoint of the convergence of infinite numerical series obtained by expansion of functions $`tf_d\left(t\right)`$ into Fourier trigonometric series. Note that the conditions of Dirichlet’s theorem are only sufficient conditions for convergence of Fourier trigonometric series of functions $`tf_d\left(t\right)`$.
A nature of Fourier trigonometric series convergence of a class of real valued meromorphic functions $`tf_m\left(t\right)`$, whose finitely many isolated singularities lie onto the segment $`[\pi ,\pi ]`$, can be said to be different from a case to a case. The same holds also for Fourier’s coefficients of a function $`tf_m\left(t\right)`$ in the limit as $`k+\mathrm{}`$. Namely, in the general case of real meromorphic functions $`tf_m\left(t\right)`$, at all points of the interval $`(\pi ,\pi )`$ at which a function $`tf_m\left(t\right)`$ is continuous, Fourier trigonometric series of the function $`tf_m\left(t\right)`$ is summable, more exactly it has defined sum $`F_m\left(t\right)`$: $`F_m\left(t\right)=f_m\left(t\right)`$. The concept of the sum of Fourier trigonometric series, in this case, is generalization of the concept of the sum in the Caushy’s sense. At the break points $`t_{si}`$ of the function $`tf_m\left(t\right)`$ in the segment $`[\pi ,\pi ]`$, the sum of Fourier trigonometric series of the real valued meromorphic function $`tf_m\left(t\right)`$, in the general case, is not defined.$`\mathrm{}`$
## 3. Examples
### 3.1. Example 1
An expansion of the function $`t\frac{1}{2}\frac{\mathrm{sin}t}{1\mathrm{cos}t}`$ into a Fourier trigonometric series in the segment $`[\pi ,\pi ]`$. The function $`f\left(t\right)=\frac{1}{2}\frac{\mathrm{sin}t}{1\mathrm{cos}t}`$ having at the point a simple pole as a singularity is a real valued meromorphic function in the segment $`[\pi ,\pi ]`$. Caushy’s principal value ($`v.p.`$) of an improper integral of the function $`f\left(t\right)`$ is equal to:
$$v.p.\underset{\pi }{\overset{𝜋}{}}f\left(t\right)dt=\frac{1}{2}\underset{\epsilon 0^+}{lim}\left[\underset{\pi }{\overset{\epsilon }{}}\frac{\mathrm{sin}t}{1\mathrm{cos}t}dt+\underset{𝜀}{\overset{𝜋}{}}\frac{\mathrm{sin}t}{1\mathrm{cos}t}dt\right]=0\text{.}$$
(3.1)
As $`\underset{z0}{lim}\frac{1}{2}\frac{z\mathrm{sin}z}{1\mathrm{cos}z}=1`$; $`zC^1`$, the by-pass integral value in the limit as $`\epsilon 0^+`$ (Jordan’s singular value ($`v.s.`$) of the improper integral) is equal to:
$$\underset{\epsilon 0^+}{lim}\underset{G_{\epsilon _\kappa }}{}\frac{1}{2}\frac{\mathrm{sin}z}{1\mathrm{cos}z}dz=\{\begin{array}{c}i\pi \text{;}\text{ }\kappa =1\hfill \\ i\pi \text{;}\text{ }\kappa =2\hfill \end{array}\text{,}$$
(3.2)
dependently on the choice of a circular arc $`G_{\epsilon _\kappa }`$ bypassing the singularity of the function $`f\left(t\right)`$ in the complex plane: $`G_{\epsilon _\kappa }=\{zz=\epsilon e^{i\theta _\kappa }\text{;}\text{ }\theta _\kappa \left\{\begin{array}{c}[\pi ,0]\text{;}\text{ }\kappa =1\hfill \\ [\pi ,0]\text{;}\text{ }\kappa =2\hfill \end{array}\right\}`$.
The total value ($`v.t.`$) of the improper integral, as a sum of Caushy’s principal value ($`v.p.`$) and Jordan’s singular value ($`v.s.`$):
$$\frac{1}{\pi }v.t.\underset{\pi }{\overset{𝜋}{}}f\left(t\right)dt=\frac{1}{2\pi }v.t.\underset{\pi }{\overset{𝜋}{}}\frac{\mathrm{sin}t}{1\mathrm{cos}t}dt=\{\begin{array}{c}i\hfill \\ i\hfill \end{array}$$
(3.3)
absolutely exists in this case and as one can see is not unique.
On the other hand, since $`\frac{1}{\pi }\underset{0}{\overset{𝜋}{}}\frac{\mathrm{sin}\left[\left(k+\frac{1}{2}\right)t\right]}{\mathrm{sin}\left(\frac{1}{2}t\right)}dt=1`$; $`kN`$ \- Formula (6) Section 6.2, Chapter 6, p. 95, \- and $`\frac{1}{2}\underset{\pi }{\overset{𝜋}{}}\mathrm{cos}\left(kt\right)dt=0`$; $`kN`$, as well as $`\underset{z0}{lim}z\frac{\mathrm{sin}z\mathrm{cos}\left(kz\right)}{2\left(1\mathrm{cos}z\right)}=1`$, then it follows that
$$B_k=\frac{1}{\pi }v.t.\underset{\pi }{\overset{𝜋}{}}\frac{\mathrm{sin}t\mathrm{sin}\left(kt\right)}{2\left(1\mathrm{cos}t\right)}dt=\frac{1}{2\pi }v.t.\underset{\pi }{\overset{𝜋}{}}\mathrm{cot}\frac{t}{2}\mathrm{sin}\left(kt\right)dt=1\text{}kN\text{,}$$
(3.4)
as well as
$$A_k=\frac{1}{\pi }v.t.\underset{\pi }{\overset{𝜋}{}}\frac{\mathrm{sin}t\mathrm{cos}\left(kt\right)}{2\left(1\mathrm{cos}t\right)}dt=\{\begin{array}{c}i\hfill \\ i\hfill \end{array}\text{;}\text{ }kN\text{,}$$
(3.5)
Accordingly, and by the Fourier formula (2.33), a Fourier trigonometric series of the function $`t\frac{1}{2}\frac{\mathrm{sin}t}{1\mathrm{cos}t}`$ can be expressed by the following functional form
$$\frac{1}{2}\frac{\mathrm{sin}t}{1\mathrm{cos}t}=\underset{k=1}{\overset{+\mathrm{}}{}}\mathrm{sin}\left(kt\right)\frac{i}{2}\left[1+2\underset{k=1}{\overset{+\mathrm{}}{}}\mathrm{cos}\left(kt\right)\right]\text{,}$$
(3.6)
that is, the equalities
$$\frac{1}{2}\frac{\mathrm{sin}t}{1\mathrm{cos}t}=\underset{k=1}{\overset{+\mathrm{}}{}}\mathrm{sin}\left(kt\right)\text{ and }1+2\underset{k=1}{\overset{+\mathrm{}}{}}\mathrm{cos}\left(kt\right)=0\text{ i.e.}$$
(3.7)
$$\frac{1}{2}\frac{1e^{\pm it}}{1\mathrm{cos}t}=\underset{k=1}{\overset{+\mathrm{}}{}}e^{\pm ikt}$$
hold for each $`t(\pi ,\pi )`$ and $`t0`$, respectively.
By relation (2.15), for $`t=\pm \pi `$, it follows that
$$\underset{k=1}{\overset{+\mathrm{}}{}}\mathrm{sin}\left(k\pi \right)=0\text{ and }1+2\underset{k=1}{\overset{+\mathrm{}}{}}\mathrm{cos}\left(k\pi \right)=1+2\underset{k=1}{\overset{+\mathrm{}}{}}\left(1\right)^k=0\text{.}\mathrm{}$$
(3.8)
### 3.2. Example 2
An expansion of the function $`t\frac{1}{2}\frac{1}{1\mathrm{cos}t}`$ into a Fourier trigonometric series in the segment $`[\pi ,\pi ]`$. The real valued meromorphic function $`f\left(t\right)=\frac{1}{2}\frac{1}{1\mathrm{cos}t}`$ in the segment $`[\pi ,\pi ]`$ has the second order pole at the point $`t=0`$ as a singularity. The improper integral $`\frac{1}{2}\underset{\pi }{\overset{𝜋}{}}\frac{dt}{1\mathrm{cos}t}`$ absolutely exists and reduces to the indefinite expression of difference of infinities $`\mathrm{}\mathrm{}`$. Namely, independently on the choice of the circular arc bypassing the singularity $`z=0`$ of the function $`z\frac{1}{2}\frac{1}{1\mathrm{cos}z}`$ in the complex plane (for example $`G_\epsilon =\{zz=\epsilon e^{i\theta }\text{;}\text{ }\theta [\pi ,0]\}`$), on the one hand it holds
$$\frac{1}{2}v.t.\underset{\pi }{\overset{𝜋}{}}\frac{dt}{1\mathrm{cos}t}=\frac{1}{2}\underset{\epsilon 0^+}{lim}\left[\underset{\pi }{\overset{\epsilon }{}}\frac{dt}{1\mathrm{cos}t}+\underset{G_\epsilon }{}\frac{dz}{1\mathrm{cos}z}+\underset{𝜀}{\overset{𝜋}{}}\frac{dt}{1\mathrm{cos}t}\right]=$$
(3.9)
$$=\underset{\epsilon 0^+}{lim}\left[\frac{\mathrm{sin}\epsilon }{1\mathrm{cos}\epsilon }+\frac{1}{2}\underset{G_\epsilon }{}\frac{dz}{1\mathrm{cos}z}\right]\text{.}$$
On the other hand, since - see Definition 4, Section 2.2, Chapter 2, p. 82,
$$\underset{G_\epsilon }{}\frac{dz}{1\mathrm{cos}z}=\underset{\pi }{\overset{0}{}}\frac{\frac{d}{d\theta }\left[z\left(\theta \right)\right]}{1\mathrm{cos}\left[z\left(\theta \right)\right]}d\theta =\underset{\pi }{\overset{0}{}}\frac{i\epsilon e^{i\theta }}{1\mathrm{cos}\left(\epsilon e^{i\theta }\right)}d\theta \text{,}$$
(3.10)
that is
$$\underset{G_\epsilon }{}\frac{dz}{1\mathrm{cos}z}=\frac{\mathrm{sin}\left[z\left(\theta \right)\right]}{1\mathrm{cos}\left[z\left(\theta \right)\right]}|{}_{\pi }{}^{0}=\frac{\mathrm{sin}\left(\epsilon e^{i\theta }\right)}{1\mathrm{cos}\left(\epsilon e^{i\theta }\right)}|{}_{\pi }{}^{0}=\frac{2\mathrm{sin}\epsilon }{1\mathrm{cos}\epsilon }\text{,}$$
(3.11)
then finally it follows that the total value ($`v.t.`$) of the improper integral $`\underset{\pi }{\overset{𝜋}{}}\frac{dt}{2\left(1\mathrm{cos}t\right)}`$ is equal to the value zero:
$$v.t.\underset{\pi }{\overset{𝜋}{}}\frac{dt}{2\left(1\mathrm{cos}t\right)}=0\text{.}$$
(3.12)
The complex function $`z\frac{z^k}{z1}`$; $`kN`$ of complex variable $`z`$ is a meromorphic function having at the point $`z=1`$ a simple pole as singularity. Since $`\underset{z1}{lim}\left[\left(z1\right)\frac{z^k}{z1}\right]=1`$ then, according to the result (2.1) in the Section 2.1 of the paper, Cauchy’s principle value ($`v.p.`$) of the improper integral $`\underset{𝐺}{}\frac{z^k}{z1}dz`$ over the circular contour of integration $`G`$: $`G=\{zz=\epsilon e^{i\theta }\text{;}\text{ }\theta [\pi ,\pi ]\}`$ is equal to: $`v.p.\underset{𝐺}{}\frac{z^k}{z1}dz=i\pi `$.
With regard to the fact that $`z=\epsilon e^{i\theta }`$ onto the integration contour $`G`$, it follows that
$$v.p.\underset{𝐺}{}\frac{z^k}{z1}dz=v.p.\underset{\pi }{\overset{𝜋}{}}\frac{ie^{ik\theta }}{1e^{i\theta }}d\theta =v.p.[\underset{\pi }{\overset{𝜋}{}}\frac{i\mathrm{cos}\left(k\theta \right)}{2\left(1\mathrm{cos}\theta \right)}d\theta $$
(3.13)
$$\underset{\pi }{\overset{𝜋}{}}\frac{i\mathrm{cos}\left[\left(k+1\right)\theta \right]}{2\left(1\mathrm{cos}\theta \right)}d\theta ]=i\pi \text{,}$$
since $`v.p.\underset{\pi }{\overset{𝜋}{}}\frac{\mathrm{cos}\left(k\theta \right)\mathrm{sin}\theta }{1\mathrm{cos}\theta }d\theta =0`$ and $`\underset{\pi }{\overset{𝜋}{}}\mathrm{sin}\left(k\theta \right)d\theta =0`$; for each $`kN`$.
On the other hand, for $`kN`$: $`\underset{z0}{lim}\left[z\mathrm{cos}\left(kz\right)\right]=0`$ and $`\underset{z0}{lim}\frac{z\mathrm{sin}\left(kz\right)\mathrm{sin}z}{1\mathrm{cos}z}=0`$. According to the result of Jordan’s lemma - Theorem 1, Subsection 3.1.4, Section 3.1, Chapter 3, p. 52, \- it holds
$$\frac{1}{2\pi i}\underset{\epsilon 0^+}{lim}\underset{G_{\epsilon _\kappa }}{}\left\{\frac{\mathrm{cos}\left(kz\right)\mathrm{cos}\left[\left(k+1\right)z\right]}{2\left(1\mathrm{cos}z\right)}\right\}dz=$$
(3.14)
$$=\frac{1}{4\pi i}\underset{\epsilon 0^+}{lim}\underset{G_{\epsilon _\kappa }}{}\left[\mathrm{cos}\left(kz\right)\frac{\mathrm{sin}\left(kz\right)\mathrm{sin}z}{\left(1\mathrm{cos}z\right)}\right]dz=0\text{,}$$
where the singularity $`z=0`$ of the meromorphic function $`z\frac{1}{2}\frac{1}{1\mathrm{cos}z}`$ is bypassed by the parts $`G_{\epsilon _\kappa }`$: $`G_{\epsilon _\kappa }=\{zz=e^{i\theta _\kappa }\text{;}\text{ }\theta _\kappa \left\{\begin{array}{c}[\pi ,0]\text{;}\text{ }\kappa =1\hfill \\ [\pi ,0]\text{;}\text{ }\kappa =2\hfill \end{array}\right\}`$, of a circular path of integration: $`G_\epsilon =\{zz=e^{i\theta }\text{;}\text{ }\theta [\pi ,\pi ]\}`$.
Finally, from results: (3.13) and (3.14), the integral equation is obtained
$$\frac{1}{2\pi }v.t.\left[\underset{\pi }{\overset{𝜋}{}}\frac{\mathrm{cos}\left(k\theta \right)}{2\left(1\mathrm{cos}\theta \right)}d\theta \underset{\pi }{\overset{𝜋}{}}\frac{\mathrm{cos}\left[\left(k+1\right)\theta \right]}{2\left(1\mathrm{cos}\theta \right)}d\theta \right]=\frac{1}{2}\text{;}\text{ }kN\text{,}$$
(3.15)
and that is in agreement with result (3.4).
Further, since $`\mathrm{cos}\left(2\theta \right)=12\left(\mathrm{sin}\theta \right)^2`$, it follows that
$$v.t.\underset{\pi }{\overset{𝜋}{}}\frac{\mathrm{cos}\left(2\theta \right)}{2\left(1\mathrm{cos}\theta \right)}d\theta =v.t.\left[\underset{\pi }{\overset{𝜋}{}}\frac{d\theta }{2\left(1\mathrm{cos}\theta \right)}\underset{\pi }{\overset{𝜋}{}}\frac{\left(\mathrm{sin}\theta \right)^2}{1\mathrm{cos}\theta }d\theta \right]\text{,}$$
(3.16)
that is $`\frac{1}{\pi }v.t.\underset{\pi }{\overset{𝜋}{}}\frac{\mathrm{cos}\left(2\theta \right)}{2\left(1\mathrm{cos}\theta \right)}d\theta =2`$, in view of the results: (3.4) and (3.12).
Consequently, considering (3.15) it has been just proved by a method of mathematical induction that for each $`kN`$ it holds
$$A_k=\frac{1}{\pi }v.t.\underset{\pi }{\overset{𝜋}{}}\frac{\mathrm{cos}\left(kt\right)}{2\left(1\mathrm{cos}t\right)}dt=k\text{.}$$
(3.17)
As for an improper integral $`\underset{\pi }{\overset{𝜋}{}}\frac{\mathrm{sin}\left(kt\right)}{2\left(1\mathrm{cos}t\right)}dt`$, its total value ($`v.t.`$) reduces to the value of by-pass integrals $`\underset{G_{\epsilon _\kappa }}{}\frac{\mathrm{sin}\left(kz\right)}{2\left(1\mathrm{cos}z\right)}dz`$ in the limit as $`\epsilon 0^+`$. Taking into account that for each $`kN`$: $`\underset{z0}{lim}\frac{z\mathrm{sin}\left(kz\right)}{2\left(1\mathrm{cos}z\right)}=k`$, it follows that
$$B_k=\frac{1}{\pi }v.t.\underset{\pi }{\overset{𝜋}{}}\frac{\mathrm{sin}\left(kt\right)}{2\left(1\mathrm{cos}t\right)}dt=$$
(3.18)
$$=\frac{1}{\pi }\underset{\epsilon 0^+}{lim}\underset{G_{\epsilon _\kappa }}{}\frac{\mathrm{sin}\left(kz\right)}{2\left(1\mathrm{cos}z\right)}dz=\{\begin{array}{c}ik\text{;}\text{ }\kappa =1\hfill \\ ik\text{;}\text{ }\kappa =2\hfill \end{array}\text{}kN\text{}$$
Therefore, Fourier trigonometric series of the function $`t\frac{1}{2}\frac{1}{1\mathrm{cos}t}`$ in the segment $`[\pi ,\pi ]`$, according to results: (3.12) and (3.17) as well as (3.18), can be expressed by the following functional form
$$\frac{1}{2}\frac{1}{1\mathrm{cos}t}=\underset{k=1}{\overset{+\mathrm{}}{}}k\mathrm{cos}\left(kt\right)i\underset{k=1}{\overset{+\mathrm{}}{}}k\mathrm{sin}\left(kt\right)=\underset{k=1}{\overset{+\mathrm{}}{}}ke^{\pm ikt}\text{,}$$
(3.19)
that is, the equalities
$$\frac{1}{2}\frac{1}{1\mathrm{cos}t}=\underset{k=1}{\overset{+\mathrm{}}{}}k\mathrm{cos}\left(kt\right)\text{ and }\underset{k=1}{\overset{+\mathrm{}}{}}k\mathrm{sin}\left(kt\right)=0\text{,}$$
(3.20)
hold for each $`t(\pi ,\pi )`$ and $`t0`$, respectively.
In the extreme points of the segment $`[\pi ,\pi ]`$, from (2.15), it follows that
$$\underset{k=1}{\overset{+\mathrm{}}{}}k\mathrm{cos}\left(k\pi \right)=\underset{k=1}{\overset{+\mathrm{}}{}}k\left(1\right)^k=\frac{1}{4}\text{ and }\underset{k=1}{\overset{+\mathrm{}}{}}k\mathrm{sin}\left(k\pi \right)=0\text{.}\mathrm{}$$
(3.21)
By an expansion of the real valued functions of real variable $`t`$:
$$f\left(t\right)=\{\begin{array}{c}\frac{\mathrm{sin}t}{2\left(1\mathrm{cos}t\right)}\text{;}\text{ }\tau _0\left|t\right|\pi \hfill \\ 0\text{;}\text{ }\left|t\right|<\tau _0\hfill \end{array}\text{ and }g\left(t\right)=\{\begin{array}{c}b\text{;}\text{ }\tau _0t\pi \hfill \\ 0\text{;}\text{ }\left|t\right|<\tau _0\hfill \\ a\text{;}\text{ }\pi t\tau _0\hfill \end{array}\text{;}\text{ }\tau _0>0\text{,}$$
satisfying Dirichlet’s conditions in the segment $`[\pi ,\pi ]`$, into Fourier trigonometric series, it is obtained that
$$f\left(t\right)=\underset{k=1}{\overset{+\mathrm{}}{}}\left[\frac{1}{\pi }\underset{\tau _0}{\overset{𝜋}{}}\frac{\mathrm{sin}\tau }{1\mathrm{cos}\tau }\mathrm{sin}\left(k\tau \right)d\tau \right]\mathrm{sin}\left(kt\right)\text{}\left|t\right|(\tau _0,\pi )\text{,}$$
$$g\left(t\right)=\frac{1}{2\pi }(\underset{\pi }{\overset{\tau _0}{}}ad\tau +\underset{\tau _0}{\overset{𝜋}{}}bd\tau )+\underset{k=1}{\overset{+\mathrm{}}{}}\frac{1}{\pi }\{[\underset{\pi }{\overset{\tau _0}{}}a\mathrm{sin}\left(k\tau \right)d\tau +\underset{\tau _0}{\overset{𝜋}{}}b\mathrm{sin}\left(k\tau \right)d\tau ]\mathrm{sin}\left(kt\right)+$$
$$+[\underset{\pi }{\overset{\tau _0}{}}a\mathrm{cos}\left(k\tau \right)d\tau +\underset{\tau _0}{\overset{𝜋}{}}b\mathrm{cos}\left(k\tau \right)d\tau ]\mathrm{cos}\left(kt\right)\left\}\text{;}\text{ }\right|t|(\tau _0,\pi )\text{,}$$
that is<sup>3</sup><sup>3</sup>3By the well-known trigonometrical equalities: $`\mathrm{sin}\left[\left(k+1\right)t\right]=\mathrm{sin}\left(kt\right)\mathrm{cos}t+\mathrm{cos}\left(kt\right)\mathrm{sin}t`$ and $`\frac{\mathrm{sin}t\mathrm{sin}\left[\left(k+1\right)t\right]}{1\mathrm{cos}t}=\frac{\mathrm{sin}t\mathrm{cos}t}{1\mathrm{cos}t}\mathrm{sin}\left(kt\right)+\left(1+\mathrm{cos}t\right)\mathrm{cos}\left(kt\right)`$, as well as $`\mathrm{sin}\left(kt\right)\mathrm{cos}t=\frac{1}{2}\left\{\mathrm{sin}\left[\left(k1\right)t\right]+\mathrm{sin}\left[\left(k+1\right)t\right]\right\}`$ and $`\mathrm{cos}t\mathrm{cos}\left(kt\right)=\frac{1}{2}\left\{\mathrm{cos}\left[\left(k1\right)t\right]+\mathrm{cos}\left[\left(k+1\right)t\right]\right\}`$, the following recurrent formula for Fourier’s coefficients of the function $`f\left(t\right)`$ in the segment $`[\pi ,\pi ]`$ is obtained
$$\frac{1}{\pi }\underset{\tau _0}{\overset{𝜋}{}}\frac{\mathrm{sin}t\mathrm{sin}\left[\left(k+1\right)t\right]}{1\mathrm{cos}t}dt=\frac{1}{\pi }\underset{\tau _0}{\overset{𝜋}{}}\frac{\mathrm{sin}t\mathrm{sin}\left[\left(k1\right)t\right]}{1\mathrm{cos}t}dt\frac{2\mathrm{sin}\left(k\tau _0\right)}{k\pi }$$
$$\frac{\mathrm{sin}\left[\left(k+1\right)\tau _0\right]}{\left(k+1\right)\pi }\frac{\mathrm{sin}\left[\left(k1\right)\tau _0\right]}{\left(k1\right)\pi }\text{}kN$$
On the other hand, for $`k=1`$, that is, for $`k=2`$, it holds: $`\frac{1}{\pi }\underset{\tau _0}{\overset{𝜋}{}}\frac{\left(\mathrm{sin}t\right)^2}{1\mathrm{cos}t}dt=\frac{1}{\pi }\underset{\tau _0}{\overset{𝜋}{}}\left(1+\mathrm{cos}t\right)dt=1\frac{\tau _0}{\pi }\frac{\mathrm{sin}\tau _0}{\pi }`$, that is, $`\frac{1}{\pi }\underset{\tau _0}{\overset{𝜋}{}}\frac{\mathrm{sin}\left(2t\right)\mathrm{sin}t}{1\mathrm{cos}t}dt=\frac{2}{\pi }\underset{\tau _0}{\overset{𝜋}{}}\left(1+\mathrm{cos}t\right)\mathrm{cos}tdt=1\frac{\tau _0}{\pi }2\frac{\mathrm{sin}\tau _0}{\pi }`$ $`\frac{\mathrm{sin}\left(2\tau _0\right)}{2\pi }`$, respectively.
$$f\left(t\right)=\underset{k=1}{\overset{+\mathrm{}}{}}\left[1+\frac{\tau _0}{\pi }2\underset{\kappa =0}{\overset{k1}{}}\frac{\mathrm{sin}\left(\kappa \tau _0\right)}{\kappa \pi }\frac{\mathrm{sin}\left(k\tau _0\right)}{k\pi }\right]\mathrm{sin}\left(kt\right)\text{;}\text{ }\left|t\right|(\tau _0,\pi )\text{,}$$
(3.22)
$$g\left(t\right)=\frac{a+b}{2}\frac{a+b}{2\pi }\left[\frac{1}{2}+\underset{k=1}{\overset{+\mathrm{}}{}}\frac{\mathrm{sin}\left(k\tau _0\right)}{k\tau _0}\mathrm{cos}\left(kt\right)\right]\tau _0+$$
(3.23)
$$+\frac{ba}{\pi }\underset{k=1}{\overset{+\mathrm{}}{}}\left[\mathrm{cos}\left(k\tau _0\right)\left(1\right)^k\right]\frac{\mathrm{sin}\left(kt\right)}{k}\text{;}\text{ }\left|t\right|(\tau _0,\pi ).$$
From the functional relation (3.23) it follows for $`a=b`$, $`\left|t\right|(\tau _0,\pi )`$ and $`\tau _0>0`$ that
$$\frac{1}{2}+\underset{k=1}{\overset{+\mathrm{}}{}}\frac{\mathrm{sin}\left(k\tau _0\right)}{k\tau _0}\mathrm{cos}\left(kt\right)=0\text{.}$$
(3.24)
Thus, for $`t(\tau _0,\pi )`$ and $`\tau _0>0`$ it holds
$$\frac{\pi }{2}=\underset{k=1}{\overset{+\mathrm{}}{}}\left[\mathrm{cos}\left(k\tau _0\right)\left(1\right)^k\right]\frac{\mathrm{sin}\left(kt\right)}{k}\text{,}$$
(3.25)
that is
$$\underset{k=1}{\overset{+\mathrm{}}{}}\mathrm{cos}\left(k\tau _0\right)\frac{\mathrm{sin}\left(kt\right)}{k}=\frac{\pi }{2}\frac{t}{2}\text{;}\text{ }t(\tau _0,\pi )\text{,}$$
(3.26)
since for $`t(\pi ,\pi )`$, ,
$$\underset{k=1}{\overset{+\mathrm{}}{}}\left(1\right)^k\frac{\mathrm{sin}\left(kt\right)}{k}=\frac{t}{2}\text{.}$$
(3.27)
On the other hand, taking into account the fact that
$$\underset{k+\mathrm{}}{lim}\left[1+\frac{\tau _0}{\pi }2\underset{\kappa =0}{\overset{k1}{}}\frac{\mathrm{sin}\left(\kappa \tau _0\right)}{\kappa \pi }\frac{\mathrm{sin}\left(k\tau _0\right)}{k\pi }\right]=0\text{,}$$
see the last Comment in the preceding Section of this paper, finally it follows for $`\tau _0(0,\pi )`$ that
$$\underset{k=0}{\overset{+\mathrm{}}{}}\frac{\mathrm{sin}\left(k\tau _0\right)}{k}=\frac{\pi }{2}+\frac{\tau _0}{2}\text{.}$$
(3.28)
Since a real parameter $`\tau _0`$ takes any value from the interval $`(0,\pi )`$, even if that is finitely small value, it would be reasonable to ask: Whether the functional expressions: (3.22) and (3.24) as well as (3.25) and (3.26), hold in the limit as $`\tau _00^+`$? In other words: Are the limiting values of sums, in these emphasized cases, equal to sums of limiting values of the functional expressions, as $`\tau _00^+`$, respectively? On the basis of previously derived results in the Example 1 and of the above obtained result (3.28) as well as of the well-known result of the series theory: $`\frac{\pi }{4}=\underset{k=1}{\overset{+\mathrm{}}{}}\frac{\mathrm{sin}\left[\left(2k1\right)t\right]}{2k1}`$; for $`t(0,\pi )`$, , an answer to the former questions is yes. However, the problem of generalization of a preceding conclusion stays open and can be a subject of a separate analysis.
Similarly, since $`\frac{d}{dt}\left[\frac{\mathrm{sin}t}{2\left(1\mathrm{cos}t\right)}\right]=\frac{1}{2\left(1\mathrm{cos}t\right)}`$ and $`\frac{d}{dt}\left\{\frac{1}{2}\mathrm{ln}\left[2\left(1\mathrm{cos}t\right)\right]\right\}=\frac{\mathrm{sin}t}{2\left(1\mathrm{cos}t\right)}`$ for $`\left|t\right|(0,\pi )`$, then closely related to results: (3.7) and (3.20), of the paper, as well as to the well-known result of the series theory: $`\frac{1}{2}\mathrm{ln}\left[2\left(1\mathrm{cos}t\right)\right]=\underset{k=1}{\overset{+\mathrm{}}{}}\frac{\mathrm{cos}\left(kt\right)}{k}`$, for $`\left|t\right|(0,\pi )`$, , is the following question: In which general cases the derivative of a sum of infinite functional series is equal to the sum of the derivative of any series member, separately?
This question also stays open for a separate analysis.$`\mathrm{}`$
## 4. Conclusion
Taking into consideration the fact that obtained results are theoretical news, one can say that the certain possibilities for expansion of some mathematical analysis knowledge connecting to the problems to which a proper attention has been paid in this paper are opening up. Thus, from viewpoint of results, derived in the Subsubsection 2.2.4 of the paper for instance, and having in mind the fact that causality related to the area of Fourier trigonometric series of real valued functions is the theory of partial differential equations, it is obvious which possibilities are opening up in this area of mathematics.
On the other hand, disregarding the fact that the results of the paper are, in a certain sense, the theoretical news, some of them have been predictable. So, the alternative numerical series: $`\underset{k=0}{\overset{+\mathrm{}}{}}\left(1\right)^k`$, has the defined sum, more exactly it is summable and its sum is equal to $`\frac{1}{2}`$, just as it has been assumed yet by Euler and Dalamber. Making use of this assumption they obtained absolutely exact results. It is nothing other to be left than to prove validity of this assumption. As for the results: (3.20) and (3.21), from the Example 2, one can say that they are theoretical news and causality related to the result (3.8). Namely, since $`\underset{k=1}{\overset{+\mathrm{}}{}}k\mathrm{sin}\left(kt\right)=0`$ for $`t=\frac{\pi }{2}`$, that is $`\underset{k=0}{\overset{+\mathrm{}}{}}\left(2k+1\right)\left(1\right)^k=0`$, it follows that $`\underset{k=0}{\overset{+\mathrm{}}{}}2k\left(1\right)^k=\underset{k=0}{\overset{+\mathrm{}}{}}\left(1\right)^k=\frac{1}{2}`$.
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# Reionization by Hard Photons: I. X-rays from the First Star Clusters
## 1 Introduction
While the theoretical literature on the epoch of reionization is large and increasing rapidly, our empirical knowledge of this period in the history of the universe is scant and may be succintly summarized: (i) The universe is likely to have been reionized in the period $`5.8<z_r<35`$, where the lower bound arises from the lack of Gunn-Peterson absorption in the spectra of high redshift quasars (Fan et al 2000), and the upper bound comes from the observation of small scale power in Cosmic Microwave Background anistropies (Griffiths et al 1999). In CDM cosmologies, non-linear objects above the cosmological Jeans mass $`10^{45}M_{}`$ first collapse during this period. (ii) The presence of metals in the Ly$`\alpha `$ forest implies that significant star formation took place at high redshift (Songailia & Cowie 1996, Songailia 1997). (iii) COBE constraints on the Compton y-distortion of the CMB (Wright et al, 1994, Fixsen et al 1996) implies that the IGM was not heated to high temperatures. This means it is unlikely that the hydrogen and helium in the IGM were collisionally reionized.
Thus, the current state of observations is consistent with scenarios in which the universe was reionized by an early generation of stars or quasars. While we do not know whether stars or quasars were the dominant source of ionizing photons, the observational and theoretical case for quasars is somewhat more uncertain. Extrapolation of empirical quasar luminosity functions to high redshift do not yield enough ionizing photons to maintain the observed lack of Gunn-Peterson absorption at $`z5`$ (Madau, Haardt & Rees 1999). The requisite steepening of the faint end slope at high redshift necessary to boost ionizing photon production is constrained by the lack of red, point-like sources in the Hubble Deep Field (Haiman, Madau & Loeb 1999); the authors find that AGN formation must have been suppressed in halos with $`v_c<5075\mathrm{km}\mathrm{s}^1`$. We do not have a sufficiently firm understanding of the formation and fueling of supermassive black holes to assert on theoretical grounds that AGNs must have been present at high redshift. On the other hand, a minimal level of high-redshift star formation is guaranteed by the observed metal pollution of the IGM.
Theoretical scenarios in which stars or quasars figure predominantly have been calculated in detail. Our ignorance of the efficiency of gas fragmentation, and star/black hole formation as a function of halo mass, make the prediction of observable differences between these two scenarios very uncertain. Indeed, if one normalizes assumed emissivities to a fixed reionization epoch, differences between the two scenarios boil down to: (i) stars result in supernovae, which inject dust, metals, and entropy into the host galaxy and surrounding IGM, which affects subsequent chemistry and cooling, (ii) quasars have a significantly harder spectrum than stars. In particular, they produce X-rays.
In this paper, I emphasize a hitherto neglected fact: high redshift supernova also produce X-rays, both by thermal emission from the hot supernova remnant, and inverse Compton scattering of soft photons by relativistic electrons accelerated by the supernova. Considerable X-ray emission is already observed in starburst galaxies at low redshift (e.g., Rephaeli et al 1995), and the efficiency of most proposed X-ray production mechanisms should increase with redshift (e.g., explosions take place in a denser medium at high redshift, hardening expected thermal emission; inverse Compton scattering becomes more efficient since the CMB provides a ready supply of soft photons $`U_{CMB}(1+z)^4`$). Thus, the SED of high-redshift star forming regions is considerably harder than has been previously assumed. This blurs the distinction between stellar/quasar reionization scenarios, and has a number of important physical consequences:
* Escape fraction The escape fraction of UV ionizing photons in the local universe is small, $`36\%`$ (Leitherer et al 1995, Bland-Hawthorne & Maloney 1999, Dove, Shull & Ferrara 2000) and is expected to decrease with redshift (Wood & Loeb 1999, Ricotti & Shull 1999). On the other hand, X-rays can escape freely from the host galaxy. Thus, processing by the host ISM may imply that the universe was reionized by a significantly harder spectrum than previously assumed.
* Reionization topology Photons from stellar spectra have a short mean free path and thus a sharply defined ionization front. This fact gives rise to the conventional picture of expanding HII bubbles embedded in the neutral IGM. The spectra of quasars is significantly harder and exert an influence over a larger distance (which is why it is much more difficult to perform numerical simulations of reionization by quasars; see Gnedin 1999). Nonetheless, for quasars, $`\nu L_\nu \nu ^{0.8}`$ (Zheng et al 1997, although note that the observations were only in the radio-quiet AGN subsample at energies up to 2.6 Ry) and most of the energy for ionization lies just above the Lyman edge. Thus, there is still a sharply defined HII region and a thin ionization front where the ionization fraction drops sharply. By contrast, for the inverse Compton case $`\nu L_\nu \mathrm{const}`$, there is equal power per logarithmic interval, and thus there is no preferred energy scale. In particular, there is no preferred scale for the mean free path of ionizing photons. When the universe is largely neutral, it is optically thick even to hard photons and all photons with energies $`E<E_{thick}=1.5\left(\frac{1+z}{10}\right)^{0.5}x_{HI}^{1/3}\mathrm{keV}`$ (where $`x_{HI}`$ is the mean neutral fraction) are absorbed across a Hubble volume. While the (more numerous) soft photons can only travel a short distance before ionizing neutral HI and HeI, the (less numerous but more energetic) hard photons will be able to travel further and ionize an equivalent number of photons by secondary ionizations. Thus, even if sources are distributed very inhomogeneously, reionization will be a fairly homogeneous event, with a largely uniform ionizing background and fluctuations in ionization fraction determined mainly by gas clumping. Instead of an two-phase medium in whose HII filling fraction increases with time, the early IGM may have been a single phase medium whose ionization fraction increases with time.
* Increased reheating A hard spectrum can reheat the IGM to considerably higher temperatures than soft stellar spectra, both through photoionization heating and Compton heating. A soft spectrum loses thermal contact with the IGM once HI and HeI are completely ionized (at the mean IGM density, HI has a recombination time longer than the Hubble time for $`z<10`$), and the gas cools adiabatically due to the expansion of the universe (Hui & Gnedin 1997). By contrast, a hard spectrum can continually transfer large amounts of energy from the radiation field to the IGM by ionizing HeII, which recombines rapidly (Miralda-Escude & Rees 1994). This feedback mechanism is important in increasing the Jeans mass, a proposed mechanism for preventing excessive cooling and star formation at high redshift (e.g., Prunet & Blanchard 1999). The higher IGM temperatures may also explain why observed Ly$`\alpha `$ forest line widths are commonly in excess of that predicted by numerical simulations (Theuns et al 1999, Ricotti et al 2000).
* Early universe chemistry $`\mathrm{H}_2`$ is an extremely important coolant in the metal-free early universe. While the neutral IGM is optically thick to UV ionizing photons, it is optically thin to photons longward of the Lyman limit (except at wavelengths corresponding to higher order hydrogen Lyman resonance lines, as well as $`H_2`$ resonance lines). In particular, photons in the 11.2-13.6 eV range quickly establish a soft UV background which photodissociates $`\mathrm{H}_2`$ via the Solomon process, shutting down subsequent star formation (Haiman, Rees & Loeb 1997, Cicardi, Ferrara & Abel 1998), unless the $`H_2`$ opacity is sufficient to reduce the photo-dissociation rate (Ricotti, Gnedin & Shull 2000). If X-rays are present in the early universe, they can counter this $`H_2`$ destruction. The IGM is also optically thin to X-rays, which can penetrate dense clouds of gas and promote gas phase $`\mathrm{H}_2`$ formation $`\mathrm{H}+\mathrm{e}^{}\mathrm{H}^{}+\gamma `$ and $`\mathrm{H}^{}+\mathrm{H}\mathrm{H}_2+\mathrm{e}^{}`$ by increasing the abundance of free electrons. Haiman, Abel & Rees (1999) show that if quasars were the dominant ionizing sources in the early universe, gas cooling and thus star formation can continue unabated. In Paper II, I show that in fact even if only stars were present, a self-consistent treatment of the stellar SED incorporating X-rays produced by supernovae favours $`\mathrm{H}_2`$ formation over destruction in dense regions.
In this paper, I study the emission mechanisms and observational signatures of X-ray bright star clusters at high redshift. In Paper II (Oh 2000a), I address the changes in reionization topology, reheating and early universe chemistry mentioned above due to these X-rays.
In all numerical estimates, I assume a background cosmology given by the ’concordance’ values of Ostriker & Steinhardt (1995): $`(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda },\mathrm{\Omega }_b,h,\sigma _{8h^1},n)=(0.35,0.65,0.04,0.65,0.87,0.96)`$. This corresponds to $`\mathrm{\Omega }_bh^2=0.017`$, compared with $`\mathrm{\Omega }_bh^2=0.020\pm 0.002(95\%\mathrm{c}.\mathrm{l}.)`$ (Burles, Nollett & Turner 2000), and $`\mathrm{\Omega }_bh^2=0.0205\pm 0.0018`$ (O’Meara et al 2000) from Big Bang Nucleosynthesis, and the significantly higher values $`0.022<\mathrm{\Omega }_bh^2<0.040(95\%\mathrm{c}.\mathrm{l}.)`$ (Tegmark & Zaldarriaga 2000) preferred by recent CMB anisotropy data, such as Boomerang and Maxima.
## 2 Emission mechanisms
### 2.1 Star formation at $`z>3`$
What fraction of present day stars formed at high redshift? Estimates of the comoving star formation rate as a function of redshift (Madau et al, 1996) should be regarded as lower bounds, particularly at high redshift, due to the unknown effects of dust extinction, and star formation in faint systems below the survey detection threshold. Indeed, Lyman break survey results for $`3.6<z<4.5`$ (Steidel et al 1999) suggest that after correction for dust extinction, the comoving star formation rate for $`z>1`$ is constant, rather than falling sharply as previously believed. Furthermore, in recent years compelling evidence has emerged that the majority of stars in ellipticals and bulges formed at high redshift, $`z>3`$. This comes from the tightness of correlations between various global properties of ellipticals which indicate a very small age dispersion and thus a high redshift of formation, unless their formation was synchronized to an implausible degree. The evidence includes the tightness of the fundamental plane and color magnitude relations for ellipticals, and the modest shift in zero-points for these relations with redshift (Renzini 1998 and references therein). Since spheroids contain $`30\%`$ of all stars in the local universe (King & Ellis 1985, Schechter & Dressler 1987), this would imply that $`30\%`$ of all stars have formed at $`z>3`$. Since $`20\%`$ of baryons have been processed into stars by the present day (Fukugita, Hogan & Peebles 1998), this implies that $`0.2\times 0.36\%`$ of baryons have been processed into stars by $`z3`$. Assuming widespread and uniform enrichment and 1 $`\mathrm{M}_{}`$ of metals per 100 $`\mathrm{M}_{}`$ of stars formed, this translates into an IGM metallicity of $`6\times 10^43\times 10^2Z_{}`$. At $`z3`$, the metallicity of damped Ly$`\alpha `$ systems appears to be $`0.05Z_{}`$ (Pettini et al 1997), in reasonable agreement. The observed metallicity of the Ly$`\alpha `$ forest at $`z3`$, to which most models of reionization have been normalized, is between $`10^2Z_{}`$ and $`10^3Z_{}`$ (Songailia & Cowie 1996, Songailia 1997), which would imply that only $`0.22\%`$ of present day stars formed at $`z>3`$. However, its metallicity may be more representative of low density regions, rather than the mean cosmological metallicity (Cen & Ostriker 1999). Note that normalization of high redshift star formation to Lyman $`\alpha `$ forest metallicities assume efficient metal ejection (which underestimates star formation if a significant fraction of metals are retained) and mixing (which overestimates star formation if Ly$`\alpha `$ lines are preferentially observed in overdense regions which are sites of star formation). In this context, it is worth mentioning claims that Ly$`\alpha `$ lines with $`10^{13.5}\mathrm{cm}^2<\mathrm{N}_{\mathrm{HI}}<10^{14.5}\mathrm{cm}^2`$ reveal lower metallicities by a factor of 10 than clouds with $`\mathrm{N}_{\mathrm{HI}}>10^{14.5}\mathrm{cm}^2`$ (Lu et al 1999). Ellison et al (2000) find no break in the power law column density distribution for C IV down to log N(C IV) =11.7, and Schaye et al (2000) detect O VI down to $`\tau _{HI}10^1`$ in underdense gas, so it appears that metal pollution was fairly widespread. I regard $`\mathrm{Z}10^32.5\times 10^2Z_{}`$ at $`z3`$ as a fairly firm bracket on the range of possibilities. For inverse Compton radiation, this corresponds to an energy release per IGM baryon of $`ϵ^{SN}10\left(\frac{Z}{10^2Z_{}}\right)\left(\frac{ϵ}{0.1}\right)`$eV (where $`ϵ`$ is the efficiency of conversion of supernova energy to X-rays), which is comparable to the energy release in stellar UV radiation for the low escape fractions expected, $`ϵ^{stellar}10\left(\frac{Z}{10^2Z_{}}\right)\left(\frac{f_{esc}}{0.01}\right)`$eV, where $`f_{esc}`$ is the escape fraction of ionizing photons from the source.
### 2.2 X-ray emission in local starbursts
Most models of reionization use population synthesis codes to estimate the spectral energy distribution of starbursts. However, there are many processes associated with star formation that generate UV and X-rays, beside stellar radiation: massive X-ray binaries, thermal emission from supernova remnants and hot gas in galactic halos and winds, inverse Compton scattering of soft photons by relativistic electrons produced in supernovae. Indeed, X-ray emission appears to be ubiquitous among starbursts (e.g., Rephaeli, Gruber, & Persic 1995), and starburst galaxies may account for a significant portion of the XRB (Bookbinder et al 1980, Rephaeli et al 1991, Moran, Lehnert & Helfand 1999). The X-ray emission from these processes, which hardens the spectrum of starbursts and changes both the topology and chemistry of reionization, has to date been neglected in studies of the $`z>5`$ universe.
The X-ray luminosity of starbursts correlates well with other star formation indicators; for example, David et al (1992) find a roughly linear relation between $`L_{FIR}`$ and $`L_X`$. As a very rough empirical calibration, the starburst galaxies M82 & NGC 3256 observed with ROSAT and ASCA (Moran & Lehnert 1997, Moran, Lehnert & Helfand, 1999) follow the relations (after correction for absorption): $`L_{X,0.210\mathrm{keV}}=8\times 10^4L_{IR}`$, and $`L_{X,5\mathrm{keV}}=1.2\times 10^4L_{R,5\mathrm{GHz}}`$ (note that 5 keV flux density is relatively unaffected by photoelectric absorption or soft thermal emission). The starburst model of Leitherer and Heckman (1995) yields $`L_{bol}L_{\mathrm{FIR}}1.5\times 10^{10}(\mathrm{SFR}/1\mathrm{M}_{}\mathrm{yr}^1)L_{}`$; as a cross-check, the empirical relation for radio emission is (Condon 1992) $`L_R=1.4\times 10^{28}(\nu /\mathrm{GHz})^\alpha (\mathrm{SFR}/\mathrm{M}_{}\mathrm{yr}^1)\mathrm{erg}\mathrm{s}^1\mathrm{Hz}^1`$, where $`\alpha 0.8`$. Together I obtain:
$$L_X=5\times 10^{40}\left(\frac{SFR}{1\mathrm{M}_{}\mathrm{yr}^1}\right)\mathrm{erg}\mathrm{s}^1$$
(1)
One should not regard this as more than a rough order of magnitude estimate; a large scatter is expected in this relation. By way of comparison, Rephaeli et al (1995) obtain from the mean of 51 starbursts observed with Einstein and HEAO, $`L_{X,230\mathrm{keV}}8\times 10^3L_{IR}`$, an order of magnitude greater; and David et al (1992) obtain from a sample of 71 normal and starburst galaxies a ratio lower by about an order of magnitude, largely due to the inclusion of normal galaxies (this is consistent with an inverse Compton origin for X-rays, since normal galaxies have much lower radiation field energy densities and would not be expected to show significant inverse Compton emission).
The typical observed X-ray spectrum is a power-law, $`L_\nu \nu ^{0.8}`$, consistent with a non-thermal origin. An obscured AGN is not likely to be the source of these X-rays, as several observations suggest that the X-ray emission is powered primarily by massive stars. In NGC 3256, observations by ISO fail to detect high excitation emission lines (Rigopoulou et al 1996). In M82, the optical spectrum is HII-like (Kennicutt 1992), discrete nuclear radio sources are spatially resolved (Muxlow et al 1994), its nuclear X-ray emission is extended (Bregman et al 1995), and the expected broad H$`\alpha `$ emission is not detected (Moran & Lehnert 1997); the HEX continuum and Fe-K line emission of NGC 253 as observed by BeppoSAX is extended (Cappi et al 1999). Moran & Lehnart (1997) and Moran, Lehnart, & Helfand (1999) have modelled the X-ray emission of M82 and NGC 3256, and find inverse-Compton emission to be the most likely mechanism, rather than an obscured AGN or massive X-ray binaries. The ratio between the observed radio and X-ray fluxes (note that $`L_X/L_{syn}U_{IR}/U_B`$, where $`U_{IR}`$ is the energy density of the infra-red radiation field, and $`U_B`$ is the energy density of the magnetic field) is consistent with an inverse-Compton origin for the X-rays. Furthermore, the X-ray and radio emission have the same spectral slope, as is expected if both types of emission are non-thermal, arising from the same population of electrons. The X-ray luminosity is energetically consistent with a star formation origin. Assuming a Salpeter IMF and that each supernova explosion yields $`10^{51}`$ erg in kinetic energy yields an energy injection rate into the ISM:
$$\dot{E}_{SN}3\times 10^{40}\left(\frac{\mathrm{SFR}}{1\mathrm{M}_{}\mathrm{yr}^1}\right)\left(\frac{ϵ}{0.1}\right)\mathrm{erg}\mathrm{s}^1$$
(2)
where $`ϵ`$ is the fraction of energy injected into relativistic electrons; consistency with the observed value, (1), implies that $`ϵ10\%`$. Hereafter, I shall use the empirical equation (1) as a fiducial conversion between X-ray luminosity and star formation rate.
Is an acceleration efficiency of $`ϵ10\%`$ reasonable? The acceleration mechanism for relativistic electrons is poorly understood. While first-order Fermi acceleration in shocks is widely accepted as the acceleration mechanism for cosmic rays (Blandford & Eichler 1987, Jones & Ellison 1991), the electron acceleration is thought to be more problematic, due to the smaller electron gyroradius (which leads to greater difficulties in bouncing an electron back and forth across a shock of finite thickness), and the difficulty of initially boosting the electron to relativistic speeds, where Fermi acceleration can operate (Levinson 1994). From measurements of cosmic rays energy density it is inferred that $`10\%`$ of the supernova kinetic energy, or $`10^{50}`$erg per explosion, is liberated as cosmic rays (Volk, Klein, & Wielebinski 1989), but the division between electrons and protons at the source is not known. Since the measured ratio of cosmic ray protons to electrons is $`75`$ (e.g., Gaisser 1990), it might well be that relativistic electrons only constitute $`ϵ10^3`$ of the supernova energy budget. Thus, the reader should be cautioned that $`ϵ0.1`$, which corresponds roughly equal energy division between protons and electrons (e/p=1), may be an overly optimistic estimate of the energy injection into relativistic electrons. Theoretical models of shock acceleration, in which e/p is a free parameter, often set e/p$`15\%`$ for consistency with cosmic-ray experiments (Ellison & Reynolds 1991, Ellison et al 2000). However, this is somewhat model-dependent: in models where electrons are injected directly from the thermal pool, $`5\%`$ of the energy in the shock must go to non-thermal electrons in order to match gamma-ray observations (Bykov et al 2000). Furthermore, note that the observed cosmic-ray e/p ratio could equally well be the result of different transport processes and energy loss mechanisms for electrons and protons. In particular, cosmic ray electrons are subject to loss processes which operate on much longer timescales for cosmic ray protons (inverse Compton, synchrotron losses, etc); the cosmic ray flux at earth for electrons could arise from a much smaller effective volume than that for protons. At the source, the energy division between protons and electrons could range between 1 and 100. Perhaps the most reliable means of inferring the proton/electron energy division is by direct observations of supernova remnants. In modelling the observed production of gamma-rays in the supernova remnants IC 443 and $`\gamma `$ Cygni observed by the EGRET instrument on the Compton Gamma Ray Observatory, Gaisser, Protheroe & Stanev (1998) find that a proton to electron ratio of 3–5 gives the best fit to the observed spectra, implying $`ϵ0.020.03`$. Similarly, in modelling the $`\gamma `$-ray flux from 2EG J1857+0118 associated with supernova remnant W44, de Jager & Mastichiadis (1997) find $`ϵ0.09`$. They speculate that electron injection by the pulsar may be responsible for the increased electron energy content. Given the large uncertainties, henceforth I shall simply use the empirical relation (1).
Thus, barring non-standard IMFs, type II detonation energies or alternate sources of relativistic electrons, the empirical relation (1) implies an acceleration efficiency $`ϵ0.1`$ which is plausible but certainly lies at the upper limit of theoretical expectations. Another possibility is that the X-ray emission cannot be wholly attributed to inverse Compton emission alone (this assumption rests on the arguments of Moran & Lehnart (1997) and Moran, Lehnart, & Helfand (1999) with regards to the slopes and relative intensities of the observed non-thermal radio and X-ray emission). The X-ray emission may instead be due to X-ray binaries, thermal emission from supernova remnants or starburst driven superwinds (e.g., see Natarajan & Almaini 2000). It should be noted that only soft X-rays are relevant for reionization, since the universe is optically thin to photons with energies $`E>E_{thick}=1.5\left(\frac{1+z}{10}\right)^{0.5}x_{HI}^{1/3}\mathrm{keV}`$ (where $`x_{HI}`$ is the mean neutral fraction of the IGM). Since equation (1) is calibrated with the soft bands observed by ROSAT, this implies that even if $`ϵ0.1`$ and the observed X-rays are not predominantly due to inverse Compton emission, the importance of X-rays for reionization (in particular, for changing the topology, for increased reheating, and increased $`H_2`$ production) may still hold. However, observational predictions which focus specifically on the inverse Compton mechanism (e.g., the gamma-ray background (section (3.2)), and detecting synchrotron emission with the SKA (section (3.4)) will no longer be valid. Since the acceleration efficiency is the most uncertain parameter in this paper, wherever relevant I insert the scaling factor $`\left(\frac{ϵ}{0.1}\right)`$ into numerical estimates.
How does the X-ray luminosity compare with stellar UV ionizing radiation? Assuming a Salpeter IMF with solar metallicity, the Bruzual & Charlot (1999) population synthesis code yields an energy output of $`L_{ion}=3.2\times 10^{42}(\mathrm{SFR}/1\mathrm{M}_{}\mathrm{yr}^1)\mathrm{erg}\mathrm{s}^1`$ in ionizing photons, which translates into $`\dot{\mathrm{N}}_{\mathrm{ion}}=10^{53}(\mathrm{SFR}/1\mathrm{M}_{}\mathrm{yr}^1)\mathrm{photons}\mathrm{s}^1`$. However, note that most of these ionizing photons are absorbed locally with the ISM of the star cluster; the escape fraction of ionizing photons into the IGM is expected to be small. Leitherer et al (1995) have observed four starburst galaxies with the Hopkins Ultraviolet Telescope (HUT). Their analyis suggests an escape fraction of only $`3\%`$, based on a comparison between the observed Lyman continuum flux and theoretical spectral energy distributions. For our own Galaxy, Dove, Shull & Ferrara (2000) find an escape fraction for ionizing photons of $`6\%`$ and $`3\%`$ (for coeval and Gaussian star formation histories respectively) from OB associations in the Milky Way disk. Bland-Hawthorn & Maloney (1999) find an escape fraction of $`6\%`$ is necessary for consistency with the observed H$`\alpha `$ emission from the Magellanic stream and high velocity clouds. On the other hand, a recent composite spectrum of 29 Lyman break galaxies (LBGs) with redshifts $`z=3.40\pm 0.09`$ shows significant detection of Lyman continuum flux (Steidel, Pettini & Adelberger 2000); for typical stellar synthesis models, the observed flux ratio L(1500)/L(900)=$`4.6\pm 1.0`$ implies little or no photoelectric absorption. The fraction of 900 $`\mathrm{\AA }`$ photons which escape, $`f_{esc}1520\%`$, is modulated almost entirely by dust absorption. Nonetheless, the authors themselves stress this result should be treated as preliminary; the result could be due to a large number of uncertainties or selection effects, among them the fact that these galaxies were selected from the bluest quartile of LBGs. On the theoretical side, radiative transfer calculations by Woods & Loeb (1999) find that the escape fraction at $`z10`$ is $`<1\%`$ for stars; calculations by Ricotti & Shull (1999) find that the escape fraction decreases strongly with increasing redshift and halo mass; for a $`10^9M_{}`$ halo at $`z=9`$, the escape fraction is $`10^3`$ (note that in the Ricotti & Shull (1999) models, the escape fraction rises towards low masses, and can be considerable for the halos with $`M<10^7M_{}`$. Since such halos have $`T_{vir}<10^4`$K, their contribution to reionization depends on whether $`H_2`$ formation and cooling can take place despite photodissociative processes (Haiman, Rees & Loeb 1997, Cicardi, Ferrara & Abel 1998, Ricotti, Gnedin & Shull 2000)). For low escape fractions, the energy release in UV photons is roughly comparable to that in inverse-Compton X-rays :
$$L_{UV}=3\times 10^{40}\left(\frac{f_{esc}}{0.01}\right)\left(\frac{\mathrm{SFR}}{1\mathrm{M}_{}\mathrm{yr}^1}\right)\mathrm{erg}\mathrm{s}^1$$
(3)
Note that to first order the ratio of stellar UV to inverse Compton X-rays is not sensitive to uncertainties in the IMF, as the same massive stars with $`M>20M_{}`$ that dominate the Lyman continuum of a stellar population also explode as supernovae (however, see section (2.4) for some caveats).
Will X-ray emission still be efficient at high redshift? The following changes are expected to take place at high redshift: (i) The ISM is initially free of dust and metals, although rapid enrichment could occur on fairly short timescales ($`t10^610^7`$yr). (ii) The average ISM density is significantly higher, $`n_{halo}n_018\pi ^2(1+z)^3=0.02(\frac{1+z}{10})^3\mathrm{cm}^3`$ in a halo and $`n_{disc}\lambda ^3n_{halo}160(\frac{1+z}{10})^3\mathrm{cm}^3(\lambda /0.05)^3`$ in a disk (where $`\lambda `$ is the spin parameter). A supernova remnant at high $`z`$ expands into a denser ISM: since it spends a shorter time in the Taylor-Sedov phase, most of its energy is radiated at a smaller radius, where the effective temperature is higher. This implies a harder spectrum for thermal emission. In addition, the density and temperature of gas in star forming regions is determined by the properties of $`\mathrm{H}_2`$ cooling (which saturates at $`n10^4\mathrm{cm}^3`$ and $`T300`$K), rather than metal cooling as in the local universe. (iii) Potential wells are significantly shallower, so pressurised regions (hot gas, strong magnetic fields) cannot be efficiently confined. (iv) The CMB energy density $`U_{CMB}(1+z)^4`$, so inverse Compton radiation becomes particularly efficient at high redshift. Because of (iii) and (iv), relativistic electrons cool predominantly by inverse Compton scattering rather than synchrotron emission. Since inverse Compton emission is likely to be the most promising mechanism for X-ray emission, I shall consider it at length.
### 2.3 Inverse Compton emission
The energy loss rate for a relativistic electron with Lorentz factor $`\gamma `$ is given by:
$$\dot{\mathrm{E}_{IC}}=\frac{4}{3}\sigma _Tc\gamma ^2U_{rad}=1.12\times 10^{16}\left(\frac{1+z}{10}\right)^4\left(\frac{\gamma }{10^3}\right)^2\mathrm{erg}\mathrm{s}^1$$
(4)
for $`U_{rad}=U_{\mathrm{CMB}}`$. The loss rate by synchrotron radiation is given by substiting $`U_B`$ for $`U_{rad}`$, $`\dot{\mathrm{E}_{synch}}=\frac{4}{3}\sigma _Tc\gamma ^2U_B`$. In the local universe, galaxies with relatively quiescent star formation emit most of their electron energy in synchrotron radiation. Starbursts in the local universe can radiate efficiently in IC, as the energy density in the local radiation field is sufficiently high (typically, $`U_r10^8\mathrm{erg}\mathrm{cm}^3`$ as opposed to $`U_r10^{12}\mathrm{erg}\mathrm{cm}^3`$ in our Galaxy). Seed photons are provided by IR emission from dust grains. By contrast, at high redshift, all star forming regions will emit in inverse Compton radiation, as the CMB provides a universal soft photon bath of high energy density, $`U_{CMB}=4\times 10^9(\frac{1+z}{10})^4\mathrm{erg}\mathrm{cm}^3`$. An electron with Lorentz factor $`\gamma `$ will boost a CMB photon of frequency $`\nu _o`$ to a frequency $`\nu =\gamma ^2\nu _o`$. Thus, the rest frame frequency of a CMB photon (at the peak of the blackbody spectrum) which undergoes inverse Compton scattering is:
$$\mathrm{E}_{IC}=600\left(\frac{\gamma }{300}\right)^2\left(\frac{1+z}{10}\right)\mathrm{eV}$$
(5)
Note that the observed frequency is independent of source redshift, since the higher initial frequency and redshifting effects cancel out.
Below, I examine in detail the mechanisms by which relativistic electrons lose energy, to see if indeed inverse Compton radiation will predominate.
#### 2.3.1 Energy loss processes for relativistic electrons
Once relativistic electrons are produced, they can cool via a variety of mechanisms. Let us examine them in turn (for more details see Pacholczyk 1970, Daly 1992).
Synchrotron radiation The relative emission rate in inverse Compton and synchrotron emission is given simply by the relative energy densities in the radiation and magnetic fields, $`\dot{E}_{IC}/\dot{E}_{syn}=U_{rad}/U_m`$. This yields:
$$\frac{\dot{E}_{IC}}{\dot{E}_{syn}}=1.1\times 10^3\left(\frac{B}{10\mu G}\right)^2\left(\frac{1+z}{10}\right)^4$$
(6)
Note that in energy loss terms the CMB may be characterized as having an effective magnetic field strength $`B_{CMB,eff}=3.24\times 10^2\left(\frac{1+z}{10}\right)^2\mu \mathrm{G}`$. Could magnetic fields in proto-galaxies possibly reach these high values? A reasonal assumption is that $`P_BP_{rel}<P_{gas}`$, where $`P_B`$ is the magnetic field pressure, $`P_{rel}`$ is the pressure in relativistic particles, and $`P_{gas}`$ is the thermal gas pressure. Local observations of synchrotron and inverse Compton emission from radio galaxies are consistent with equipartition $`P_BP_{rel}`$ (Kaneda et al 1995). I have assumed that the energy injection into relativistic particles is $`10\%`$ of the total kinetic energy of a supernova, so $`P_{rel}<P_{gas}`$ should be a strict upper bound. Thus, $`P_B<P_{gas}`$ gives the upper bound:
$$B<6\left(\frac{n}{1\mathrm{c}\mathrm{m}^3}\right)^{1/2}\left(\frac{T}{10^4K}\right)\mu G$$
(7)
where $`n`$ is the baryon number density (note that gas with temperatures $`>10^{45}`$K will escape from the shallow potential wells of the first proto-galaxies). If the magnetic field exceeds the above value, the over-pressurised lobe will expand on the dynamical time scale until the magnetic pressure drops. Thus, for $`z>5`$, it seems likely that synchrotron energy losses will be unimportant.
Ionization & Cherenkov losses Interactions with the non-relativistic gas will result in energy losses via ionization and Cherenkov emission of plasma waves at a rate independent of the Lorentz factor, $`\dot{E}_{ion}9\times 10^{19}n\mathrm{erg}\mathrm{s}^1`$, which implies that the relative energy loss rate is:
$$\frac{\dot{E}_{IC}}{\dot{E}_{ion}}=120\left(\frac{n}{1\mathrm{c}\mathrm{m}^3}\right)^1\left(\frac{1+z}{10}\right)^4\left(\frac{\gamma }{10^3}\right)^2$$
(8)
Thus, for
$$\gamma >\gamma _{break}100\left(\frac{n}{1\mathrm{c}\mathrm{m}^3}\right)^{1/2}(\frac{1+z}{10})^2$$
(9)
inverse Compton losses are more important than ionization losses. The cutoff Lorenz factor at the lower end $`\gamma _{co}`$, is of interest since the lower end accounts for most of the electrons, both in terms of number and energy: $`N_{electrons}(>\gamma )\gamma ^{2\alpha }=\gamma _{co}^{1.6}`$; $`E_{electrons}(>\gamma )\gamma ^{2\alpha +1}=\gamma _{co}^{0.6}(\alpha =0.8)`$. Note that the electrons that cool via ionization losses do not drop out completely, but merely form a flattened distribution with $`\gamma _f\gamma _i350n(\mathrm{t}/10^7\mathrm{yr})`$. These electrons with lower Lorentz factors could scatter CMB photons to optical and UV frequencies.
Free-free radiation Free-free radiation results in an energy loss rate $`\dot{E}_{freefree}=6\times 10^{22}n\gamma \mathrm{erg}\mathrm{s}^1`$. Thus, free-free radiation only dominates over ionization and Cherenkov radation for $`\gamma >1500`$. However, in this regime inverse Compton losses dominate, since
$$\frac{\dot{E}_{IC}}{\dot{E}_{ff}}=190\left(\frac{n}{1cm^3}\right)^1\left(\frac{1+z}{10}\right)^4\left(\frac{\gamma }{10^3}\right)$$
(10)
Thus, free-free emission is never important in cooling relativistic electrons at high redshift.
In summary, the inverse Compton emission is the dominant energy loss mechanism between a lower and upper energy cutoff. From equation (9), the lower frequency break is determined by the competition between the inverse Compton loss rate and ionization and atomic cooling losses, below
$$E_{lower}=\gamma _{break}^2h\nu _{\mathrm{CMB}}=70\left(\frac{\mathrm{n}}{1\mathrm{cm}^3}\right)\left(\frac{1+z}{10}\right)^3\mathrm{eV},$$
(11)
where I assume the seed photon $`\nu _{\mathrm{CMB}}=1.6\times 10^{12}\left(\frac{1+z}{10}\right)`$GHz lies at the peak of the CMB blackbody spectrum. The competition between inverse Compton cooling losses and the rate of energy injection by Fermi acceleration determines the upper energy cutoff. The timescale for losses by inverse Compton radiation is:
$$t_{life}=\frac{E}{\dot{E}}=7.9\times 10^5\left(\frac{1+z}{10}\right)^4\left(\frac{\overline{\gamma }}{300}\right)^1\mathrm{years}$$
(12)
Equating the Fermi acceleration timescale $`t_{acc}r_Lc/v_{sh}^2=1.3\times 10^3\left(\frac{\gamma }{300}\right)\left(\frac{B}{10\mu G}\right)^1\left(\frac{v_{sh}}{2000\mathrm{k}\mathrm{m}\mathrm{s}^1}\right)^2`$yr (where $`r_L`$ is the Larmour gyroradius, and $`v_{sh}`$ is the typical shock velocity) to $`t_{life}`$ (equation (12), I obtain for the maximum Lorentz factor $`\gamma _{max}=7.4\times 10^6\left(\frac{\mathrm{B}}{10\mu \mathrm{G}}\right)^{1/2}\left(\frac{v_{sh}}{2000\mathrm{km}\mathrm{s}^1}\right)\left(\frac{1+z}{10}\right)^2`$ which corresponds to an upper energy cutoff:
$$E_{upper}360\left(\frac{\mathrm{B}}{10\mu \mathrm{G}}\right)\left(\frac{v_{sh}}{2000\mathrm{k}\mathrm{m}\mathrm{s}^1}\right)^2\left(\frac{1+z}{10}\right)^3\mathrm{GeV}.$$
(13)
In section (2.3.2), I derive the form of the spectrum. For the power law spectrum obtained, $`L_\nu \nu ^1`$, the specific luminosity depends only logarithmically on the energy cutoffs: $`L_\nu =\frac{L_{tot}}{\mathrm{log}(\nu _{\mathrm{upper}}/\nu _{\mathrm{lower}})}\nu ^1`$.
#### 2.3.2 Inverse Compton spectrum
The Fermi shock acceleration mechanism for cosmic rays involves a steady growth of particle energy as a particle scatters back and forth across the shock front. This naturally produces a power law electron energy spectrum with energy spectrum $`dn/d\gamma \gamma ^p`$, where $`p=(\chi +2)/(\chi 1)`$ and $`\chi `$ is the compression ratio for the shock (e.g., Jones & Ellison 1991). The final emission spectrum $`L_\nu \nu ^\alpha `$ depends on the electron energy spectrum; the synchrotron or inverse Compton spectral index is $`\alpha =(p1)/2`$. Thus, the spectral index is determined by the shock structure rather than the details of the scattering process. For a strong adiabatic shock $`\chi =4`$ and $`dn/d\gamma \gamma ^2`$; for an isothermal shock, $`\chi 1`$ and $`dn/d\gamma \gamma ^1`$. The assumption of adiabaticity is most appropriate for the Taylor-Sedov phase, when most of the electrons are accelerated. Galactic SNR show a mean radio spectral index $`\alpha =0.5\pm 0.15`$ (Droge et al 1987), which agrees with $`dn/d\gamma \gamma ^2`$. I shall use this as the canonical electron injection spectrum in this paper. The steepening of the diffuse synchrotron emission to $`\alpha 0.8`$ is likely to be due to energy losses as the electrons age. I examine the emission spectrum in detail below. I assume that there is one supernova for every 100 $`M_{}`$ of stars formed, that each supernova liberates $`10^{51}`$erg in kinetic energy, and $`1M_{}`$ of metals.
The equation for the evolution of the electron population is given by (Ginzburg & Syrovatskii 1964, Sarazin 1999):
$$\frac{N(\gamma )}{t}=\frac{}{\gamma }[b(\gamma )N(\gamma )]+Q(\gamma )$$
(14)
where $`N(\gamma )`$ is the number of electrons in the range $`\gamma `$ to $`\gamma +d\gamma `$, the rate of production of new relavistic electrons is given by $`Q(\gamma )`$, and the rate of energy loss of an individual particle is given by $`b(\gamma )d\gamma /dt`$. From equation (12), we see that the electron lifetime is shorter than both the Hubble time for the redshifts under consideration $`z<30`$ and typical timescales for starbursts ($`10^7`$ years). We can thus assume that the electron population quickly reaches a steady state where the energy losses due to inverse Compton scattering (for high $`\gamma `$) and Coulomb collisions (for low $`\gamma `$) balance the injection rate due to supernova explosions, and set the time derivative to zero. Since injected electrons survive for less than a Hubble time, I ignore evolution in loss rate due to the evolution of the CMB energy density. Assuming that each supernova injects a population of electrons with $`Q(\gamma )=Q_o\gamma ^p`$, where $`p2`$, then the steady state solution to equation (14) is given by:
$`N(\gamma )`$ $`=`$ $`4\times 10^{61}\left({\displaystyle \frac{1+z}{10}}\right)^4\left({\displaystyle \frac{\mathrm{SFR}}{1\mathrm{M}_{}\mathrm{yr}^1}}\right)\gamma ^{(p+1)};\gamma >\gamma _{break}`$ (15)
$`N(\gamma )`$ $`=`$ $`4\times 10^{57}\left({\displaystyle \frac{\mathrm{n}}{1\mathrm{cm}^3}}\right)^1\left({\displaystyle \frac{\mathrm{SFR}}{1\mathrm{M}_{}\mathrm{yr}^1}}\right)\gamma ^{(p1)};\gamma <\gamma _{break}`$
where $`\gamma _{break}`$ (equation (9)) is the transition energy between the regimes where ionization and inverse Compton losses dominate. Thus, the electron population flattens by one power at low energies and steepens by one power at high energies, as compared with the distribution function for the injected population. Since $`\alpha =(p1)/2`$, the emitted spectrum flattens by a half power and steepens by a half power at low and high energies respectively. In particular, for $`p=2`$, as is appropriate for adiabatic shocks, $`L_\nu \mathrm{const}`$ at low energies and $`L_\nu \nu ^1`$ at high energies.
The spectrum of inverse Compton radiation is given by (Sarazin 1999):
$$L_\nu =12\pi \sigma _T_1^{\mathrm{}}N(\gamma )𝑑\gamma _0^1J(\frac{\nu }{4\gamma ^2x})F(x)𝑑x$$
(16)
where $`N(\gamma )`$ is the electron energy spectrum, $`J(\nu )=B_\nu (T_{CMB,o}(1+z))`$ is the CMB blackbody spectrum, and $`F(x)=1+x+2x\mathrm{ln}(x)2x^2`$. This differs only by a normalization correction from assuming that all photons reside at the peak of the blackbody spectrum, $`\nu =1.6\times 10^{11}(1+z)`$Hz. This is because the seed photon spectrum is narrow compared to the electron energy distribution. In figure (1), I show example spectra computed with equations (15) and (16). Note that the efficiency of Coulomb cooling depends on the assumed electron density; I assume that $`n_e\delta n_o(1+z)^3`$, where I assume the overdensity to be either $`\delta 200`$ (to mimic the overdensity at the virial radius) or $`\delta 10^4`$ (to mimic the overdensity for collapsed gas in dense star forming regions). At $`z10`$, this assumption corresponds to $`n_e4\times 10^2\mathrm{cm}^3`$ and $`n_e2\mathrm{cm}^3`$ respectively. Also plotted for reference is the spectrum of the same starburst with a Salpeter IMF, assuming an escape fraction of $`1\%`$, and the spectrum of a mini-quasar with spectrum $`L_\nu \nu ^{1.8}`$, normalized to have the same energy release above the Lyman limit as the starburst. Note that the inverse Compton case has the hardest spectrum of all. Since the spectrum is so hard, photoelectric absorption by the host galaxy does not significantly attenuate the ionizing flux. Most UV photons produced will not escape the host galaxy, whereas X-ray photons with $`E>270\left(\frac{N_{HI}}{10^{21}cm^2}\right)\mathrm{eV}`$ where $`N_{HI}`$ is the column density in the host galaxy, will escape unimpeded. Since $`\nu L_\nu \mathrm{const}`$, most of the energy in ionizing photons escape. Note, incidentally, that the escape fraction for UV photons produced by inverse Compton may be significantly higher than photons of the same frequencies produced by stars, as relativistic electrons can disperse from the star forming region, where gas densities are highest and most photon captures take place.
### 2.4 Zero-metallicity star formation
The stellar IMF at high redshift under conditions of low or zero metallicity is unknown. Up to now, I have assumed a Salpeter IMF, in which one supernova expodes for every $`100M_{}`$ of stars formed, and each supernova deposits on average $`10^{50}\left(\frac{ϵ}{0.1}\right)`$erg in relativistic electrons and $`1M_{}`$ of metals. Since the estimated amount of star formation at high redshift is calibrated to the observed IGM metallicity at $`z3`$, it is worth asking whether the energy injected into X-rays per solar mass of metals produced, $`ϵ_Z10^{50}\mathrm{erg}\mathrm{M}_{}^1`$, could change significantly at high redshift. I have argued that variations in the IMF will not effect significant changes in $`ϵ_Z`$ or the ratio of energy emitted in stellar UV ionizing photons to inverse-Compton X-rays, since the same massive OB stars which produce UV ionizing photons also explode as supernova, producing both relativistic electrons and metals.
However, zero-metallicity star formation may result in very different stellar populations from that seen at low redshift. The lack of efficient cooling mechanisms could result in extremely top heavy stellar IMFs (Larson 1998, Larson 1999) and in particular the production of “Very Massive Objects” (VMOs) in the range $`10^210^5M_{}`$ (Carr, Bond & Arnett 1984). Such stars have high effective temperatures and produce a much harder spectrum than ordinary stars (Tumlinson & Shull, 2000, Bromm et al, 2000). Moreover, the distribution of stellar endpoints is very different from a normal IMF (Heger, Woosley & Waters 2000). Stars with masses between 10 and 35 $`M_{}`$ explode as type II supernovae. While it is not known whether metal free stars with masses between 35 and 100 $`M_{}`$ will explode–they might collapse to form black holes–stars with masses $`100250M_{}`$ are disrupted by the pair production instability, once again producing an energetic supernova event and dispersing metals. Stars more massive than $`250M_{}`$ should collapse completely to black holes, without ejecting any metals (unless they eject their envelopes during hydrogen shell burning, in which case it is possible for them to explode). The latter provides an obvious mechanism for seeding supermassive black holes to form AGN. A number of studies of zero-metallicity star formation suggest that the Jeans/Bonner-Ebert mass and hence the lower mass cutoff for star formation is very high, $`M_{}>10^210^3M_{}`$ (Abel, Bryan & Norman 1999, Padoan, Nordlund & Jones 1997).
There are three main observations to make: (i) VMOs with initial stellar masses $`100M_{}<M<250M_{}`$ which are disrupted by the pair-production stability show $`ϵ_Z(PP)ϵ_Z(\mathrm{type}\mathrm{II})`$, which imply that even if these objects are abundant in the early universe, our estimates of the level of inverse Compton X-ray emission do not change strongly. In particular, $`E_{\mathrm{explosion}}=6.3\times 10^{52}\left(\frac{M}{10^2M_{}}\right)^{2.8}\frac{\mathrm{min}[0.13,(1\varphi _L)^{2.8}]}{0.13}\mathrm{erg}`$ while the yield in elements heavier than helium is $`Z_{\mathrm{ej}}=\mathrm{min}[(1\varphi _L),0.5]`$, where $`\varphi _L`$ is the fraction of the initial mass lost during hydrogen burning (Carr, Bond & Arnett, 1984). Thus, $`ϵ_Z(PP)=1.3\times 10^{50}\left(\frac{M}{10^2M_{}}\right)^{1.8}\frac{\mathrm{min}[0.25,(1\varphi _L)^{1.8}]}{0.25}\mathrm{erg}\mathrm{M}_{}^1ϵ_Z(\mathrm{type}\mathrm{II})`$. (ii) Stars which directly collapse to form black holes represent an additional, unaccounted source of ionizing photons, since they are not included in the metal pollution budget (of course, their most important contribution could lie in seeding AGN formation). (iii) Zero metallicity is a singularity: true zero-metallicity stellar populations differ greatly from low metallicity ones. Even the introduction of trace amounts of metals $`Z10^4Z_{}`$ introduces important changes in stellar structure and evolution (Heger, Woosley & Waters 2000). Furthermore, trace amounts of metals drastically reduces the abundance of VMOs by allowing efficient cooling past the 300K barrier imposed by $`\mathrm{H}_2`$ cooling, reducing the Jeans mass and thus the minimum mass for star formation. Prompt initial enrichment is quite plausible: the lifetime of massive stars before they explode to pollute the IGM with metals is $`10^6`$ years, which is a short fraction of the Hubble time $`t_H8\times 10^8(\frac{1+z}{10})^{3/2}`$yrs even at high redshift. Although the degree of metal mixing is uncertain (e.g. Gnedin & Ostriker 1997, Nath & Trentham 1997, Ferrara, Pettini & Shchekinov 2000), note that first star forming regions are very highly biased, and subsequent generations of the halo hierarchy collapse in proto-clusters and filaments very close to the first star clusters. Thus, stars of finite metallicity could quickly predominate, even if the mean metallicity of the universe is close to primordial (Cen & Ostriker 1999). It is therefore possible and even likely that true zero metallicity star formation was confined to a very small fraction of the stars formed at high redshift, and thus negligible in terms of the energy budget for reionization.
## 3 Observational signatures
In this section, to estimate number counts I use the Press-Schechter based high redshift star formation models of Haiman & Loeb (1997), which are normalised to the observed metallicity of the z=3 IGM, $`\mathrm{Z}=10^310^2\mathrm{Z}_{}`$. In these models, in every halo capable of atomic cooling (i.e., with a virial temperature $`T_{vir}>10^4`$K), a fixed fraction of the gas $`f_{star}=1.7,17\%`$ (for $`\mathrm{Z}(z=3)=10^3,10^2\mathrm{Z}_{}`$ respectively) fragments in a starburst lasting $`t_o10^7`$yrs. The correspondence between halo mass and star formation rate is thus $`\mathrm{SFR}2\left(\frac{\mathrm{M}_{\mathrm{halo}}}{10^9\mathrm{M}_{}}\right)\left(\frac{f_{star}}{0.17}\right)\mathrm{M}_{}\mathrm{yr}^1`$, and each halo is only visible for some fraction of the time $`\frac{t_o}{t_H(z)}`$.
### 3.1 HeII recombination lines
One possible signature of the hard spectrum produced by inverse Compton X-rays would be HeII recombination lines from the host galaxy. Indeed, such lines may well be detectable from the first luminous objects with sufficiently hard spectra, such as mini-quasars or metal-free stars (Oh, Haiman & Rees 2000). The different sources may perhaps be distinguished on the basis of line widths and line ratios (Tumlinson, Giroux & Shull, 2000). However, inverse Compton emission produces too few HeII ionizing photons for such recombination lines to be detectable with NGST at high redshift. For a Salpeter IMF, $`\dot{N}_{ion,HI}=10^{53}\left(\frac{\mathrm{SFR}}{1\mathrm{M}_{}\mathrm{yr}^1}\right)\mathrm{photons}\mathrm{s}^1`$ from stellar UV radiation. On the other hand, since secondary ionizations of HeII are negligible (Shull & van Steenberg 1985), the production rate of HeII ionizing photons from non-thermal emission is $`\dot{N}_{ion,HeII}=_{4\nu _L}^{\nu _{thin}}\frac{L_\nu }{h\nu }2.5\times 10^{49}\left(\frac{SFR}{1M_{}yr^1}\right)\mathrm{photon}\mathrm{s}^1`$, where $`\nu _{thin}`$ is the frequency at which the halo becomes optically thin. Thus, $`Q\dot{N}_{\mathrm{ion}}^{\mathrm{HeII}}/\dot{N}_{\mathrm{ion}}^{\mathrm{HI}}2.5\times 10^4`$, as compared with $`Q0.05`$ for stellar emission from metal-free stellar population, and $`Q4^\alpha =0.080.25`$ for a QSO, where $`L_\nu \nu ^\alpha `$ and $`\alpha =11.8`$, implying that the relative flux in HeII and H$`\alpha ,\beta `$ recombination lines is much smaller for inverse Compton emission than metal-free stars or AGN. The luminosity in a helium recombination line $`i`$ may be estimated as $`L_\mathrm{i}=Qf_\mathrm{i}L_{\mathrm{H}\alpha }`$, where $`L_{\mathrm{H}\alpha }`$ is the Balmer $`\alpha `$ line luminosity, $`f_\mathrm{i}\frac{j_i}{j_{H\alpha }}\left(\frac{\alpha _B(HI)n_{HII}}{\alpha _B(HeII)n_{HeIII}}\right)`$, $`ϵ`$ is the fraction of SN energy which emerges as IC emission, and $`j_i,j_{H\alpha }`$ may be obtained from Seaton (1978). The observed flux is $`J_i=\frac{L_\mathrm{i}}{4\pi d_\mathrm{L}^2}\frac{1}{\delta \nu }=0.04\left(\frac{q_\mathrm{i}}{0.5}\right)\left(\frac{1+z}{10}\right)^1\left(\frac{ϵ}{0.1}\right)\left(\frac{R}{1000}\right)\left(\frac{\mathrm{SFR}}{1\mathrm{M}_{}\mathrm{yr}^1}\right)\mathrm{nJy}`$ where $`\mathrm{R}`$ is the spectral resolution, and $`q_\mathrm{i}f_\mathrm{i}\nu _{\mathrm{H}\alpha }/\nu _\mathrm{i}`$, where $`f_{\lambda 4686}=0.74,f_{\lambda 1640}=4.7,f_{\lambda 3203}=0.30`$. This is undetectable by NGST, which in $`10^5`$s integration time requires for a 10 $`\sigma `$ detection a flux $`F(10\sigma )30`$nJy at these frequencies (observed wavelengths $`1<\lambda <5.5\mu `$m) and spectral resolution $`R1000`$ (see Oh, Haiman & Rees 2000 for details).
It is worth mentioning that a hard source may produce relatively little recombination line flux and yet play an important role in reionizing the universe. The recombination line flux from a source is $`(1f_{esc})`$, whereas the ionizing flux escaping into the IGM is $`f_{esc}`$ (where $`f_{esc}`$ is significantly higher for hard sources: all photons with $`E>E_{halo,thin}=270\left(\frac{N_{HI}}{10^{21}\mathrm{cm}^2}\right)`$eV can escape freely from the host). Secondary ionizations are unimportant within a host halo: much of the gas in fully ionized, in which case the energetic photoelectron deposits its energy as heat. In any case, the halo is optically thin to the most energetic photons $`E>E_{halo,thin}`$ which are important for secondary ionizations. Thus, in recombination line flux what matters is the total number of ionizing photons produced, whereas for the reionization of the IGM what matters is the total output energy (since in a largely neutral medium with $`x_e<0.1`$ the total number of ionizations for a photon of energy $`E_{photon}`$, including secondary ionizations, is $`N_{ion}E_{photon}/37\mathrm{e}\mathrm{V}`$ (Shull & van Steenberg 1985). Note, however, that as the medium becomes more ionised an increasing fraction of the energy is deposited as heat, and an additional source of soft photons is needed for reionization to proceed (Oh 2000a)). Thus, for instance, for $`f_{esc}1\%`$, inverse Compton radiation makes a negligible contribution to the H$`\alpha `$ luminosity, $`\frac{L_{H\alpha }^{IC}}{L_{H\alpha }^{stellar}}\frac{\dot{N}_{ion}^{IC}}{\dot{N}_{ion}^{stellar}}=2.5\times 10^4`$, but the ionizing IC and stellar radiation escaping from the host are energetically comparable (from equations (1) and (3), $`L_X^{IC}L_{UV}^{stellar}`$), and thus they produce roughly equal number of ionizations in the IGM.
### 3.2 X-rays and gamma-rays
Unfortunately, direct detection of inverse Compton X-rays from high-redshift star clusters is unlikely. From equation (1), the observed flux in X-rays from a star forming region is:
$$f=\frac{L}{4\pi d_L^2}=5\times 10^{20}\left(\frac{ϵ}{0.1}\right)\left(\frac{\mathrm{SFR}}{1\mathrm{M}_{}\mathrm{yr}^1}\right)\left(\frac{1+z}{10}\right)^2\mathrm{erg}\mathrm{s}^1\mathrm{cm}^2$$
(17)
By contrast, the Chandra X-ray Observatory (CXO) sensitivity (see http://chandra.harvard.edu/) for a 5$`\sigma `$ detection of a point source in an integration time of $`10^5`$s in the 0.2–10 keV range is $`F_X=4\times 10^{15}\mathrm{erg}\mathrm{s}^1\mathrm{cm}^2`$. The next generation X-ray telescope, Constellation-X (http://constellation.gsfc.nasa.gov/) is optimised for X-ray spectroscopy and will not have significantly greater point source sensitivity: for a similar bandpass and integration time it will be able to detect objects out to a flux limit of $`F_X=2\times 10^{15}\mathrm{erg}\mathrm{s}^1\mathrm{cm}^2`$. Thus, in a week CXO and Constellation-X will at best be able to detect starbursts (with seed photons provided by the IR radiation field from dusty star forming regions) with star formation rates SFR $`100\mathrm{M}_{}\mathrm{yr}^1`$ (corresponding to $`L_X4.8\times 10^{42}\mathrm{erg}\mathrm{s}^1`$) out to redshift $`z1`$ (note that the Lyman-break galaxies detected at $`z3`$ (Steidel et al 1996) typically have have inferred star formation rates $`\mathrm{SFR}100\mathrm{M}_{}\mathrm{yr}^1`$), while the very brightest starbursts, with star formation rates of SFR$`1000\mathrm{M}_{}\mathrm{yr}^1`$, might be detectable out to $`z3`$. As discussed in section 3.4, simultaneous detection of synchrotron radiation with the Square Kilometer Array will then constrain magnetic field strengths in these objects.
Detection in gamma ray emission is also unlikely. The upcoming Gamma-ray Large Area Space Telescopy (GLAST) (http://glastproject.gsfc.nasa.gov/) will have a flux sensitivity of $`2\times 10^9\mathrm{photons}\mathrm{cm}^2\mathrm{s}^1`$ in the 20MeV–300GeV range, whereas star-forming regions at high redshift will have fluxes of at most $`F_\gamma 10^{15}\left(\frac{ϵ}{0.1}\right)\left(\frac{\mathrm{SFR}}{1\mathrm{M}_{}\mathrm{yr}^1}\right)\left(\frac{1+z}{10}\right)^2\mathrm{photons}\mathrm{cm}^2\mathrm{s}^1`$.
What fraction of the unresolved X-ray and gamma-ray background could be due to star formation at high redshift $`z>3`$? Observations of absorption in CIV and other metals in Ly$`\alpha `$ forest absorption lines with $`10^{14.7}\mathrm{cm}^2<N_{HI}<10^{16}\mathrm{cm}^2`$ indicate that $`Z10^2Z_{}2\times 10^4`$. Each supernova produces about $`1M_{}`$ of metals (Woosley & Weaver 1995). Assuming $`(\mathrm{\Omega }_b,h)=(0.04,0.65)`$, this implies one supernova every $`1000\mathrm{kpc}^3`$ (comoving). If each supernova injects $`10^{50}`$ ergs in hard X-rays, the comoving X-ray energy density was $`U_X\frac{c}{4\pi }\frac{U_X}{(1+\overline{z})}2\times 10^6\mathrm{eV}\mathrm{cm}^3`$. If the mean source redshift was $`\overline{z}5`$, the distant sources produced a diffuse flux of $`J1\mathrm{k}\mathrm{e}\mathrm{V}\mathrm{s}^1\mathrm{cm}^2\mathrm{sr}^1`$. One can make a more detailed estimate using the equation of cosmological transfer (Peebles 1993):
$$J(\nu _o,z_o)=\frac{1}{4\pi }_{z_o}^{\mathrm{}}𝑑z\frac{dl}{dz}\frac{(1+z_o)^3}{(1+z)^3}ϵ(\nu ,z)e^{\tau _{\mathrm{eff}}(\nu _o,z_o,z)}$$
(18)
where $`z_o`$ is the observer redshift, $`\nu =\nu _o(1+z)/(1+z_o)`$, and $`ϵ(\nu ,z)`$ is the comoving X-ray emissivity. Using the star formation model of Haiman & Loeb (1997), I obtain the estimate:
$$\nu J_\nu =0.8\left(\frac{ϵ}{0.1}\right)\left(\frac{Z(z=3)}{10^2Z_{}}\right)\mathrm{keV}\mathrm{s}^1\mathrm{cm}^2\mathrm{sr}^1\mathrm{for}z_{source}3$$
(19)
In Fig (2), I display this predicted level of emission against the observed X-ray and gamma-ray background in the 1 keV –100 GeV range, from analytic fits to the ASCA and HEAO A2,A4 data in the 3-60 keV range (Boldt 1987), the HEAO 1 A-4 data in the 80-400 keV range (Kinzer et al 1997), the COMPTEL data in the 800 keV – 30MeV range (Kappadath et al 1996). Data points from EGRET in the 30 MeV – 100 GeV range (Sreekumar et al 1998) are also shown. The level of the unresolved background of course depends on the sensitivity of the instrument; recently CXO resolved $`80\%`$ of the hard X-ray background in the 2-10 keV range into point sources (Mushotzky et al 2000). This agrees well with the predictions of XRB synthesis models (e.g., Madau, Ghisellini & Fabian 1994), which use AGN unification schemes to reproduce the observed spectral shape of the XRB. A prediction of the IC scenario presented here is that a non-trivial fraction of the X-ray/gamma-ray background will not be resolved into point sources with upcoming missions, due to the extreme faintness of high redshift sources.
It is particularly intriguing that both the amplitude and spectral shape of the gamma-ray background as observed by EGRET is well-matched by the predicted level of gamma-ray emission in this model. This raises the exciting possibility that the majority of the observed gamma-ray background comes from inverse Compton emission at high redshifts. At present, the origin of the gamma-ray background is still unknown. The most favoured scenario for some time was that it is due to unresolved gamma-ray blazars (Bignami et al 1979, Kazanas & Protheroe 1983, Stecker & Salamon 1996): the observed blazar $`\gamma `$-ray spectrum has an average spectral index compatible with the observed GRB (Chiang & Mukherjee 1998). However, extrapolation of the observed EGERT blazar luminosity function implies that unresolved blazars can account for at most $`25\%`$ of the diffuse $`\gamma `$-ray background (Chiang & Mukherjee 1998). The unresolved blazar model will most likely be decisively tested by GLAST (Stecker & Salamon 1999), which will be two orders of magnitude more sensitive than EGRET. A host of other models include pulsars expelled into the halo by asymmetric supernova explosions (Dixon et al 1998, Hartmann 1995), primordial black hole evaporation (Page & Hawking 1976), supermassive black holes at very high redshift (Gnedin & Ostriker 1992), annihilation of weakly interactive big bang remnants (Silk & Srednicki 1984, Rudaz & Stecker 1991) and finally, inverse Compton radiation from cosmic ray electrons in our own Galaxy (Strong & Moskalenko 1998, Dar & De Rujula 2000), and from collapsing clusters (Loeb & Waxman 2000). However, to date the possibility of inverse Compton emission from high redshift supernovae has not been discussed.
A scenario in which the majority of the gamma-ray background comes from inverse Compton emission at high redshift make a number of firm predictions: (i) As previously mentioned, the majority of the GRB will remain unresolved by GLAST, due to the extreme faintness of the contributing sources. (ii) After removal of the Galactic contribution, which is correlated with the structure of our Galaxy and our position within it (Dixon et al 1998, Dar et al 1999), a highly isotropic component of the GRB will still be present. (iii) The GRB should be extremely smooth, and exhibit significant fluctuations only at extremely small angular scales. The fluctuations should be dominated by the Poisson rather than the clustering contribution (see Oh 1999). Using $`S^n=_0^{S_c}\frac{dN}{dS}S^n`$, where $`S_c`$ is the cut-off flux for point source removal, and the star formation model of Haiman & Loeb (1997), I find that for $`z>3`$ sources (which are too faint to be removed as point sources), $`\frac{I^2^{1/2}}{I}=3.4\times 10^2\left(\frac{\theta }{5^{}}\right)^1`$, where $`I`$ is surface brightness (the angular resolution of GLAST is expected to be of order $`15`$ arcmin). (iv) The gamma-ray background at $`E>100`$GeV should be attenuated, due to pair production opacity against IR/UV photons (e.g., Salamon & Stecker 1998, Oh 2000b). High energy photons initiate an electromagnetic cascade which transfers energy from high energy photons to the lower energy portion of the spectrum, where the universe is optically thin (Coppi & Aharonian 1997). Note that the EGRET spectrum was directly determined with data only up to 10 GeV; beyond 10 GeV larger uncertainties exist due to backsplash in the NaI calorimeter, and Monte-Carlo simulations were used to determine the differential flux in the 10-30,30-50 and 50-120 GeV range (Sreekumar et al 1998). Thus, the EGRET data points with $`E>10`$GeV are less reliable. It would be intriguing to see if GLAST (sensitive out to 300 GeV) indeed shows an absorption edge to the gamma-ray background at higher energies. It would also be interesting to look for absorption of the gamma-ray background in a line of sight passing through a massive cluster (indicating that the gamma-rays come from higher redshifts than the cluster); the pair-production optical depth through a cluster is of order $`\tau 2n_\gamma \sigma _Tr_{vir}0.4`$ (assuming $`n_\gamma 0.1\mathrm{cm}^3`$ and $`r_{vir}1`$Mpc).
### 3.3 CMB constraints
Will the upscattering of CMB photons by relativistic electrons at high redshift produce an observable signal in the CMB, or violate any present observational constraints? When the electrons are relativistic, CMB photons are inverse Compton scattered to such high energies (completely out of the detector bandpass) the process may simply be thought of as absorption, $`\delta I/I\tau _e^{rel}`$, where $`\tau _e^{rel}`$ is the optical depth of relativistic electrons, and $`I`$ is the number flux of CMB photons. The flux decrement due to the absorption of CMB photons may be estimated as:
$$\mathrm{\Delta }S_\nu J_{CMB}𝑑\mathrm{\Omega }𝑑ln_e^{rel}\sigma _T=J_{CMB}\sigma _T\frac{N_e^{rel}}{d_A^2}$$
(20)
where $`J_{CMB}`$ is the blackbody surface brightness of the CMB, and $`N_e^{rel}`$ is the total number of relativistic electrons in the system. The steady state number of relativistic electrons is given by $`N_e^{rel}=𝑑\gamma N(\gamma )`$ where $`N(\gamma )`$ is given by equation (15). Since the CMB flux peaks at $`\nu 10^{11}`$Hz, the absolute magnitude of the flux decrement is maximized by going to similar frequencies. At 20 GHz, the highest frequency detectable by SKA, the SKA has an rms sensitivity of $`6`$nJy for a $`10^5`$s integration. On the other hand, the flux decrement is $`\mathrm{\Delta }S_\nu (20\mathrm{GHz})3.5\times 10^5(\frac{1+z}{10})^2(N_e^{rel}/3\times 10^{58}\mathrm{electrons})`$ nJy, which is unobservably small.
One might hope to detect the mean signal of all the CMB photons upscattered by relativistic electrons. In particular, since the number of CMB photons is no longer conserved, the absorption might be detectable as a chemical potential distortion of the CMB. Let us estimate the number of CMB photons destroyed. Each supernova upscatters at most $`N_{scattered}ϵE_{SN}/E_X10^{60}`$ CMB photons, where I have set the average photon energy $`E_X100`$eV (note that since the number of photons $`N_\nu \nu ^1`$, most of the upscattered photons are of low energy). For a metallicity of $`Z10^2Z_{}`$ at $`z3`$, one supernova has gone off every comoving $`V_{SN}1000\mathrm{kpc}^3`$, and thus the comoving number density of upscattered CMB photons is $`\delta nN_{scattered}/V_{SN}4\times 10^8\mathrm{cm}^3`$. Since $`n_\gamma 400\mathrm{cm}^3`$, we have $`\delta n/n10^{10}`$, which results in an undetectably small chemical potential distortion. Thus, the upscattering of CMB photons at high redshift does not violate any distortion constraints on the CMB.
If the IGM is reionized inhomogeneously, as in canonical models, then secondary CMB anistropies will be created by CMB photons Thompson scattering off moving ionized patches (Agahanim et al 1996, Grusinov & Hu 1998, Knox et al 1998). The power spectrum is generally white noise, with $`\mathrm{\Delta }T/T10^610^7`$, peaking at arc-minute to sub arc-minute scales. However, if the IGM is reionized fairly homogeneously by X-rays then over a line of sight the positive and negative contributions of the velocity field will cancel out. In this case, only the second-order Ostriker-Vishniac effect due to coupling between density and velocity fields will be present. A null detection of the inhomogeneous reionization anisotropy could place upper limits on the patchiness of reionization, although this will be a difficult measurement as the inhomogeneous reionization and Ostriker-Vishniac signals are likely to be of comparable strength (Haiman & Knox 1999).
### 3.4 Radio observations
The supernovae that exploded generate magnetic turbulence and magnetic fields, which allow Fermi acceleration to take place (Jones & Ellison 1991, Blandford & Eichler 1987); observations of local radio galaxies in non-thermal radio and X-ray emission yield field strengths consistent with equipartition between relativistic particles and the magnetic field (Kaneda et al 1995). If such magnetic fields are present in the first star clusters, the relativistic electron population is a source of synchrotron radio emission as well as inverse-Compton emission. Below I find that for a given relativistic electron population, the Square Kilometer Array (SKA) will be much more sensitive to non-thermal radio emission than CXO or Constellation-X will be to inverse Compton X-ray emission. Radio observations will thus allow one to establish the presence of relativistic electrons in objects too faint to observe directly in X-ray emission. Since one knows the CMB energy density exactly as a function of redshift, given reasonable assumptions for the magnetic field the observed radio emission allows one to immediately estimate the amount of inverse Compton X-ray emission which must be taking place,
$$\mathrm{L}_\mathrm{X}=\frac{\mathrm{U}_{\mathrm{CMB}}(\mathrm{z})}{\mathrm{U}_\mathrm{B}}\mathrm{L}_{\mathrm{synch}}$$
(21)
The transition Lorentz factor $`\gamma _{break}`$ at which electron energy losses are dominated by inverse Compton rather than ionization losses is given by equation (9). Above an observed radio frequency $`\nu _{break}=\nu _L\gamma _{break}^2/(1+z)=2.8\times 10^4\left(\frac{\mathrm{B}}{10\mu \mathrm{G}}\right)\left(\frac{\mathrm{n}}{1\mathrm{cm}^3}\right)\left(\frac{1+z}{10}\right)^5`$ Hz (where $`\nu _L\frac{eB}{2\pi m_ec}`$ is the electron gyrofrequency), the steady state electron population is determined by balance between supernova injection and inverse Compton losses, and is given by equation (15). From standard formulae for synchrotron emissivity (Rybicki & Lightman 1979) $`ϵ_\nu ^{synch}=\sigma _T\gamma ^2\beta ^2U_{mag}c\mathrm{n}(\mathrm{E})\mathrm{dE}/\mathrm{d}\nu `$ (where n(E) is the number density of emitting electrons in $`dE`$) this yields a synchrotron luminosity:
$$L_\nu ^{sync}=2.7\times 10^{28}\left(\frac{\mathrm{B}}{10\mu \mathrm{G}}\right)^{1+\alpha }\left(\frac{\nu }{1\mathrm{GHz}}\right)^\alpha \left(\frac{\mathrm{SFR}}{1\mathrm{M}_{}\mathrm{yr}^1}\right)\left(\frac{ϵ}{0.1}\right)f(z,U_\gamma )\mathrm{erg}\mathrm{s}^1\mathrm{Hz}^1$$
(22)
where $`f(z,U_\gamma )=\mathrm{min}[\left(\frac{1+z}{10}\right)^4,\left(\frac{U_\gamma }{4.2\times 10^9\mathrm{erg}\mathrm{cm}^3}\right)^1]`$ (the latter term in brackets is used if the stellar radiation field has a higher energy density than the CMB). By contrast, the thermal free-free emission is given by (Oh 1999):
$$L_\nu ^{ff}=1.2\times 10^{27}\left(\frac{\mathrm{SFR}}{1\mathrm{M}_{}\mathrm{yr}^1}\right)\mathrm{erg}\mathrm{s}^1\mathrm{Hz}^1$$
(23)
Thus, assuming $`\alpha =1`$, at observed frequencies $`\nu <\nu _{trans}=2.3\left(\frac{1+z}{10}\right)^1\left(\frac{\mathrm{B}}{10\mu \mathrm{G}}\right)^2\left(\frac{ϵ}{0.1}\right)f(z,U_\gamma )\mathrm{GHz}`$ synchrotron emission dominates over free-free emission. This is well within the 0.1–20 GHz capability of the SKA. However, this is only true if the power law for electron population extends to high energies. If it is truncated at some maximum Lorentz factor $`\gamma _{max}`$, then synchrotron emission is only observable for $`\nu <\nu _{max}=\nu _L\gamma _{max}^2=0.28\left(\frac{1+z}{10}\right)^1\left(\frac{\mathrm{B}}{10\mu \mathrm{G}}\right)\left(\frac{\gamma _{max}}{10^4}\right)^2\mathrm{GHz}`$. If $`\nu _{max}<\nu _{trans}`$, this will be manifested by an abrupt drop in radio flux at $`\nu _{max}`$, beyond which the emission takes the flat spectrum free-free emission form (this does not take place for expected values of $`\gamma _{max}=7.4\times 10^6\left(\frac{\mathrm{B}}{10\mu \mathrm{G}}\right)^{1/2}\left(\frac{v_{sh}}{2000\mathrm{km}\mathrm{s}^1}\right)\left(\frac{1+z}{10}\right)^2`$; see equation (13)). Observation of such a drop will yield valuable constraints on $`B,\gamma _{max}`$ in the first objects; the value of $`\gamma _{max}`$ in turn constrains the upper energy cutoff in the inverse Compton X-ray/gamma-ray spectrum, which could be important in determining whether high-redshift starbursts make a significant contribution to the gamma-ray background.
Can radio emission from star clusters at high redshift be detected by the proposed Square Kilometer Array? Non-thermal emission can be distinguished from free-free emission with multi-frequency observations to identify frequency regimes where the spectral slope is steep.The SKA detector noise may be estimated as:
$$S_{\mathrm{instrum}}=\frac{2kT_{\mathrm{sys}}}{A_{\mathrm{eff}}\sqrt{2t\mathrm{\Delta }\nu }}=25\left(\frac{\mathrm{\Delta }\nu }{160\mathrm{MHz}}\right)^{1/2}\left(\frac{t}{10^5s}\right)^{1/2}\mathrm{nJy}$$
(24)
where I have used $`A_{\mathrm{eff}}/T_{\mathrm{sys}}=2\times 10^8\mathrm{cm}^2/\mathrm{K}`$ for the SKA (Braun et al 1998), and I assume a bandwidth $`\mathrm{\Delta }\nu 0.5\nu `$. The flux density due to non-thermal emission from a high-redshift star cluster, assuming $`\alpha =1`$, is:
$$S_{\nu _o}=\frac{L_\nu (\nu _o(1+z)}{4\pi d_L^2}(1+z)=8.5\left(\frac{\mathrm{B}}{10\mu \mathrm{G}}\right)^{1+\alpha }\left(\frac{\nu }{320MHz}\right)^\alpha \left(\frac{\mathrm{SFR}}{1\mathrm{M}_{}\mathrm{yr}^1}\right)\left(\frac{ϵ}{0.1}\right)\left(\frac{1+z}{10}\right)^2f(z,U_\gamma )\mathrm{nJy}$$
(25)
Thus, a source with SFR$`10\mathrm{M}_{}\mathrm{yr}^1`$ at $`z9`$ can be detected as a 10$`\sigma `$ detection in 10 days. To estimate the number of sources detectable by SKA, I use the Press-Schechter based high redshift star formation models of Haiman & Loeb (1997), and define an efficiency factor $`f_{radio}=\left(\frac{\mathrm{B}}{10\mu \mathrm{G}}\right)^{1+\alpha }\left(\frac{ϵ}{0.1}\right)\left(\frac{f_{star}}{0.17}\right)`$, where $`f_{star}`$ is the fraction of halo gas which fragments to form stars. In Figure (3), I display the number of sources above a given redshift which may be detected in non-thermal emission in the $`1^{}`$ SKA field of view, assuming $`f_{radio}=1,0.01`$. Also shown is the number of sources which can be detected in free-free emission at 4 GHz. One should be able to detect a large number of sources at high redshift, $`z>5`$. Thus, radio observations of non-thermal emission can serve as a useful proxy for X-ray observations in allowing one to estimate $`L_X`$, and thus the overal level of X-ray emission at high redshift.
### 3.5 Multi-wavelength observations
While a detailed study of high-redshift multi-wavelength campaigns is beyond the scope of this paper, below I describe some possible follow-up observations if synchrotron emission is detected at high redshift.
* Redshift estimates Thus far the most efficient way to select high-redshift radio galaxies has proven to be the observation of steep radio spectra $`\alpha <1.3`$ (Chambers et al 1990, van Breugel et al 1999). This is at least partially due to the fact that extremely bright radio galaxies have SEDs which steepen with frequency; the k-correction then implies that sources at increasing redshift have steeper observed spectra. Steepening with redshift due to inverse Compton losses, as well as selection effects (i.e., brighter sources have stronger magnetic fields, and thus more rapid synchrotron losses) could also play a role (Krolik & Chen 1991). This technique may fail for fainter sources at high redshift since for $`\nu >\mathrm{min}(\nu _{trans},\nu _{max})`$, the radio flux will be dominated by free-free emission and the spectra will appear flat. For instance, for $`B<3\left(\frac{1+z}{10}\right)^{2.5}\left(\frac{ϵ}{0.1}\right)^{0.5}\mu `$G, we have $`\nu _{trans}<160`$MHz and SKA will only detect free-free emission within its frequency coverage. An efficient way to select high-redshift objects prior to reionization would be to perform broad-band deep field imaging with NGST and select Lyman-break dropouts as has been done at $`z3`$ (Steidel et al 1996); before the epoch of reionization one must select ’Gunn-Peterson dropouts’, i.e. galaxies with no flux shortward of rest-frame HI Ly$`\alpha `$. Yet another method of selecting high-redshift objects would be to perform a joint survey in the submm with the Atacama Large Millimeter Array (ALMA); at low redshift the submm dust emission scales almost linearly with other star formation indicators, such as radio and UV emission. However, for $`z0.510`$ and a given dust emission SED the K-correction almost balances the cosmological dimming of a source, implying a flux density almost independent of redshift (e.g., Blain et al 2000). This has spawned suggestions to obtain approximate redshifts by the flux density ratio between submm and radio wavebands (Carilli & Yun 1999, 2000), although uncertainties include an AGN contribution to the radio flux and the dust temperature in the galaxy, with a degeneracy between hotter galaxies at high redshift and cooler ones nearby (Blain 1999).
* Relative importance of X-ray and UV stellar emission for reionization As previously noted, from a measurement of the radio synchrotron flux we can use equation (21) to estimate the level of inverse Compton X-ray emission which must be taking place. Redshifts and thus $`\mathrm{U}_{\mathrm{CMB}}(\mathrm{z})`$ can be determined with Balmer line spectroscopy with NGST, while $`\mathrm{U}_\mathrm{B}`$ can be estimated by assuming a magnetic field strength required to minimise the total energy density of the system (this is close to the value for energy equipartition between relativistic particles and magnetic fields). Observations of local radio galaxies in non-thermal radio and X-ray emission yield field strengths consistent with the minimum energy value (Kaneda et al 1995). NGST can constrain the rest-frame UV emission, and a joint measurement of the rest-frame UV flux longward of Ly$`\alpha `$ and the Balmer line flux (where $`\mathrm{L}_{\mathrm{H}\alpha }(1f_{esc})\dot{N}_{ion}`$) can constrain $`f_{esc}`$, the escape fraction of ionizing photons, where one roughly expects $`J_{H\alpha }40(1f_{esc})J_{IR}(R/1000)`$. Thus, joint SKA/NGST observations could place limits on whether inverse Compton X-rays or stellar UV photons were energetically dominant and thus more important in reionizing the universe.
* Measuring magnetic fields in bright sources The very brightest sources out to $`z3`$ will also be visible in X-ray emission with CXO; the fact that the same population of relativistic electrons is responsible for non-thermal X-ray and radio emission can be confirmed by comparing respective spectral slopes. In this case, one can estimate the strength of magnetic fields, $`\mathrm{U}_\mathrm{B}\frac{L_X}{L_{sync}}\mathrm{U}_{\mathrm{CMB}}(\mathrm{z})`$. Since dust obscuration is unimportant at both hard X-ray and radio wavelengths, the measurement should be fairly robust.
* Distinguishing between synchrotron emission from AGN and supernovae To confirm that such emission arises from star-forming regions rather than an AGN, one might look for signs of diffuse emission (note that the angular resolution of the SKA is $`0.1^{\prime \prime }`$, while the angular scale of the virial radius of typical objects will be $`\theta _{vir}0.5^{\prime \prime }(M/10^9M_{})^{1/3}`$). AGNs can be selected on the basis of color, as has been successfully carried out at lower redshifts (Fan 1999), using broad band NIR and MIR imaging with NGST. Finally, the line widths of the H$`\alpha `$, H$`\beta `$ lines as observed with NGST may be considerably broader for an AGN ($`\sigma 1000\mathrm{km}\mathrm{s}^1`$), due to line broadening by the accretion disc.
## 4 Conclusions
X-ray emission from a early star forming regions is predicted to be large and energetically comparable to UV emission. Non-thermal inverse Compton emission, which provides a good fit to local observations and should become increasingly important at high redshift, due to the evolution of the CMB energy density, is predicted to be the dominant source of X-rays. It introduces a whole host of physical consequences: the topology of reionization changes, becoming more homogeneous with much fuzzier delineation between ionized and neutral regions; reheating temperatures increase, with implications for feedback on structure formation and the observed width of Ly$`\alpha `$ forest lines; the abundance of free electrons in dense regions increases, promoting gas phase $`H_2`$ formation, cooling, and star formation. These effects will be considered in a companion paper (Oh 2000a). While direct detection of individual sources with CXO or Constellation-X appears difficult, we can hope to confirm the presence of relativistic electrons in high redshift objects by detecting non-thermal radio emission with the Square Kilometer Array. Given the CMB energy density at that epoch, this yields a minimal level of inverse Compton X-ray emission. Combined with NGST observations of rest frame UV emission, this will determine if stellar radiation or inverse Compton X-rays were the dominant factor in reheating and reionizing the universe. In addition, in this scenario a non-trivial fraction of the hard X-ray and gamma-ray background comes from inverse Compton emission at high redshift. In particular, it is possible to reproduce both the shape and amplitude of the gamma-ray background observed by EGRET, and predict that the majority of the smooth, isotropic gamma-ray background will remain unresolved by GLAST, which should display attenuation above $`100`$GeV, from the pair production opacity due to ambient UV/IR radiation fields.
Many of the conclusions in this paper depend upon a scenario in which the escape fraction of UV ionizing photons is small, $`f_{esc}<10\%`$. This assumption is well supported by observations in the local universe, as well as theoretical radiative transfer calculations of high redshift star forming regions, which predict very low escape fractions, $`f_{esc}0.01`$ (Wood & Loeb 1999, Ricotti & Shull 1999, although note that the latter authors predict a substantial escape fraction at low masses $`10^7M_{}`$). However, note that the latter ignore gas clumping and the multi-phase structure of the ISM. If for some reason the escape fraction is unexpectedly high (e.g. the supernovae blow holes in the ISM through which UV photons can escape), then the X-ray component is energetically subdominant and stellar UV radiation dominates the reionization of the universe. Even in this regime, the X-ray component still plays a role in heating the gas above 15 000 K, He II reionization, and promoting gas phase $`\mathrm{H}_2`$ formation, as these tasks cannot be accomplished by soft photons. The escape fraction of UV ionising photons in high redshift objects may eventually be deduced by comparing the IR and $`H\alpha `$ fluxes observed by NGST (Oh 1999).
The most uncertain aspect of this paper is the assumed level of X-ray luminosity, $`L_X0.1\dot{E}_{SN}`$. I have calibrated the conversion rate via local X-ray and gamma-ray observations of starburst galaxies and individual supernova remnants within our galaxy, and argued that efficiency of all proposed X-ray production mechanisms either remains constant or increases with redshift. If the X-ray emission in local starbursts is primarily due to inverse Compton scattering of soft IR photons (Moran & Lehnart 1997, Moran, Lehnart & Helfand 1999), the empirical relation between star formation rate and X-ray luminosity (equation (1)) implies an electron acceleration efficiency of $`ϵ10\%`$, which lies at the upper limit of theoretical expectations. If instead the empirical relation (1) is approximately correct but the X-rays arise from a variety of emission mechanisms, the importance of X-rays for reionization still holds, but specific observational tests which rely on the inverse-Compton mechanism, such as the gamma-ray background observations(section (3.2)) and observations of radio synchrotron emission with the SKA (section (3.4)) will fail.
## 5 Acknowledgements
I am very grateful to my advisor David Spergel for his encouragement and advice. I also thank Roger Blandford, Andrea Ferrara, Zoltan Haiman and Ed Turner for helpful conversations, and Bruce Draine and Michael Strauss for detailed and helpful comments on an earlier manuscript, and the anonymous referee for helpful comments. I thank the Institute of Theoretical Physics, Santa Barbara for its hospitality during the completion of this work. This work is supported by the NASA ATP grant NAG5-7154, and by the National Science Foundation, grant number PHY94-07194.
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# Galaxy Distances via Rotational Parallaxes
## 1 Introduction
At the moment, astrometry is undergoing a quiet revolution. ESA’s Hipparcos mission set the stage. Currently there are two Hipparcos++ missions in preparation: USNO leads a collaboration to build FAME<sup>1</sup><sup>1</sup>1http://aa.usno.navy.mil/fame/, a MIDEX type mission to be launched in 2004, and Germany is planning the DIVA<sup>2</sup><sup>2</sup>2http://www.aip.de/groups/DIVA/ mission. These spacecraft will extend the reach of astrometry to about two kpc from 100 parsec. Space interferometry is the next step, with the ultimate goal to detect and characterize Earth-like planets around nearby stars (NASA’s TPF and ESA’s DARWIN missions). Before that, NASA’s Space Interferometry Mission (SIM) can determine the distances to virtually any star in the Milky Way with an accuracy of a few percent. Many important problems in Milky Way research can be solved with micro-arcsecond astrometric data.
SIM will also contribute significantly to the field of cosmology. Globular clusters play an important role in cosmology in that they contain the oldest stars known to mankind. However, precise age determinations are currently hampered by the lack of accurate distances. SIM will provide such distances, and will hence establish a firm lower limit to the age of the universe. SIM would contribute more directly to cosmology if its data could be used to establish the parallax of nearby spirals and hence provide a calibration of the zero-point of the Tully-Fisher (TF) relation. The Tully-Fisher relation is one of the tools to measure the expansion rate (H<sub>0</sub>) of the universe. A SIM-based determination of the zero-point of the TF relation would thus directly yield a determination of H<sub>0</sub> with an accuracy of several percent, an order of magnitude improvement over the current state of affairs.
Unfortunately, SIM’s phenomenal precision is still not quite good enough to achieve this. However, other, only slightly less direct methods can be applied. In this paper we describe one of those techniques: the rotational parallax method.
## 2 The Rotational Parallax Method: The Circular Orbit Case
Imagine a nearby spiral galaxy at distance, $`D`$ (in Mpc), inclined with respect to the line of sight (by $`i`$ degrees) that has a rotational speed of $`V_c`$ $`\text{km\hspace{0.17em}s}^1`$. The Andromeda galaxy (M 31, NGC 204) is such a galaxy ($`D0.77,i77^o`$). Its rotational speed ($`V_c270`$) induces a proper motion of $`\mu =\frac{V_c}{\kappa D}73.9\mu \text{as\hspace{0.17em}yr}^1`$. Here $`\kappa `$ is a constant that arises from the choice of units, and $`\kappa 4.74`$ if $`D`$ is measured in Mpc and $`\mu `$ in micro-arcsec per year. Such proper motion is easily resolvable with an instrument such as SIM.
In Fig. Galaxy Distances via Rotational Parallaxes we present a sketch of how a typical nearby spiral might appear on the sky (bottom panel). In the plane of the galaxy (top panel) we use two coordinate systems: rectangular ($`x`$ and $`y`$) and polar ($`R`$ and $`\theta `$). Projected on the plane of the sky, the $`x`$ and $`y^{}`$ coordinates axes coincide with the major and minor axes, respectively. The $`y^{}`$ coordinate is the foreshortened $`y`$ coordinate.
In this section we will discuss the simplified case that all stellar orbits are circular, more realistic situations are discussed in § 4. The following elementary relations between the coordinates and the various projections of the orbital speed $`\overline{V_c}`$ are derived with the aid of Fig. Galaxy Distances via Rotational Parallaxes:
$`x`$ $`=`$ $`R\mathrm{cos}\theta `$ (1)
$`y`$ $`=`$ $`R\mathrm{sin}\theta `$ (2)
$`y^{}`$ $`=`$ $`{\displaystyle \frac{y}{\mathrm{cos}i}}={\displaystyle \frac{R\mathrm{sin}\theta }{\mathrm{cos}i}}`$ (3)
$`\mathrm{tan}\theta `$ $`=`$ $`{\displaystyle \frac{y}{x}}={\displaystyle \frac{y^{}}{x\mathrm{cos}i}}`$ (4)
$`V_x`$ $`=`$ $`s_\mathrm{\Omega }V_c\mathrm{sin}\theta ^{}`$ (5)
$`V_y`$ $`=`$ $`s_\mathrm{\Omega }V_c\mathrm{cos}\theta ^{}={\displaystyle \frac{V_x}{\mathrm{tan}\theta ^{}}}`$ (6)
$`V_r`$ $`=`$ $`V_y\mathrm{sin}i`$ (7)
$`V_y^{}`$ $`=`$ $`V_y\mathrm{cos}i=\kappa \mu _y^{}D`$ (8)
$`V_x`$ $`=`$ $`=\kappa \mu _xD`$ (9)
$`\mathrm{cos}\theta ^{}`$ $`=`$ $`{\displaystyle \frac{V_y}{\sqrt{V_x^2+V_y^2}}}={\displaystyle \frac{\mu _y^{}}{\sqrt{\mu _y^{}^2+\mu _x^2\mathrm{cos}^2i}}}`$ (10)
where $`V_x`$ and $`V_y`$ are the projections of $`\overline{V_c}`$ on the $`x`$ and $`y`$ axes. The angle $`\theta ^{}`$ between $`\overline{V}_c`$ and $`V_y`$ equals $`\theta \mathrm{arctan}(y/x)`$ for circular orbits. $`V_r`$ is the radial velocity along the line of sight and $`V_x`$ and $`V_y^{}`$ are the components of $`\overline{V_c}`$ along the apparent major and minor axes. The requirement that velocities are positive in the $`+x`$ and $`+y`$ directions leads to $`s_\mathrm{\Omega }=+1`$ for counter-clockwise rotation, and $`s_\mathrm{\Omega }=1`$ for clockwise motions. The observable proper motions along $`x`$ and $`y^{}`$ are symbolized by $`\mu _x`$ and $`\mu _y^{}`$, respectively.
### 2.1 The Principal Axes Method
The derivation of distances from observed proper motions and radial velocities is a well established practice in astronomy (e.g. orbital parallax), and are among the most reliable distance measures available. Since the proper motion is the ratio of a velocity (in $`\text{km\hspace{0.17em}s}^1`$) and the distance, the latter can be determined from observations. A common problem is that the angle between the space velocity and the line-of-sight is typically not known, so that the distance is determined modulo $`\mathrm{tan}i`$. For external galaxies, the inclination is well established from the axis ratio of the image and/or from the analysis of H I or H $`\alpha `$ radial velocity fields.
The principal axis method is particularly appealing for the case of circular orbits<sup>3</sup><sup>3</sup>3We do not consider out-of-the plane components here.. For elliptical orbits we have to introduce two additional a priori unknown angles: the angle $`\mathrm{\Delta }\theta _M`$ between the orbital velocity ($`\overline{V}_o`$) and the tangent to a circular orbit at the major axis, and the angle between a circular and an elliptical orbit at the minor axis, $`\mathrm{\Delta }\theta _m`$. These two angular differences are related to the eccentricity ($`e`$) and position angle ($`\varphi `$) of the elliptical orbit. Before proceeding in the derivation of the distance, it is profitable to re-write equations (7)-(9) for stars located on the principal axes:
$$\begin{array}{ccccccc}\hfill V_{r,M}& =& \hfill V_r(\mathrm{cos}\theta =\pm 1)& =& \hfill V_y\mathrm{sin}i& =& s_\mathrm{\Omega }V_o\mathrm{cos}\mathrm{\Delta }\theta _M\mathrm{sin}i\hfill \\ \hfill \mu _M& =& \hfill \mu _y^{}(\mathrm{cos}\theta =\pm 1)& =& \hfill \frac{V_y^{}}{\kappa D}& =& s_\mathrm{\Omega }\frac{V_o\mathrm{cos}\mathrm{\Delta }\theta _M\mathrm{cos}i}{\kappa D}\hfill \\ \hfill \mathrm{tan}\mathrm{\Delta }_{\theta _M}& =& \hfill \frac{\mu _x(\mathrm{cos}\theta =\pm 1)}{\mu _y^{}(\mathrm{cos}\theta =\pm 1)}& & & & \\ \hfill \mu _m& =& \hfill \mu _x(\mathrm{sin}(\theta )=\pm 1)& =& \hfill \frac{V_x}{\kappa D}& =& s_\mathrm{\Omega }\frac{V_o\mathrm{cos}\mathrm{\Delta }\theta _m}{\kappa D},\hfill \end{array}$$
(11)
where $`V_{r,M}`$ and $`\mu _M`$ are the radial and proper motion on the major axis and $`\mu _m`$ equals the proper motion on the minor axis. With an assumed inclination the first two equations of (11) can be solved for the “major-axis” distance $`D_M`$ \[eqn. (12) below\]. Further simplifications can be made in case the orbits are circular, so that the $`\mathrm{cos}\mathrm{\Delta }\theta _M`$ and $`\mathrm{cos}\mathrm{\Delta }\theta _m`$ terms in (11) equal unity. $`D_M`$ and the “minor-major axes” inclination and distance are given by:
$`D_M`$ $`=`$ $`{\displaystyle \frac{V_c\mathrm{cos}i}{\kappa \mu _M}}={\displaystyle \frac{V_{r,M}}{\kappa \mathrm{tan}i\mu _M}}`$ (12)
$`\mathrm{sin}i_{mM}`$ $`=`$ $`\sqrt{1{\displaystyle \frac{\mu _M}{\mu _m}}}`$ (13)
$`D_{mM}`$ $`=`$ $`{\displaystyle \frac{V_{r,M}}{\kappa \sqrt{\mu _m^2\mu _M^2}}}.`$ (14)
A clear advantage of the $`mM`$ method is that the systemic motion of the galaxy as a whole can be taken out easily by considering points symmetric with respect to the center of the galaxy. Rotation-induced proper motions at symmetric points have opposite sign, while the radial and planar components of the systemic velocity ($`V_{sys}`$) have the same orientation and magnitude. The rotation-induced proper motion component can thus be easily computed. For example, when taking out the systemic motion, $`\mu _m,\mu _M`$ and $`V_{r,M}`$ in equations (12)-(14) should be replaced by $`\mu _m=\frac{1}{2}[\mu _m^{}(\theta =90)\mu _m^{}(\theta =+90)]`$ and similar relations for $`V_r`$ and $`\mu _M`$. The primed relations indicate the observed motions that include the projection of the systemic motion.
One potential disadvantage of this principal axes or $`nM`$ method \[cf. eqn. (14)\] is that the measurements of $`\mu _m`$ and $`\mu _M`$ are to be taken at the same galactocentric radius, or from radii with the same circular velocity and inclination. In the discussion above we did not consider the possibility that the actual distances to the target stars may differ significantly, which may be the case for the nearest galaxies (up to 5% for M 31 & M 33). However, application of the major-axis method (12) is unaffected by this problem. We present a more general distance solution in later sections.
### 2.2 The Single Star Method
Another potential disadvantage of the $`mM`$ method is that the available surface area in the galaxy for suitable stars is relatively small, since the target stars are limited to regions close to the principal axes and because suitably bright targets are quite rare. Fortunately, other rotational parallax varieties exist, and we discuss their merits below.
In fact, we can determine the distance of a single stars if we neglect, for now, the systemic motion of the galaxy and the random motions of the stars. Using some of the relations (1)-(9), we solve for the inclination of the stellar orbit:
$`\mathrm{cos}i`$ $`=`$ $`{\displaystyle \frac{\mu _y^{}}{\mu _x}}\mathrm{tan}\theta ^{}`$ (15)
$`\mathrm{cos}^2i`$ $`=`$ $`{\displaystyle \frac{\mu _y^{}}{\mu _x}}{\displaystyle \frac{y^{}}{x}},`$ (16)
Equation (15) holds for any type of orbit, circular or elliptical. Further, if the orbit is circular we can equate $`\theta `$ and $`\theta ^{}`$ and arrive at equation (16). The orbital speed can be computed from the radial velocity \[eqn. (7)\]: $`V_o=\frac{V_r}{\mathrm{cos}\theta ^{}\mathrm{sin}i}`$. This equation can be re-written for the elliptical and circular orbit cases with the aid of eqns. (10) and (16), respectively:
$`V_o`$ $`=`$ $`V_r\sqrt{1+{\displaystyle \frac{1}{\mathrm{tan}^2i}}}\left(1+{\displaystyle \frac{\mu _x^2}{\mu _y^{}^2}}\right)`$ (17)
$`V_c`$ $`=`$ $`V_r\sqrt{{\displaystyle \frac{\mu _x}{\mu _y^{}}}{\displaystyle \frac{x\mu _y^{}y^{}\mu _x}{x\mu _x+y^{}\mu _y^{}}}}`$ (18)
In the last equation we expressed the trigonometric terms in the primary observables.
A rotational parallax distance can now be obtained by combining the radial velocity and either the proper motion perpendicular or parallel to the major axis. As was the case for the principal axis method, we need to make a distinction between the case of elliptical and circular orbits. If the orbit is elliptical, the distance must be calculated from $`V_r`$ and $`\mu _y^{}`$, and depends only on the assumed inclination. In the circular orbit case the expressions obtained from the $`\mu _x`$ and $`\mu _y^{}`$ are identical and can be re-written in terms of observable quantities only. For the two single-star distances we find:
$`D_y^{}`$ $`=`$ $`{\displaystyle \frac{V_r}{\kappa \mu _y^{}\mathrm{tan}i}}`$ (19)
$`D_{PP}`$ $`=`$ $`{\displaystyle \frac{V_r}{\kappa }}\sqrt{{\displaystyle \frac{y^{}/\mu _y^{}}{x\mu _x+y^{}\mu _y^{}}}}`$ (20)
We will refer to distances calculated from the perpendicular and parallel proper motions of individual stars \[eqn.(20)\] as “$`PP`$ distances.”
### 2.3 Error Estimates
It is instructive to estimate the attainable errors using SIM astrometry. This is most easily done for the case of circular orbits. Below we present the expressions for the errors in the inferred distance, inclination and circular velocity, where we express all trigonometric terms in the primary observables:
$`{\displaystyle \frac{\mathrm{\Delta }D_y^{}}{D_y^{}}}`$ $`=`$ $`\sqrt{\left({\displaystyle \frac{\mathrm{\Delta }V_r}{V_r}}\right)^2+\left({\displaystyle \frac{\mathrm{\Delta }\mu _y^{}}{\mu _y^{}}}\right)^2}`$ (21)
$`{\displaystyle \frac{\mathrm{\Delta }D_{mM}}{D_{mM}}}`$ $`=`$ $`\sqrt{\left({\displaystyle \frac{\mathrm{\Delta }V_r}{V_r}}\right)^2+\left({\displaystyle \frac{\mu _m\mathrm{\Delta }\mu _m}{\mu _m^2\mu _M^2}}\right)^2+\left({\displaystyle \frac{\mu _M\mathrm{\Delta }\mu _M}{\mu _m^2\mu _M^2}}\right)^2}`$ (22)
$`{\displaystyle \frac{\mathrm{\Delta }\mathrm{cos}{}_{}{}^{2}i}{\mathrm{cos}{}_{}{}^{2}i}}`$ $`=`$ $`\sqrt{\left({\displaystyle \frac{\mathrm{\Delta }\mu _y^{}}{\mu _y^{}}}\right)^2+\left({\displaystyle \frac{\mathrm{\Delta }\mu _x}{\mu _x}}\right)^2}`$ (23)
$`\left({\displaystyle \frac{\mathrm{\Delta }V_c}{V_c}}\right)^2`$ $`=`$ $`\left({\displaystyle \frac{\mathrm{\Delta }V_r}{V_r}}\right)^2+\left({\displaystyle \frac{y^{}}{2}}{\displaystyle \frac{x(\mu _x^2\mu _y^{}^2)+2y^{}\mu _x\mu _y^{}}{(x\mu _x+y^{}\mu _y^{})(x\mu _y^{}y^{}\mu _x)}}\right)^2\times `$ (24)
$`\left(\left({\displaystyle \frac{\mathrm{\Delta }\mu _x}{\mu _x}}\right)^2+\left({\displaystyle \frac{\mathrm{\Delta }\mu _y^{}}{\mu _y^{}}}\right)^2\right)`$
$`\left({\displaystyle \frac{\mathrm{\Delta }D_{PP}}{D_{PP}}}\right)^2`$ $`=`$ $`\left({\displaystyle \frac{\mathrm{\Delta }V_r}{V_r}}\right)^2+{\displaystyle \frac{1}{4(x\mu _x+y^{}\mu _y^{})^2}}\times `$ (25)
$`\left((x\mu _x)^2\left({\displaystyle \frac{\mathrm{\Delta }\mu _x}{\mu _x}}\right)^2+(x\mu _x+2y^{}\mu _y^{})^2\left({\displaystyle \frac{\mathrm{\Delta }\mu _y^{}}{\mu _y^{}}}\right)^2\right)`$
At large inclination, where $`x\mu _x>>y^{}\mu _y^{}`$, the errors on $`V_c`$ and $`D_{PP}`$ simplify to
$`\left(\frac{\mathrm{\Delta }V_r}{V_r}\right)^2+\frac{1}{4}\left(\left(\frac{\mathrm{\Delta }\mu _x}{\mu _x}\right)^2+\left(\frac{\mathrm{\Delta }\mu _y^{}}{\mu _y^{}}\right)^2\right)`$.
Typical values for proper motions may be $`\mu _{x,y^{}}16`$ $`\mu \text{as\hspace{0.17em}yr}^1`$ (40 for M31), $`\mathrm{\Delta }\mu _{x,y^{}}4`$ $`\mu \text{as\hspace{0.17em}yr}^1`$, while the radial velocity uncertainty would be less than 10 $`\text{km\hspace{0.17em}s}^1`$. For M31, the fractional errors in the radial velocity and the proper motions equal about, 4%, $`25`$% and $`10`$%, for $`V_r`$, $`\mu _y^{}`$, and $`\mu _x`$, respectively. The total fractional errors on $`D_{PP}`$$`V_c`$ and $`D_y^{}`$ are of order 15, 15, and 25%, respectively. These values compare well with the exact values presented in Fig. Galaxy Distances via Rotational Parallaxes.
These considerations indicate that proper motion errors dominate the final uncertainty in the derived distances. The $`mM`$ distance can have substantially smaller errors. Also note that, in the pure circular motion case, the $`mM`$ distance estimate requires astrometry for two stars, while the $`PP`$ distances can be obtained for one star only<sup>4</sup><sup>4</sup>4When comparing the final errors on the distances, we multiply the $`mM`$ errors by $`\sqrt{2}`$ so as to compare the distance errors per target (i.e. per unit integration time). . We compare the distance error budgets in figures (Galaxy Distances via Rotational Parallaxes)-(Galaxy Distances via Rotational Parallaxes), for M 31, M 33 and M 81 for a range of plausible astrometric errors.
These figures illustrate an important aspect of the single-star method: the attainable distance and inclination errors are not very sensitive to the galactic azimuth $`\theta `$. As a result, targets from all over the galaxy can be incorporated in the analysis. Furthermore, because of its proximity and large rotation speed, M 31 clearly offers the best opportunity for reliable rotational parallax measurements, followed by M 33 and M 81.
The same figures indicate that the $`PP`$ method works best at low inclinations and that the $`mM`$ method is superior at large inclinations. The latter fact can be understood as follows: for a strongly inclined galaxy, the proper motion perpendicular to the major axis is small ($`\mu _M=\mu _y^{}0`$), so that the $`\mu _M`$ term in equation (14) for $`D_{mM}`$ goes to zero. In that case $`D_{mM}V_{r,M}/(\kappa \mu _m`$). As a result, the third error-term in equation (22) becomes negligible. For the $`PP`$ method however, both proper motions are required to derive the distance, so that the proper motion error comes in twice.
In figure Galaxy Distances via Rotational Parallaxes we present the dependence of the distance errors on both radial velocity and proper motion errors for M 31. In that figure we contour, from top to bottom, the fractional errors on $`D_y^{}`$, $`D_{PP}`$ and $`D_{mM}`$. For the Local Group galaxies it will pay off to decrease the radial velocity errors to as low values as possible. As expected, proper motion errors dominate the error budget for distant galaxies such as M 81 ($`D3.6`$ Mpc, not shown).
Naively, one might expect that of our three distance estimators, $`D_y^{}`$ would be determined most accurately. After all, for this calculation we assume an inclination and we use only one proper potion, while $`D_{mM}`$ and $`D_{PP}`$ require an additional proper motion measurement. Figure (Galaxy Distances via Rotational Parallaxes) shows that (for $`\theta =30\mathrm{°}`$) this is not always the case. In fact, the actual expected errors are complex functions of location in the galaxy, inclination, rotation speed and distance. For distant galaxies, the errors are dominated by proper motion errors, so that all three methods are equivalent (not shown). For M 33 all methods yield approximately equally well determined distances, while distance errors for M 31 are minimized if we can use the principal axes method.
For $`\mathrm{\Delta }\mu _i1`$ $`\mu \text{as\hspace{0.17em}yr}^1`$ and $`\mathrm{\Delta }V_r10`$ $`\text{km\hspace{0.17em}s}^1`$, the errors on the inferred distance equal about 6, 15, and 10% for M 31, M 33 and M 81, respectively. For the Local Groups galaxies the random motions of stars introduce additional uncertainties. Assuming a dispersion of 10 $`\text{km\hspace{0.17em}s}^1`$, the effective proper motion accuracies are about 2.5 $`\mu \text{as\hspace{0.17em}yr}^1`$ at 0.8 Mpc. The internal dispersion at the distance of M 81 is about five times smaller and insignificant with respect to the likely measurement errors that SIM can attain in “reasonable” integration times.
## 3 Deviations from the Gedanken Galaxy
### 3.1 Systematic Motions
In section 2.1 we have seen that systematic motions induced by the space motion of the galaxy ($`\overline{V}_{sys}`$) can be corrected for rather easily in the principal axes method. Likewise, a global fit to the observed proper motions and radial velocities would allow for the determination of $`\overline{V}_{sys}`$ in the $`PP`$ method. For example, the radial velocities of M 31, M 33 and M 81 equal -300 $`\pm `$4, -200 $`\pm `$6 and -34 $`\pm `$4 $`\text{km\hspace{0.17em}s}^1`$. These velocities are comparable to their internal velocities of $`270`$, $`100`$ and $`213`$ $`\text{km\hspace{0.17em}s}^1`$. We will discuss this issue in later subsections.
### 3.2 The Effects of Stellar Warps
It is known from radial velocity observations that, for most galaxies, the inclination gradually changes as a function of radius. This so-called “warping” is most pronounced beyond the optical disk (Briggs, 1990). Although such inclination changes are hard to detect inside the optical disk, some galaxies are known to warp in the the outskirts of the optical disk<sup>5</sup><sup>5</sup>5e.g., The Milky Way (Drimmel, Smart & Lattanzi, 2000; Porcel, Battaner & Jimenez-Vicente, 1997), and references therein, NGC 7814 (Lequeux, Dantel-Fort & Fort, 1995), and others (Sanchez-Saavedra, Battaner & Florido, 1990; Florido et al., 1991), see Reshetnikov & Combes (1998) for a recent review. These authors measured the warps at a distance ($`R_{25}`$) where the B-band surface brightness is roughly 25 magnitudes per square arcseconds. This surface brightness is reached at radii of approximately 90 and 31 arc-minutes in M 31 and M 33, respectively (de Vaucouleurs et al., 1991).
Cursory inspection of large-scale optical pictures of the Andromeda galaxy suggests that the plane of the galaxy is warped in the outer parts of the galaxy. In fact, the orientation of the outer isophotes starts changing –an indication of a warp– beyond about 90 arc-minutes from the center (Walterbos & Kennicutt, 1988).
Also, an analysis of the spiral structure and H I radial velocity field of M 31 suggests that the inclination is about 60° at 10′ from the center. Due to the almost linear warping of $`+0\mathrm{°}.3`$ per kpc, an inclination of 80° is reached at a radius of 100′ (Braun, 1991). Braun’s determination of the radial inclination gradient is based on the association of H I with spiral arms, but does not self-consistently include the non-axisymmetric velocity component that is likely to be induced by the spiral density wave. It thus may be that the actual warp is less severe than claimed by Braun.
Similarly, an analysis of the optical spiral arms in M 33 suggests that the inclination of M 33 changes from forty to sixty-three degrees across the disk (Sandage & Humphreys, 1980), although this result is disputed (Maucherat et al., 1984). On the other hand, the H I warping of M 33 Corbelli & Schneider (1997) sets in around $`R_{25}`$ (i.e., 30 arcmin, or 8.2 kpc at 0.84 Mpc).
### 3.3 Spiral Structure
Non-circular orbits, or more generally, streaming motions with respect to the simplified picture of the circular orbit might significantly complicate the determination of rotational parallaxes. Streaming motions induced by spiral density waves can reach tens of kilometers per second, depending on the galaxy, and the locations of the target with respect to the spiral arms. To first order, such deviations would induce systematic errors in the inferred distances that equal the systematic velocities induced by spiral density waves, or about $`\mathrm{\Delta }V/V_c20`$%.
The theory of spiral arm density waves has been developed extensively \[e.g. (Lin, Yuan & Shu, 1969)\]. Applications of this theory to radial velocity fields of external galaxies show that streaming motions of the order of 20-50 $`\text{km\hspace{0.17em}s}^1`$ \[e.g. (Visser, 1980; Vogel, Kulkarni & Scoville, 1988; Boulanger & Viallefond, 1992; Tilanus & Allen, 1993)\]. The analysis of the spiral structure of M 31 by Braun (1991) indicates streaming motions of 40 $`\text{km\hspace{0.17em}s}^1`$ for radii smaller than 40$`\mathrm{}`$, of 20 $`\text{km\hspace{0.17em}s}^1`$ between 30$`\mathrm{}`$ and 75$`\mathrm{}`$ and of order 10 $`\text{km\hspace{0.17em}s}^1`$ in the outermost parts. M 31 exhibits a well developed two-armed spiral pattern with a small, but radially varying pitch angle Braun (1991). In the inner region, $`R27\mathrm{}`$, the spiral has a larger pitch angle ($`\varphi 16\mathrm{°}`$) than in the outer region ($`\varphi 7\mathrm{°}`$).
Stars will respond to the potential induced by the spiral density wave. If the pattern is logarithmic, the radial and tangential velocities \[$`\stackrel{~}{V}_R(R)`$ and $`\stackrel{~}{V}_\theta (R)`$\] induced by the perturbation can be found analytically. These perturbation velocities vary with (extra) galactocentric radius $`R`$. So do the pitch angle and the spiral phase ($`\chi `$). $`\stackrel{~}{V}_R`$, $`\stackrel{~}{V}_\theta `$, $`\varphi `$ and $`\chi `$ can be calculated from the rotation curve, the stellar velocity dispersion, and the pattern-speed and amplitude of the perturbation, where the last three properties are not well established observationally for our target galaxies. This procedure has been followed for the Milky Way (Crézé & Mennessier, 1973; Amaral & Lepine, 1997; Mishurov et al., 1997; Mishurov & Zenina, 1999; Lepine, Mishurov & Dedikov, 2000), and can be readily generalized to external galaxies.
It is also possible that the spiral structure is generated in response to a bar or nearby companions. For the Andromeda galaxy, the companions, M 32 and NGC 205, have estimated peri-centers between 13 and 35 kpc \[e.g. (Byrd, 1976,1977,1978; Sato & Sawa, 1986; Cepa & Beckman, 1988)\]. In fact, radial velocities and SIM-based proper motions can conceivably be used to determine the parameters of the spiral pattern, and allow for detailed tests of various theories of spiral structure.
However, these spiral-structure theories have a large number of parameters. In the density-wave theory, almost all factors that determine the spiral pattern depend on galactocentric radius, so that these “factors” are really functions with many more unknowns. Further, on small scales, the details of spiral structure may deviate from the large-scale $`n`$-armed spiral. For example, it has been suggested that the large star-formation complex NGC 206 in Andromeda resulted from the recent interaction between two spiral-arm segments, moving with relative velocity of about 30 $`\text{km\hspace{0.17em}s}^1`$ (Magnier et al., 1997). Alternatively, the often seen bifurcations in the spiral patterns of spiral galaxies may arise from the superposition of spiral modes with different multiplicity. The Milky Way may be an example of two-plus-four armed spiral galaxy (Amaral & Lepine, 1997).
Clearly, we would like to avoid the complications that arise due to streaming motions on small and large scales. We will discuss the effects of non-axisymmetric streaming motions in more detail in section 5.
## 4 Accurate Distances, Notwithstanding Perturbations?
In the previous sections we discussed the case of circular orbits and alluded occasionally to the case of elliptical orbits. A complication of elliptical orbits is that the angle between the orbital velocity and the line-of-sight ($`\theta ^{}`$) can not be deduced from the position of the target and the inclination of the galaxy. That is to say, $`\theta ^{}\theta `$ \[c.f. eqns. (4) and (10)\]. This means that predicting the orbital motion (velocity and direction) at point ($`x_2,y_2`$) given measurements at ($`x_1,y_1`$) becomes significantly more complicated, even if these points lie at the same galactocentric radius. Knowledge of the ellipticity and position angle of the orbit are required to solve this problem. Furthermore, due to the intrinsic dispersions of the target population, additional velocity vectors are added to the motions of the targets. All these parameters are likely to depend on galactocentric distance. And finally, the systemic motion of the galaxy needs to be determined and its effects subtracted from on the individual radial velocities and proper motions.
### 4.1 $`\mu _y^{}`$-$`V_r`$ Correlations
Here we extend on the procedure to determine the $`D_y^{}`$ distance \[§§ 2.2, eqn. (19)\] and make use of the fact that the radial velocity and the $`V_y^{}`$ velocity are two orthogonal components of the total space velocity in the $`ry^{}`$ plane ($`V_{tot,ry^{}}`$, see figure Galaxy Distances via Rotational Parallaxes). Like in equation (19), the distance then follows from the observed radial velocity and $`y^{}`$ proper motion: $`D_y^{}=V_r/(\kappa \mu _y^{}\mathrm{tan}i_t)`$, where $`i_t`$ is the angle between $`V_{tot,ry^{}}`$ and the plane of the sky. Generally speaking, $`i_t`$ differs from the geometric inclination of the galaxy. For example, $`i_t`$ will lie between $`i`$ and 90° if random motions are unimportant and the systemic motion is due to Hubble flow only. We illustrate the various contributions to $`V_{tot,ry^{}}`$ in figure Galaxy Distances via Rotational Parallaxes. From this figure we deduce the following relations for the observed $`y^{}`$ proper motion and the radial velocity:
$`V_y^{}`$ $`=`$ $`\left(V_{o,y}+V_{\sigma ,y}\right)\mathrm{cos}i+V_{sys,ry^{}}\mathrm{cos}i_s+V_{\sigma ,z}\mathrm{sin}i`$ (26)
$`V_{obs,r}`$ $`=`$ $`\left(V_{c,y}+V_{e,y}+V_{\sigma ,y}\right)\mathrm{sin}i+V_{sys,ry^{}}\mathrm{sin}i_sV_{\sigma ,z}\mathrm{cos}i`$ (27)
where $`i_s`$ represents the angle between the systemic motion and the plane of the sky. The orbital velocity in the $`y`$-direction ($`V_{o,y}`$) comprises both a circular orbit term ($`V_{c,y}`$) and a contribution for the ellipticity of the orbit ($`V_{e,y}`$). The random motions in the $`y`$ and $`z`$ directions are denoted as $`V_{\sigma ,y}`$ and $`V_{\sigma ,z}`$, respectively. We now proceed by expressing the observable $`y^{}`$ proper motion in terms of the observable radial velocity. To this end we solve eqn. (27) for $`(V_{o,y}+V_{\sigma ,y})`$, substitute the result in (26) and divide by $`\kappa D`$ to arrive at:
$`\mu _y^{}`$ $`=`$ $`{\displaystyle \frac{V_{obs,r}}{\kappa D\mathrm{tan}i}}+{\displaystyle \frac{1}{\kappa D}}\left(V_{sys,ry^{}}(\mathrm{cos}i_s{\displaystyle \frac{\mathrm{sin}i_s}{\mathrm{tan}i}})+{\displaystyle \frac{V_{\sigma ,z}}{\mathrm{sin}i}}\right)`$ (28)
$`=`$ $`\alpha _y^{}V_{obs,r}+\beta _y^{}.`$ (29)
These two equations show that the observable $`y^{}`$ proper motion is a linear function of the observable radial velocity. We recognize the $`\mathrm{cos}i_s`$ and $`\mathrm{sin}i_s`$ terms in (28) as the $`y`$ and radial velocity components of the systemic velocity, respectively. More importantly, equation (28) shows that the slope $`\alpha _y^{}`$ is independent of the presence of non-circular motions. This is so because in eqns. (27)-(29) we use the total $`y`$ component of the orbital speed: for the purpose of distance determination, it is not necessary to know how the total $`y`$ velocity is distributed between the circular, elliptical and random components.
### 4.2 $`\mu _x`$ Correlations
Following the method of the previous subsection, we will express the observed $`\mu _x`$ as a function of either the observed radial velocity or the $`y^{}`$ proper motion. In order to do so, we need to connect the $`x`$ and $`y`$ components of the orbital speed via a simple relation. This is straightforward for the circular component of the orbits. As above, the total observable velocity comprises circular velocities ($`V_{c,x}`$ and $`V_{c,y}`$), “elliptical velocities” ($`V_{e,x}`$ and $`V_{e,y}`$), random motion components and the systemic velocity. We also employ the $`r`$ and $`y^{}`$ motions derived above \[eqns. (26) and (27)\] and the $`x`$ motion:
$`V_{tot,x}`$ $`=`$ $`V_{c,x}+V_{e,x}+V_{sys,x}+V_{\sigma ,x}`$ (30)
First we solve eqn. (27) for $`V_{c,x}`$, substitute the result in (30) and divide by $`\kappa D`$ to find $`\mu _x`$:
$`\mu _{x,r}`$ $`=`$ $`{\displaystyle \frac{\mathrm{tan}\theta }{\mathrm{sin}i\kappa D}}\mathrm{\Delta }V_{obs,r}+{\displaystyle \frac{\mathrm{tan}\theta }{\kappa D}}\left(V_{e,y}+V_{\sigma ,y}{\displaystyle \frac{V_{\sigma ,z}}{\mathrm{tan}i}}\right)+{\displaystyle \frac{V_{e,x}+V_{sys,x}+V_{\sigma ,x}}{\kappa D}},`$ (31)
where $`\mathrm{\Delta }V_{obs,r}V_{obs,r}V_{sys,ry^{}}\mathrm{sin}i_s`$. Likewise, when solving the $`y^{}`$ proper motion for $`V_{c,x}`$ and using $`\mathrm{\Delta }\mu _y^{}\mu _y^{}V_{sys,ry^{}}\mathrm{cos}i_s`$, we find:
$`\mu _{x,y^{}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{tan}\theta }{\mathrm{cos}i}}\mathrm{\Delta }\mu _y^{}+{\displaystyle \frac{\mathrm{tan}\theta }{\kappa D}}\left(V_{e,y}+V_{\sigma ,y}+V_{\sigma ,z}\mathrm{tan}i\right)+{\displaystyle \frac{V_{e,x}+V_{sys,x}+V_{\sigma ,x}}{\kappa D}}.`$ (32)
In case the orbits are exactly circular, equations (31) and (32) simplify to:
$`\mu _{x,r}`$ $`=`$ $`\mu _x={\displaystyle \frac{y^{}/x\mathrm{\Delta }V_{obs,r}}{\mathrm{sin}i\mathrm{cos}i\kappa D}}+{\displaystyle \frac{V_{sys,x}}{\kappa D}}=\alpha _{x,r}(y^{}/x\mathrm{\Delta }V_{obs,r})+\beta _x,`$ (33)
$`\mu _{x,y^{}}`$ $`=`$ $`\mu _x={\displaystyle \frac{y^{}/x\mathrm{\Delta }\mu _y^{}}{\mathrm{cos}^2i}}+{\displaystyle \frac{V_{sys,x}}{\kappa D}}=\alpha _{x,yp}(y^{}/x\mathrm{\Delta }\mu _y^{})+\beta _{x,y^{}},`$ (34)
where we have replaced $`\mathrm{tan}\theta `$ by $`y^{}/(x\mathrm{cos}i)`$. The independent variables $`(y^{}/x\mathrm{\Delta }V_{obs,r})`$ and $`(y^{}/x\mathrm{\Delta }\mu _y^{})`$ are defined for $`x=0`$ since, for circular orbits, they are proportional to $`(V_c\mathrm{sin}\theta )`$ \[cf. eqns. (4)-(9)\]. However, the random and elliptical contributions to $`\mathrm{\Delta }V_{obs,r}`$ can “blow-up” close to the minor axis as a result of the multiplication by $`1/x`$. It is thus advisable to down-weight the regions close to the minor axis. Like the $`\alpha _y^{}`$ slope derived in the previous subsection, the slope $`\alpha _{x,r}`$ depends on both the inclination and the distance, but in an independent manner. The correlation between $`\mu _x`$ and $`\mu _y^{}`$ yields the inclination of the galaxy directly.
The proper motion equations (29), (33), and (34) represent the multi-star equivalent of the single-stars case and include the systemic motion of the galaxy: we solve them for inclination and distance in §4.3 below, and we discuss the effects of non-circular and random motions in section 5.
### 4.3 The Rotational Parallax Distance
In the previous subsections we found that the $`y^{}`$ and $`x`$ proper motions can be expressed in terms of the observed radial velocities. From these linear slopes the inclination and distance follow:
$`\alpha _y^{}`$ $`=`$ $`{\displaystyle \frac{1}{\kappa D\mathrm{tan}i}}`$ (35)
$`\alpha _{x,r}`$ $`=`$ $`{\displaystyle \frac{1}{\kappa D\mathrm{tan}i\mathrm{cos}^2i}}`$ (36)
$`\mathrm{cos}^2i`$ $`=`$ $`{\displaystyle \frac{\alpha _y^{}}{\alpha _{x,r}}}={\displaystyle \frac{1}{\alpha _{x,yp}}}`$ (37)
$`D_{\alpha i,xy^{}}`$ $`=`$ $`{\displaystyle \frac{1}{\kappa \alpha _y^{}\mathrm{tan}i}}{\displaystyle \frac{1}{\kappa \alpha _{xr}\mathrm{tan}i\mathrm{cos}^2i}}`$ (38)
$`D_{\alpha xy^{}r}`$ $`=`$ $`{\displaystyle \frac{1}{\kappa \alpha _y^{}}}\sqrt{{\displaystyle \frac{\alpha _y^{}}{\alpha _{x,r}+\alpha _y^{}}}}.`$ (39)
The above relations hold exactly for circular orbits. We will investigate the effects of non-circular orbits in the next subsection. Equations (37) and (38) for $`i`$ and $`D_{\alpha i,xy^{}}`$ are the equivalent of the relations for inclination and distance as derived in the single-star method \[(16) and (19)\]. It thus follows that $`\alpha _{x,r}`$ is equivalent to $`x\mu _x`$, $`\alpha _y^{}`$ corresponds to $`y^{}\mu _y^{}`$ and $`\alpha _y^{}\mathrm{\Delta }V_{obs,r}\mu _y^{}`$, so that eqn. (39) is obtained.
### 4.4 The Systemic Velocity Vector
A straightforward manner to determine the systemic velocity of the galaxy would be to take the appropriate averages of the observed proper motions and radial velocities of targets that have similar $`|x|`$ and $`|y|`$ coordinates. In this way, the motions induced by the internal motions of the galaxy are averaged out, so that the components of the systemic velocity of galaxy become apparent. However, due to elliptical streaming motions these values can only be considered to be “reasonable” values. A better determination of the systemic velocity is possible from a full solution of equations (28), (31) and (32).
## 5 Practical Implementation
In section 3 we have reviewed the various effects that complicate the implementation of the single-star rotational parallax method. Above we outlined the modifications required to measure the systemic motion. In the next subsections we show that neither warping of the stellar disk nor spiral structure inhibits our ability to determine accurate rotational parallaxes.
### 5.1 Inclination Effetcs
For many spiral galaxies, it is possible to determine the inclination from the axis-ratio of contours of constant optical surface brightness. For M 31 this procedure is more complicated due to its large inclination and finite thickness. However, the inclination can also be determined from H I velocity fields \[e.g., see Unwin (1983) and Brinks & Burton (1984) for M 31 and Corbelli & Schneider (1997) for M 33\]. These inclination estimates typically have an uncertainty of a few degrees. If such H I inclination estimates were to be used to determine the distance from the slope of the observed $`\mu _y^{}`$-$`V_r`$ relation, the resulting distance uncertainties would be of order 20% for M 81 and M 33, and a factor of two for M 31. To obtain distance estimates that are accurate at the percent level, the other proper motion component needs to be utilized.
From equations (29), (33), and (34) it is clear that no a-priori knowledge of the inclination is required, if the inclination does not vary significantly with radius. In that case, distance and inclination can be determined, even when the targets are arbitrarily located across the face of the galaxy, provided that the kinematic variations are sufficiently sampled. On the other hand, if the inclination does change with radius, it might be best to select targets in an elliptical annulus with an axis-ratio equal to the cosine of the best-guess inclination. The drawback of such an approach is that the range in radii sampled will increase if the a-priori inclination estimate was wrong<sup>6</sup><sup>6</sup>6By about a factor of two for a 2° error at $`i_{est}=75\mathrm{°}`$.. In principle, the best distance determination is possible when all targets have similar distances from the galaxy center so as to minimize the effects of any possible radial variations of inclination (and rotation speed).
### 5.2 Non-Circular Motions
In general, non-circular motions will be present in our target galaxies. The question is, how will those motions affect our ability to determine an accurate rotational parallax. Equations (31)-(34) above show that elliptical motions in both the $`x`$ and $`y`$ direction can adversely influence the distance determination. Below we will outline a technique that can be used to detect those motions and correct for their effects.
The contribution from elliptical motions in equations (27) and (30) can be written as the sum of two component. If the angular variation of an elliptical streaming component is identical to the azimuthal dependence of the circular velocity component we term that component “invisible.” That is to say, an invisible elliptical streaming field will induce a proper motion and radial velocity field that is indistinguishable from that of the circular streaming field. Thus, the invisible elliptical motion is given by:
$`\overline{V}_{ei}`$ $`=`$ $`V_{ei,x}\mathrm{sin}\theta \widehat{x}+V_{ei,y}\mathrm{cos}\theta \widehat{y},`$ (40)
with $`\widehat{x}`$ and $`\widehat{y}`$ the unit vectors in the $`x`$ and $`y`$ directions, respectively. Also note that, if $`V_{ei,x}`$ equals $`V_{ei,y}`$, then the resulting elliptical motion amounts to an additional circular velocity term, and should be absorbed in $`V_c`$. All elliptical streaming components orthogonal to $`\overline{V}_{ei}`$ are directly detectable in the observed stellar motions. The invisible elliptical streaming can significantly bias the inferred inclination and distances since their proper motion and radial velocity signature are indistinguishable from the circular motion terms. In section 5.4 we show how invisible streaming motions can be detected, and their distance-bias corrected for.
The exact functional form of the elliptical streaming field depends on the physical mechanism that drives such non-circular motions (§3.3). A full investigation of the dependence of the effect of non-circular streaming on the accuracy with which the rotational parallax can be determined is beyond the scope of the current work. However, we will explore the effects of non-circular motion by investigating the effects of a toy model for elliptical streaming. This illustrative model has the invisible streaming component discussed above plus an additional orthogonal, visible, component:
$`\overline{V}_{e,toy}`$ $`=`$ $`\overline{V}_{ei}+V_{ev,x}\mathrm{cos}\theta \widehat{x}+V_{ev,y}\mathrm{sin}\theta \widehat{y},`$ (41)
with $`V_{ev,x}`$ and $`V_{ev,y}`$ the visible $`x`$ and $`y`$ components of the elliptical streaming field, respectively.
### 5.3 The Works
The rotational parallax distances derived above can be improved upon if the $`\mu _y^{}`$-$`V_r`$ and $`\mu _x`$-$`V_r`$ equations \[(28) and (33)\] are solved simultaneously. In order to arrive at a unbiased solutions in case $`\overline{V}_{ei}`$ is non-zero, additional constraints are required. For example, one could demand that the inferred radial gradients of the inclination and rotation curve are small and linear. However, better, non-parametric additional constraints are available: 1) the lack of azimuthal variation of the rotation curve inferred from $`\mu _y^{}`$, $`V_r`$ and $`\mu _x`$, and 2) the requirement that the rotation speed derived from these three observable are identical, to within the errors<sup>7</sup><sup>7</sup>7Since $`V_r`$ is easier to measure than $`\mu _y^{}`$ (for M 31), and the resulting $`V_c(\theta )`$ should be equivalent, we will only use the information contained in $`V_r`$ and $`\mu _x`$.
Such a multiple non-linear regression solution as described above is beyond the scope of the current paper. However, we will investigate a poor-man’s approach to the problem and use that to derive estimates of the accuracy to which the rotational parallax can be determined with an instrument like SIM. Although the multiple regression approach is clearly advisable to make optimal use of the available data, the poor-man’s approach has the advantage that it clearly illustrates the problematic areas of the rotational parallax determination method. Further, since the poor-man’s route will not lead to the best possible solution, the error estimates so obtained are likely to be improved upon when using a multiple non-linear regression technique.
In the poor-man’s approach we split the procedure in four distinct steps. In the first step, the $`\mu _y^{}`$-$`V_r`$ correlation \[eqn. (28)\] is used to determine a solution for the product of distance and the the tangent of the inclination: $`D\times \mathrm{tan}i=(\kappa \alpha _y^{})^1`$ \[cf. eqn. (38)\]. As discussed in section §4.1, this determination of $`D\times \mathrm{tan}i`$ has no sensitivity to elliptical streaming motions at all. In the second step we determine $`D,i`$ and $`\overline{V}_{ev}`$ from the $`\mu _x`$-$`V_r`$ correlation \[eqn. (33)\] given the previously determined value for $`D\times \mathrm{tan}i`$. In the third step, we repeat the previous step $`N_{try}`$ times. We randomly selected a value for $`D\times \mathrm{tan}i`$, based on its estimated value and dispersion obtained in step #1. The final best value for the to-be-fitted parameters is obtained by averaging the $`N_{try}`$ results, where the derived errors equal the second moment of the $`N_{try}`$ values. In the fourth and final step, we check the additional constraints that the inferred rotation curve should have negligible azimuthal dependence, and that the rotation speed value inferred from the radial velocities and proper motions are equal to within the errors (see the Diagnostics section below for details). In principle, the diagnostics step can already be incorporated in step #2 to improve the estimates on the parameters. We will not do so in the spirit of deriving upper limits to the error budgets.
#### 5.3.1 Rotation Speed Diagnostics
Starting from equations (30) and (27) we write<sup>8</sup><sup>8</sup>8We set the contributions from the systemic motion and the stellar random velocities to zero.:
$`V_{c,x}^{inf}(\theta )`$ $``$ $`\left(V_{c,x}+V_{ei,x}\right)\mathrm{sin}\theta =\left[\kappa D\mu _{obs,x}V_{ev,x}(\theta )\right]`$ (42)
$`V_{c,r}^{inf}(\theta )`$ $``$ $`\left(V_{c,r}+V_{ei,y}\right)\mathrm{cos}\theta ={\displaystyle \frac{V_{obs,r}}{\mathrm{sin}i}}V_{ev,y}(\theta ),`$ (43)
where $`V_{c,x}^{inf}(\theta )`$ and $`V_{c,r}^{inf}(\theta )`$ are the circular velocities inferred from the observed $`x`$ proper motions and the radial velocities, respectively. These equations clearly illustrate that the invisible component of the elliptical streaming motions are indistinguishable from the circular velocity term. Here we retained the general expressions for the visible components of the elliptical streaming field. The right-hand sides (RHSs) of these two equations can be constructed from the observed $`V_r`$ and $`\mu _x`$ motions and the fitted values for distance, inclination and $`\overline{V}_{ev}`$. If the correct values for $`i`$ and $`D`$ are used, and if the true visible streaming field is subtracted from the RHSs, equation (42) should show a pure sine modulation, while only the cosine modulation should contain significant power in equation (43).
In general, neither the true inclination and distance, nor the correct elliptical streaming field will be used in the empirical determination of the RHSs of the $`V_c^{inf}`$ equations. Thus, higher order $`\theta `$ modulations will be observable in the RHSs of the $`V_c^{inf}`$ equations. The usage of an erroneous inclination ($`i_e`$) is most damaging because the estimated position angle, $`\theta _e`$, is determined from the estimated inclination such that $`\theta _e=\mathrm{arctan}(\frac{y^{}}{x\mathrm{cos}i_e})`$. In practice, a search for additional modulations in the $`V_c^{inf}`$ equations will be performed in terms of the estimated azimuth, not the true azimuth ($`\theta _t`$). To see what the consequences are, we suppose that the true inclination, $`i_t`$, can be expanded to first order around $`i_e`$. With $`\mathrm{\Delta }ii_ei_t`$ we find:
$`\mathrm{sin}\theta _t`$ $``$ $`\left(1{\displaystyle \frac{\mathrm{\Delta }i}{4}}\mathrm{tan}i_e\right)\mathrm{sin}\theta _e{\displaystyle \frac{\mathrm{\Delta }i}{4}}\mathrm{tan}i_e\mathrm{sin}3\theta _e`$ (44)
$`\mathrm{cos}\theta _t`$ $``$ $`\left(1+{\displaystyle \frac{\mathrm{\Delta }i}{4}}\mathrm{tan}i_e\right)\mathrm{cos}\theta _e{\displaystyle \frac{\mathrm{\Delta }i}{4}}\mathrm{tan}i_e\mathrm{cos}3\theta _e`$ (45)
At the inclination of the Andromeda galaxy, the $`3\theta _e`$ modulations have an amplitude of about $`\frac{2\mathrm{\Delta }i}{1\mathrm{degree}}`$ percent of the $`\theta _e`$ amplitudes, or several $`\text{km\hspace{0.17em}s}^1`$. We find that, with a sufficient number targets, such effects are easily detectable.
In fact, if we expand the RHSs of the $`V_c^{inf}`$ equations in Fourier series, their $`3\theta _e`$ coefficients can be used to calculate new estimates for the inclination: $`\mathrm{\Delta }i_c=A_{c3}\frac{4}{\mathrm{tan}i_e}`$ and $`\mathrm{\Delta }i_s=A_{s3}\frac{4}{\mathrm{tan}i_e}`$ with $`A_{c3}`$ and $`A_{s3}`$ the measured Fourier coefficients of the $`\mathrm{cos}3\theta _e`$, and $`\mathrm{sin}3\theta _e`$ modulations. Significant $`3\theta _e`$ coefficients can occur in two cases: 1) the correct elliptical streaming was determined but the estimated inclination is wrong, and 2) the right inclination was determined, but the wrong $`3\theta _e`$ components was fitted for the elliptical streaming field. The observed $`3\theta _e`$ terms can only be used to arrive at a better inclination estimate in case the correct $`\overline{V}_{ev}`$ has been subtracted. Generally speaking, that can not known to be the case. However, in both cases, the presence of significant $`3\theta _e`$ terms indicate that the model used is inadequate and that a better solution is possible. Further, the ambiguity of the meaning of any detected $`3\theta _e`$ modulation also illustrates the necessity of a multiple non-linear regression technique.
### 5.4 The Invisible Components
So far, the invisible component remains undetectable, even when the $`3\theta _e`$ constraints are included, and even if a multiple non-linear regression technique has been used. As it turns out, the lowest order ($`\theta _e`$) Fourier coefficients of the expansion of the $`V_c^{inf}`$ equations provide powerful additional constraints on the invisible elliptical streaming components.
For example, suppose that $`\overline{V}_{ei}`$ is non-existent, in that case, the $`A_{c\theta }`$ and $`A_{s\theta }`$ coefficients of the cosine and sine modulations should be identical and equal to the rotation speed of the galaxy. If $`\overline{V}_{ei}`$ is non-zero, $`A_{c\theta }`$ will not be equal to $`A_{s\theta }`$. Unfortunately, equations (42) and (43) are not sufficient to determine the three unknowns contained in the two equations. However, it is not necessary to know the values of $`V_{ei,x}`$ and $`V_{ei,y}`$ separately. Only the algebraic sum of the two invisible components, $`V_{ei,xy}(\theta )(V_{ei,x}\mathrm{sin}\theta +V_{ei,y}\mathrm{cos}\theta )`$, occur in equations (31) and (32). Thus, in these two equations, $`V_{ei,xy}(\theta )`$ can be replaced by $`(A_{s\theta }+A_{c\theta })V_c(\mathrm{sin}\theta +\mathrm{cos}\theta )`$. In a multiple regression technique, the circular velocity thus becomes a to-be-fitted parameter. As a result, the $`V_c^{inf}`$ equations will also allow for the determination of both invisible components of the elliptical streaming field.
The arguments presented above show that, with the additional constraint that the $`V_c^{inf}`$ equations only show a $`\theta `$ modulations, all parameters of the model can be determined. Because even the “invisible” components of the elliptical streaming field can be determined experimentally, the inferred distance, inclination rotation speed and galaxian space motion can be measured without significant systematic errors.
### 5.5 Final Accuracies
In the SIM book it is suggested one-hundred targets are observed in both M 31 and M 33 and twenty-five in M 81. Our discussion above indicate that systematic errors will be unimportant so that the final attainable distance error is inversely proportional to the square-root of the number of target stars. With such a moderate number of targets and a proper motion uncertainty of 4$`\mu \text{as\hspace{0.17em}yr}^1`$, extremely accurate distances can be determined for the nearest spirals: $`\mathrm{\Delta }D`$ 0.7%, 1.9% and 10% $`(\times \frac{\mathrm{\Delta }\mu }{4\mu \text{as\hspace{0.17em}yr}^1})`$, for M 31, M 33 and M 81, respectively.
It should be kept in mind that these considerations neglect several important aspects that will be subject of future study. First, systematic effects internal to the galaxies such as spiral arm streaming motions or runaway stars will reduce the final achievable results. On the positive side, it will be possible to identify “deviant” objects using the single star method so that they can be eliminated from the target list. Further, it may very well be possible that a global fit to the proper motions and radial velocities of the targets will produce significantly tighter results if we impose smoothness criteria for the radial variation of rotation speed and inclination. It might also be possible to use external information on the radial gradients of the inclination and rotation curve from the H I velocity field to further decrease the errors. This too will be investigated in the near future.
We have generated numerical models that describe the stellar motions in three potential SIM targets: M 31, M 81 and M 3. The stellar disks of these model galaxies are inclined by 77°, 56°, and 57° with respect to the line-of-sight, respectively. The motions of the stars have circular velocity components of 270, 213 and 97.3 $`\text{km\hspace{0.17em}s}^1`$, while we add a random component to the targets of 10 $`\text{km\hspace{0.17em}s}^1`$ in all three directions. We computed several models with either fixed inclinations or a small inclination gradient. We also varied the elliptical streaming component, from non-existent to strong (approximately 20% of the circular velocity), where we tried several angular dependencies consistent with the toy model described in section 5.2 (i.e., $`\theta `$ components only). Our method of analysis of this model data is described by the poor-man’s approach of section 5.3. We ran six different models with 100-600 targets. Each of these models were “observed” $`N_{try}=200`$ times where for each try we added random terms to the stellar space motions to simulate observational errors. To the radial velocity we added random errors from 2.5 to 10 $`\text{km\hspace{0.17em}s}^1`$, in steps of 2.5 $`\text{km\hspace{0.17em}s}^1`$, while we used a large range of proper motion errors (between 1 and 50 $`\mu \text{as\hspace{0.17em}yr}^1`$).
The results for the three galaxies are summarized in figures Galaxy Distances via Rotational Parallaxes-Galaxy Distances via Rotational Parallaxes. Each of these three figures contain three rows and four columns. The results for the smallest number of targets are displayed in the top row, those for the largest number of stars in the bottom panels. In the left two columns we plot the inferred rms and systematic distance errors, respectively. The two systemic error flags (SEFs) are plotted in the two right columns. For M 31, for which the observed proper motions are largest, the rms errors are encouragingly small, even for a small number of targets. The inferred systematic errors are also small, which is confirmed by the small SEF values. For the other two galaxies, significant systematic errors result from the poor-man’s approach. This is probably a result of the fact that these galaxies have smaller rotational proper motions. On the positive side, the SEFs are also clearly raised, indicating that the poor-man’s solutions are erroneous.
### 5.6 Beyond SIM
If the systemic velocity is zero, the systematic variation of of the distancew $`D`$<sup>9</sup><sup>9</sup>9$`D(R,\theta )=\sqrt{[D(R=0)+R\mathrm{sin}\theta \mathrm{sin}i]^2+[R\mathrm{cos}\theta ]^2}`$ along an annulus will introduce a deviation from the linear behavior of eqn. (29). However, the deviation due to the “proximity effect”, $`\mathrm{\Delta }\mu _y^{}`$, is rather small: $`\mathrm{\Delta }\mu _y^{}4,0.14,0.09`$ and $`0.01`$ $`\text{mas}\text{yr}^1`$ for the LMC, M 31, M 33 and M 81, respectively. The constant $`\beta _y^{}`$ in eqn. (29) has a distance dependence as well so that the systemic velocity terms will also contribute to $`\mathrm{\Delta }\mu _y^{}`$. In fact, because the contribution to $`\mathrm{\Delta }\mu _y^{}`$ from internal motions is down-weighted by the $`1/\mathrm{tan}i`$ term, the systemic contribution to $`\mathrm{\Delta }\mu _y^{}`$ tends to outweigh the former. However, for “reasonable” values of the systemic velocity, $`\mathrm{\Delta }\mu _y^{}`$ will be still be rather small.
With the possible exception of the Magellanic Clouds, we do not expect that it will be possible to determine the proximity effect with SIM data, in part due to the small number of stars that SIM can measure. Recall that a velocity dispersion of 10 $`\text{km\hspace{0.17em}s}^1`$ corresponds to about 3 $`\text{mas}\text{yr}^1`$, or about twenty times the proximity effect, at the distance of M 31. Assuming that measuring $`\mathrm{\Delta }\mu _y^{}`$ at ten $`V_r`$ positions would suffice to determine the proximity effect, we estimate that at least $`10\times 20^2=4000`$ targets are required. This would take about 6,000 hours or about 20% of the science time available during the SIM mission.
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# Environment-induced dynamical chaos
## Abstract
We examine the interplay of nonlinearity of a dynamical system and thermal fluctuation of its environment in the “physical limit” of small damping and slow diffusion in a semiclassical context and show that the trajectories of c-number variables exhibit dynamical chaos due to the thermal fluctuations of the bath.
The interplay of nonlinearity of a dynamical system and thermal fluctuations of its environment has been one of the major areas of investigation in the recent past. These studies have enriched our understanding of nonequilibrium processes in several contexts, such as, symmetry between the growth and decay of classical fluctuations in equilibrium , interesting topological features of patterns of paths of large fluctuations in nonlinear systems , existence of generalized nonequilibrium potential , influence of nonlinearity on dissipation in multiphoton processes and higher order diffusion in a nonlinear system , etc. While the development in these areas is largely confined to classical domain we examine a related issue in the semiclassical context. Since quantization is likely to add a new dimension to the interplay of nonlinearity and stochasticity in a weakly dissipative system, it is worthwhile to consider the physical limit of small damping and slow diffusion due to thermal fluctuations of the environment and look for the thermal fluctuations-induced features of nonlinearity in the dynamics. In this communication we specifically explore some interesting aspects of dynamical chaos in a driven bistable system whose origin lies at the fluctuations of the environmental degrees of freedom.
To describe the dissipative quantum dynamics of a system we first consider the traditional system-reservoir model developed over the last few decades . The Hamiltonian of the bare system is coupled to an environment modeled by a reservoir of harmonic oscillator modes characterized by a frequency set $`\{\mathrm{\Omega }_j\}`$. The quantum dynamics is generated by the overall Hamiltonian operator H for the system, environment and their coupling as follows;
$$H=H_0+\mathrm{}\underset{j=1}{\overset{\mathrm{}}{}}\mathrm{\Omega }_jb_{j}^{}{}_{}{}^{}b_j+\mathrm{}\underset{j=1}{\overset{\mathrm{}}{}}\left[g(\mathrm{\Omega }_j)b_j+g^{}(\mathrm{\Omega }_j)b_{j}^{}{}_{}{}^{}\right]x,$$
(1)
where
$$H_0=\frac{p^2}{2}+V(x),$$
(2)
defines the usual kinetic and potential energy terms corresponding to the system, $`x`$ and $`p`$ being the position and momentum operators, respectively. The second and the third terms in (1) specify the reservoir modes and their linear coupling to the system. $`g(\mathrm{\Omega })`$ denotes the system-reservoir coupling constant.
Systematic elimination of the reservoir modes in the usual way, using Born and Markov approximations leads one to the following standard reduced density matrix equation of the system only .
$`{\displaystyle \frac{d\rho }{dt}}`$ $`=`$ $`{\displaystyle \frac{i}{\mathrm{}}}[H_0,\rho ]+{\displaystyle \frac{\gamma }{2}}(1+\overline{n})(2a\rho a^{}a^{}a\rho \rho a^{}a)`$ (4)
$`+{\displaystyle \frac{\gamma }{2}}\overline{n}(2a^{}\rho aaa^{}\rho \rho aa^{}).`$
Here the system operator co-ordinate $`x`$ is related to the creation and the annihilation operators $`a^{},a`$ respectively as $`x=(\frac{1}{\sqrt{2\omega }})(a+a^{})`$. $`\omega `$ is the linearised frequency of the system. Also the spectral density function of the reservoir is replaced by a continuous density $`𝒟(\omega )`$. $`\gamma >0`$ is the limit of $`2\pi |g(\omega )|^2\frac{𝒟(\omega )}{\omega }`$ as $`\omega 0_+`$ and is assumed to be finite. $`\gamma `$ is the relaxation or dissipation rate and $`\overline{n}\gamma `$ is the diffusion coefficient $`D`$. $`\overline{n}(=[exp\left(\frac{\mathrm{}\omega }{kT}\right)1]^1)`$ is the average thermal photon number of the reservoir. The terms analogous to Stark and Lamb shifts have been neglected.
The first term in Eq.(3) corresponding to the dynamical motion of the system refers to Liouville flow. The terms containing $`\gamma `$ arise due to the interaction of the system with the surroundings. While $`\frac{\gamma \overline{n}}{2}`$ terms denote the diffusion of fluctuations of the reservoir modes into the system mode, $`\frac{\gamma }{2}`$ terms refer to the loss of energy from the system into reservoir. In the limit $`T0,i.e.,\overline{n}0`$ the system is influenced by pure quantum noise or vacuum fluctuations \[ Note that 1 in $`(\overline{n}+1)`$ in Eq.(3) corresponds to the vacuum\].
Our next task is to go over from a full quantum operator problem described by the Eq.(3) to an equivalent c-number problem . To this end we consider the quasi-classical distribution function W ($`x,p`$, t) of Wigner. $`x,p`$ are now c-number variables. Rewriting Eq.(3) in a quasi-classical language we obtain
$`{\displaystyle \frac{W}{t}}`$ $`=`$ $`p{\displaystyle \frac{W}{x}}+{\displaystyle \frac{V}{x}}{\displaystyle \frac{W}{p}}+{\displaystyle \frac{\gamma }{2}}\eta \left({\displaystyle \frac{xW}{x}}+{\displaystyle \frac{pW}{p}}\right)`$ (7)
$`+{\displaystyle \frac{\gamma \eta \mathrm{}}{2\omega }}(\overline{n}+1){\displaystyle \frac{^2W}{x^2}}+{\displaystyle \frac{\gamma \eta \mathrm{}\omega }{2}}(\overline{n}+1){\displaystyle \frac{^2W}{p^2}}`$
$`+{\displaystyle \underset{n1}{}}{\displaystyle \frac{\mathrm{}^{2n}(1)^n}{2^{2n}(2n+1)!}}{\displaystyle \frac{^{2n+1}V}{x^{2n+1}}}{\displaystyle \frac{^{2n+1}W}{p^{2n+1}}}.`$
$`\eta `$ in Eq.7 is a parameter used to identify the environment-induced effect on the dynamics described by Eq.(4) (kept for bookkeeping in the calculation and put $`\eta =1`$ at the end).
In the semiclassical limit $`\mathrm{}\omega <<kT`$, we have $`\overline{n}+1\overline{n}`$ and $`D\gamma kT`$ so that Eq.(7) reduces to
$`{\displaystyle \frac{W}{t}}`$ $`=`$ $`p{\displaystyle \frac{W}{x}}+{\displaystyle \frac{V}{x}}{\displaystyle \frac{W}{p}}+{\displaystyle \frac{\gamma }{2}}\eta \left({\displaystyle \frac{xW}{x}}+{\displaystyle \frac{pW}{p}}\right)`$ (10)
$`+D\eta \left({\displaystyle \frac{1}{\omega ^2}}{\displaystyle \frac{^2W}{x^2}}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{^2W}{p^2}}\right)`$
$`+{\displaystyle \underset{n1}{}}{\displaystyle \frac{\mathrm{}^{2n}(1)^n}{2^{2n}(2n+1)!}}{\displaystyle \frac{^{2n+1}V}{x^{2n+1}}}{\displaystyle \frac{^{2n+1}W}{p^{2n+1}}}.`$
The overall dynamics described above is a superposition of two contributions, i. e. , the Liouville-Wigner dynamics and the system-reservoir dissipative dynamics. That the two contributions act independently is an assumption. The master equation (3) \[or its Wigner function version (4)\] is the most popular one in quantum optics. It has been extensively used , for the strongly nonlinear processes like three-wave, four-wave mixing and strong coherent light-matter interaction phenomena. The equation has also been applied in the context of chaos, e. g., in the dissipative standard map , dissipative logistic map , semiclassical theory of quantum noise in open chaotic systems and in the studies of decoherence in relation to chaos for analysis of quantum-classical correspondence .
In the semiclassical ($`\mathrm{}0`$) limit the dissipative quantum dynamics can be conveniently described by “WKB-like” ansatz (we refer to “ WKB-like” since we are considering more that one dimension. Traditional WKB refer to one dimension only) of Eq.(5) for Wigner function of the form
$$W(x,p,t)=Z(t)\mathrm{exp}(\frac{s}{\mathrm{}}).$$
(11)
Here $`Z(t)`$ is a prefactor and $`s(x,p,t)`$ is a classical action which is a function of c-number variables $`x`$ and $`p`$ , satisfying the following Hamilton-Jacobi equation
$`{\displaystyle \frac{s}{t}}+p{\displaystyle \frac{s}{x}}{\displaystyle \frac{s}{x}}{\displaystyle \frac{s}{p}}{\displaystyle \frac{\gamma }{2}}\eta (x{\displaystyle \frac{s}{x}}+p{\displaystyle \frac{s}{p}})+\eta ({\displaystyle \frac{s}{x}})^2`$ (12)
$`+\eta \omega ^2({\displaystyle \frac{s}{p}})^2+{\displaystyle \underset{n1}{}}{\displaystyle \frac{x^{2n}(1)^{3n+1}}{2^{2n}(2n)!}}{\displaystyle \frac{^{2n+1}V}{x^{2n+1}}}{\displaystyle \frac{s}{p}}=0.`$ (13)
The derivation of Eq.(7) is based on the consideration of the “physical limit” of weak dissipation and slow fluctuations in the sense $`\frac{D_1}{\mathrm{}^2}\frac{1}{\mathrm{}}`$ where $`D_1=\frac{D}{2\omega ^2}`$ (note that $`D_1`$ and $`\mathrm{}`$ have same dimension). The above equation can be solved by integrating the Hamiltonian equations of motion,
$`\dot{x}`$ $`=`$ $`p{\displaystyle \frac{\gamma }{2}}\eta x+2\eta P`$ (14)
$`\dot{X}`$ $`=`$ $`P{\displaystyle \frac{\gamma }{2}}\eta X`$ (15)
$`\dot{p}`$ $`=`$ $`V^{}+{\displaystyle \frac{\gamma }{2}}\eta p2\omega ^2\eta X{\displaystyle \underset{n1}{}}{\displaystyle \frac{(1)^{3n+1}}{2^{2n}(2n)!}}{\displaystyle \frac{^{2n+1}V}{x^{2n+1}}}X^{2n}`$ (16)
$`\dot{P}`$ $`=`$ $`V^{\prime \prime }X+{\displaystyle \frac{\gamma }{2}}\eta P{\displaystyle \underset{n1}{}}{\displaystyle \frac{(1)^{3n+1}}{2^{2n}(2n+1)!}}{\displaystyle \frac{^{2n+2}V}{x^{2n+2}}}X^{2n+1}`$ (17)
which are derived from the following Hamiltonian $`H_{eff}`$
$`H_{eff}`$ $`=`$ $`pPV^{}X{\displaystyle \frac{\gamma }{2}}\eta (xP+pX)+\eta (P^2+\omega ^2X^2)`$ (19)
$`+{\displaystyle \underset{n1}{}}{\displaystyle \frac{(1)^{3n+1}}{2^{2n}(2n+1)!}}{\displaystyle \frac{^{2n+1}V}{x^{2n+1}}}X^{2n+1}.`$
Here we have put $`\frac{s}{x}=P`$ and $`\frac{s}{p}=X`$. The introduction of additional degree-of-freedom by incorporating the auxiliary momentum (P) the and co-ordinate (X) makes the system effectively a two-degree-of-freedom system. The origin of these two variables is the thermal fluctuations of the reservoir. It is easy to identify the environment-related terms containing $`\eta `$ in Eqs.(7-9). The auxiliary Hamiltonian (9) is therefore not to be confused with the microscopic Hamiltonian (1) which describes a system in contact with a reservoir with infinite degrees of freedom. Although the phase space trajectories concern fluctuations of c-number variables in the formal sense, because of the equations of motion (8) described by a Hamiltonian (9), the motion is strictly deterministic. The experiments on the corresponding classical version of the problem by Luchinsky and McClintock have demonstrated that a trajectory of fluctuation is indeed a part of physical reality. We emphasize, however, here a number of distinguishing features in this context. While the studies by Luchinsky and McClintock and Graham and Tel concern overdamped limit, we consider here a weakly dissipative system. Furthermore because of the quantum correction, the phase space trajectories of fluctuations are significantly modified by semiclassical features. The introduction of these quantum features at a semiclassical level through a c-number Hamiltonian description of a dissipative evolution in the physical limit of weak damping and slow diffusion due to thermal noise is the essential content of the present work. Since the Eq.(9) describes deterministic evolution under nonlinear potential, the pattern of trajectories of fluctuations may display chaotic behaviour. In what follows we investigate this dynamical aspect of the dissipative system.
The equations derived in the weak thermal noise limit for the weakly dissipative semiclassical systems are fairly general. For illustration we now consider a simple model system Hamiltonian $`H_0`$ (see Eq.2)
$$H_0=\frac{p^2}{2}+ax^4bx^2+gxcos\mathrm{\Omega }_0t$$
(20)
which describes a bistable potential driven by a time-periodic field. $`a`$ and $`b`$ are the constants of the potential $`V(x)`$. The fourth term in (10) includes the effect of coupling of the system as well as the strength of the field of frequency $`\mathrm{\Omega }_0`$. For the Hamiltonian (10) the Eqs.(8) read as
$`\dot{x}`$ $`=`$ $`p{\displaystyle \frac{\gamma }{2}}\eta x+2\eta P`$ (21)
$`\dot{X}`$ $`=`$ $`P{\displaystyle \frac{\gamma }{2}}\eta X`$ (22)
$`\dot{p}`$ $`=`$ $`4ax^32bx+g\mathrm{cos}\mathrm{\Omega }t+\eta ({\displaystyle \frac{\gamma }{2}}p2\omega ^2X)3axX^2`$ (23)
$`\dot{P}`$ $`=`$ $`(12ax^22b)X+{\displaystyle \frac{\gamma }{2}}\eta PaX^3`$ (24)
which are derivable from the auxiliary Hamiltonian
$`H_{eff}`$ $`=`$ $`Pp(4ax^32bx+g\mathrm{cos}\omega t)X{\displaystyle \frac{\gamma }{2}}\eta (xP+pX)`$ (26)
$`+\eta (P^2+\omega ^2X^2)+axX^3.`$
The system Hamiltonian (10) has served as a standard paradigm for a number of theoretical and experimental investigations over many years. For the present purpose we choose the following parameter values; $`a=\frac{1}{4}`$, $`b=\frac{1}{2}`$, $`\mathrm{\Omega }_0=6.07`$ and $`g=10`$. Since we are considering the physical limit of weak damping and small diffusion we take the initial conditions for the auxiliary variables X and P (which originate from the fluctuations due to environment) as $`P=0`$ and $`X0`$ (we have used $`X=1.5\times 10^6`$). This ensures a vanishing Hamiltonian $`H_{eff}`$ for the entire numerical investigation that follows below.
We first consider a specific trajectory with the initial condition $`p=0`$ and $`x=2.512`$ for small values of $`\gamma `$ (typical $`0.01`$). Under this condition($`\eta =1`$) the system is vanishingly coupled to the surroundings and consequently the dynamical behaviour is effectively due to the weak dissipation only. We illustrate this situation in Fig. 1 in terms of a Poincare map for the phase space which exhibits strong global chaos. On the other hand when the parameter $`\eta `$ is switched off ($`\eta =0`$) the system displays typical weak chaos (Fig.2). The similar behaviour has been observed for other sets of initial conditions for $`x`$ and $`p`$ (We have not reproduced them here for the sake of brevity).
The effect of weak dissipation and slow diffusion due to thermal fluctuations from the surroundings can be seen in the case of other sets of initial conditions also. For the initial condition $`p=0`$ and $`x=2.49`$ one observes for $`\gamma =0.01`$ dissipative strong chaos($`\eta =1`$). This is illustrated in Fig.3. It is interesting to note that for $`\eta =0`$ the same trajectory gets localized in the left well as a regular one as shown in Fig.4. The weakly chaotic and regular trajectories in Figs. 2 and 4, respectively, are purely semiclassical in nature (in the absence of any coupling to the surroundings). The strong global chaotic behaviour as shown in Figs. 1 and 3 has therefore its origin in the thermal fluctuations of the reservoir. In other words the chaotic behaviour or its enhancement is exclusively due to thermal fluctuations from the surroundings, which becomes appreciable even in the physical limit of weak damping and slow diffusion of the thermal noise within the semiclassical description. We have checked this assertion for some other values of the initial conditions for the system oscillator.
The reduction of the system-reservoir Hamiltonian description \[H in Eq.1\] for a dissipative quantum system to an auxiliary Hamiltonian description \[$`H_{eff}`$ in Eq.9\] effectively reduces the infinite-degree-of-freedom system to a two-degree-of-freedom system where the auxiliary degree of freedom characterized by X and P owe their origin in the fluctuations of the reservoir. Since X and P appear as the multiplicative factors in the auxiliary Hamiltonian $`H_{eff}`$, the weak thermal noise limit makes $`H_{eff}`$ a vanishing Hamiltonian. The observed semiclassical chaos may therefore be regarded as a dynamical manifestation of the interplay of nonlinearity and thermal fluctuations.
In this paper we have examined the weak thermal noise limit of a semiclassical dissipative nonlinear system. We have shown that the vanishing Hamiltonian method can be suitably extended to follow the phase space trajectories of fluctuations of c-number variables which exhibit dynamical chaos. In view of the accessibility of the model to analogue electronic circuits we believe that the results discussed bear further experimental relevance in the semiclassical context.
B. C. Bag is indebted to the CSIR for partial financial support and to J. Ray Chaudhuri for helpful discussions.
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# Magnetic lensing of Extremely High Energy Cosmic Rays in a Galactic wind
## 1 Introduction
The origin and nature of the highest energy cosmic rays so far detected stands as a puzzle for contemporary astrophysics . The flux of protons with energy around and above 70 EeV (1 EeV = $`10^{18}`$ eV) should be significantly attenuated over distances of order 100 Mpc due to their interaction with the cosmic microwave background . Nuclei should be attenuated over even shorter distances . Nonetheless, fourteen events with estimated energy above 100 EeV have been detected so far, largely exceeding the expected fluxes if the sources are at cosmological distances. This suggests a “local” origin of extremely high energy cosmic rays (EHECRs). The puzzle arises because the angular distribution of the fourteen EHE events is consistent with isotropy (given the limited statistics and insufficient sky coverage) and there are no known sources near their arrival directions and inside our 100 Mpc neighborhood considered to be a potential site for acceleration of cosmic rays to such enormous energies.
A potential solution to the puzzle is that EHECRs are protons or nuclei that do indeed originate in sources within a 100 Mpc neighborhood of Earth, but their arrival directions do not point to their place of origin because their trajectories are significantly bent as they traverse intervening magnetic fields. The regular component of the magnetic field in the Milky Way leads to sizeable deflections and other magnetic lensing effects upon ultra high energy charged cosmic rays, such as flux (de)magnification and multiple image formation . This is the case, however, only if the ratio $`E/Z`$ between energy and electric charge of the CRs is below approximately 30 EeV. The EHECRs would thus be severely affected by the regular component of the galactic magnetic field only if they have a significant component which is not light.
It has recently been speculated that all the events so far detected at energies above $`10^{20}`$ eV may originate from M87 in the Virgo cluster, if the Galaxy has a rather strong and extended magnetic wind. Indeed, such extreme galactic magnetic wind is compatible with an origin for all EHE events at less than $`20^{}`$ from the direction to the Virgo cluster, if two out of the thirteen events considered are He nuclei, the rest being protons. One event was excluded from the dataset due to the large uncertainties in its energy determination.
In this paper we further develop the analysis of the scenario put forward in . The determination of the deflections of CR trajectories is insufficient to test the consistency of the scenario with the observational data and to determine its generic predictions, because magnetic lensing also produces huge flux (de)magnification and multiple image formation, effects that we analyse here and which turn out to be crucial to establish the detailed features of the model.
## 2 Flux enhancement by magnetic lensing in a galactic wind
We consider an azimuthal magnetic field with strength given by
$$B=B_0\frac{r_0}{r}\mathrm{sin}\theta \mathrm{tanh}(r/r_s)$$
(1)
as a function of the radial (spherical) coordinate, $`r`$, and the angle to the north galactic pole, $`\theta `$. The distance from Earth to the galactic center is $`r_0=8.5`$kpc and the local value of this wind field is taken as $`B_0=7\mu `$G. This field has the $`\mathrm{sin}\theta /r`$ dependence adopted in ref. with an extra smoothing in the Galactic center region, given by the $`\mathrm{tanh}(r/r_s)`$ factor, introduced to avoid unphysical divergences at small radii. For definiteness we took $`r_s=5`$ kpc but the results are not dependent on the precise way in which this smoothing is performed. We also adopted a 1.5 Mpc cutoff for the extension of the field as in ref. .
The trajectories of CR protons in the Galactic wind magnetic field are obtained by backtracking antiprotons leaving the Earth. The azimuthal nature of this field bends all the trajectories towards the north galactic pole . The incident direction outside the region of influence of the wind points at less than $`15^{}`$ from the north pole for all the observed EHE events, except for the two most energetic ones. The two most energetic events could also come from that cone only if their electric charge is assumed to be larger, for example if they are He nuclei. This fact has been exploited to suggest that all the EHE events may have a common source, M87 in the Virgo cluster , which is indeed at $`b=74.4^{}`$, i.e. not far from the north galactic pole. Clearly at this level one has to be satisfied with this kind of accuracy in the pointing since the wind model is highly idealised and one is also neglecting the additional deflections that should take place near the source and in the travel through the intergalactic medium. This will also require that our conclusions be based only on the general qualitative features of the model, rather than on its specific details.
It has been pointed out that large scale magnetic fields not only deflect CR trajectories but that they also act as a giant lens (de)magnifying the fluxes received from different directions and leading to multiple image formation . For a given source this lensing effect varies with the energy and thus can enhance or attenuate the fluxes differently for different energies. This effect turns out to be quite strong in the Galactic wind model under consideration, so it has to be taken into account in a complete analysis of this scenario.
The magnification due to the lensing effect is computed as the ratio between the area subtended by a parallel bundle of particles arriving to the Galactic wind region and that subtended by the same particles at their arrival to the Earth . Fig. 1 shows the contour plots of equal magnification for three values of the $`E/Z`$ ratio. Each point denotes the arrival direction of a CR to the Earth in galactic coordinates $`(\mathrm{},b)`$. We observe that huge magnifications, in excess than a factor of 100, are attained in large regions of the sky. This is the case for most directions with $`b>15^{}`$ for $`E/Z=150`$ EeV and with $`b>15^{}`$ for $`E/Z=125`$ EeV. The critical lines (lines where the magnification diverges and that correspond to the caustic lines in the source plane) move quickly to the south as the energy decreases. Most of the sky is swept by the critical lines as $`E/Z`$ varies between 150 EeV and just below 100 EeV. Since the magnification is huge in the regions close to these lines, a strong enhancement in the detection probability of events with energies close to that at which the caustics cross the source direction is expected.
## 3 Multiple images
The presence of critical curves in the magnification maps indicates that the magnetic lensing effect of the Galactic wind produces multiple images of the source for some source directions in the range of energies considered. This fact is most clearly seen from the plots of the “sky sheets”, representing the projection of a regular grid of observing directions at Earth to the directions outside of the Galactic wind region. This kind of plots map the arrival direction(s) at which a CR is observed at Earth to the direction from where the CR arrived to the region of influence of the magnetic field. Fig. 2 shows this sheet in polar coordinates centered in the north pole for three values of the $`E/Z`$ ratio. The stretching of this sky sheet reflects the magnification for those source directions, very stretched regions being strongly demagnified and those where the sky sheet is densely contracted being very magnified. The fold lines correspond to the position of the caustics, where the magnification diverges. Sources located in regions where the sheet is folded are observed at Earth from several different directions. For $`E/Z=175`$ EeV a couple of folds ending into two cusps (a lip) develop. A source inside the folds has three different images. Note however that the folds cover only a tiny fraction of the sky and thus this is a rare effect at these energies. For $`E/Z=125`$ EeV a complete blob has developed on the sheet. The blob is connected to the rest of the surface through a diamond shaped caustic (this is due to the Earth eccentric position in the Galaxy: for observers in the symmetry axis, e.g. at the Galactic center, the diamond would shrink to a point<sup>1</sup><sup>1</sup>1These caustics are analogous to those associated to elliptical lenses, or spherical ones with external shear, in gravitational lensing , for which the central diamond shrinks to a point-like caustic in the circularly symmetric limit.). The southernmost critical line in the second panel of Fig. 1 corresponds to the diamond caustic in Fig. 2, while the northern critical line corresponds to the circular fold caustic. The huge amplifications found for all the directions in the northern sky above the southernmost critical line reflect the fact that a large fraction of the sky seen on Earth is shrinked inside the blob, and hence the sheet is highly contracted there. While the region inside the diamond caustic (leading to five images) is still tiny in the source sky, the blob (leading to three images) covers a region of $`5^{}`$ diameter. We have also drawn, with a bold solid line, the directions for which the cosmic rays would arrive to the Earth along the galactic equator. It is clear that all cosmic rays with $`E/Z=125`$ EeV with arrival directions in the northern galactic hemisphere actually enter the galactic wind from directions less than $`2.5^{}`$ away from the center of the blob. At smaller energies a still larger fraction of observing directions are swallowed into the blob (this corresponds to the motion of the critical lines towards the south in Fig. 1), which at $`E/Z=100`$ EeV has an angular diameter of the order of $`15^{}`$. At $`E/Z=75`$ EeV the diamond has nearly disappeared, and a new blob starts developing on top of the previous blob. At this energy the angular diameter of the main blob, corresponding to the source positions leading to multiple images, is already larger than $`30^{}`$.
It is important to notice that only if the blob is on top of the source position (eventually displaced by the energy dependent deflections due to the extragalactic magnetic fields) the cosmic rays are able to reach the Earth along directions above the southernmost critical line in Fig. 1, which for $`E/Z=125`$ EeV already covers the whole northern hemisphere. Hence, protons with $`E<125`$ EeV can only reach the northern hemisphere as extremely magnified ($`A>10^2`$) secondary images.
## 4 Generic predictions
We now discuss some generic features of the scenario in which all the CR events with energies above $`10^{20}`$ eV so far detected originate from M87 in the Virgo cluster, their trajectories having been significantly bent by this rather strong and extended galactic magnetic wind . These generic properties may serve to test the validity of the scenario as the data on EHECRs increase.
Detailed predictions depend upon the exact nature and strength of the galactic wind, as well as upon the precise deflections suffered by the CR trajectories from the source in M87 up to their entrance to the galactic wind. Nevertheless, the generic features of any scenario compatible with current observations, mainly determined by the focusing properties of the magnetic wind, should be similar to those in the highly idealised model analysed here.
In the model under consideration, all thirteen EHECR events enter the galactic wind region along different directions, most of them just between $`2^{}`$ and $`5^{}`$ away from the direction to the north galactic pole, and at most separated by $`12^{}`$ from the vertical direction. Intergalactic magnetic fields are assumed to be responsible for these energy-dependent deviations from the direction to M87.
This scenario predicts a strong asymmetry between the north and south galactic hemispheres. Consider for instance that all EHECRs enter the galactic halo less than $`15^{}`$ away from the direction to the north galactic pole. The lines in Figure 3 display the southernmost possible arrival directions on Earth, for several values of $`E/Z`$.
This scenario thus implies that no charged cosmic rays should arrive on Earth from directions below the lines drawn in Fig. 3 for values of $`E/Z`$ larger than that indicated along each line. Certainly CRs may arrive below these lines if they enter the galactic wind with an inclination larger than $`15^{}`$ from its symmetry axis, but then they do so with smaller magnifications (or rather with large demagnifications) as larger is the inclination. Thus, even if there were other EHECR sources as powerful as M87 in our local (less than 100 Mpc) neighborhood, their flux will not be significantly magnified (or rather they will be significantly demagnified) if the CRs enter the galactic wind far away from the direction to the north galactic pole. Only if $`E/Z`$ is above a few times $`10^{21}`$ eV does the observed flux approach its unlensed value, and the CRs arrival directions point to their true source location.
Even if there is a unique source, CRs at different energies should enter the galactic wind from different directions if they have suffered magnetic deflections in their way. It is nevertheless instructive to consider fixed incoming directions, to further illustrate some of the generic features of lensing by the galactic wind already discussed in the previous sections. Figure 4 displays the energy-dependent magnification of the flux of CRs that enter the galactic magnetic wind from directions ($`\mathrm{},b)=(270^{},88^{})`$ and $`(270^{},85^{})`$, for the principal (P) as well as for the secondary images (A, B). Figure 5 displays the change in the observed arrival directions of CRs with these entrance directions as the energy is lowered down to $`E/Z=75`$ EeV.
The principal image is magnified by a factor of order 10 at $`E/Z200`$ EeV, is further amplified at intermediate energies, and then its magnification starts to rapidly decrease while $`E/Z`$ is still above 100 EeV. Its apparent position moves south as the energy decreases. Secondary images appear at the energy at which the caustic crosses the source position, with formally divergent magnifications. One of the secondary images moves north and remains highly magnified, with an amplification factor above 100, while the other moves south and is quickly demagnified. The secondary images appear and remain in the opposite east-west hemisphere than that in which the principal image is seen.
The general features displayed in these examples are generic to different entrance directions. However, the precise energy and location at which secondary images form, the energy at which the principal image is most magnified, and the exact amount of magnification attained depend sensibly on the entrance direction. For instance, if the entrance direction is more than about $`8^{}`$ away from the north galactic pole, then secondary images appear only below 100 EeV, and the principal image acquires smaller and smaller maximum magnification as the entrance direction is farther away from the vertical.
The expected energy and angular distribution of observed EHECRs can be exemplified considering again some fixed entrance directions and assuming that the differential flux injected by the source scales for instance as $`E^{2.7}`$. We take the detecting system to have the same efficiency at all energies within the range considered, which we divide in 50 bins of equal detection probability. We consider values of $`E/Z`$ larger than 100 EeV only. Figure 5 displays the arrival directions of the events that would be detected in each case.
As already discussed, we see that there are no events at southern galactic latitudes below a certain energy threshold. Indeed, the apparent position of the principal image crosses the galactic equator at $`E/Z133`$ EeV in the left panel, $`E/Z144`$ EeV in the right, and the events in the secondary images are all below the energy of the caustic ($`E/Z127`$ EeV in the left panel, $`E/Z108`$ EeV in the right). Notice also that the divergence in the magnification at the caustic gives anyhow a finite number of events once it is convoluted with the differential spectrum of the incident CRs .
## 5 Conclusions
We have analysed flux magnification and multiple image formation in the galactic magnetic wind scenario put forward in . In this scenario all observed arrival directions of EHECRs are compatible with a common origin in M87 if intergalactic magnetic fields provide the extra deflection (of order $`20^{}`$) needed to fine-tune the incoming particles in the appropriate direction as they enter the wind. We find that magnification factors well above 100 are attained in a significant energy range, with $`E/Z`$ below 150 EeV. This reduces the energy requirements upon the source, that would need to be a factor of more than 100 less powerful than if unlensed to provide the same observed flux in this energy range.
One of the definite predictions of this model is the strong asymmetry expected between events arriving from the northern and southern galactic hemispheres. Although with the present EHE data, which involves only the northern terrestrial hemisphere and hence mainly the northern galactic one, it is not yet possible to test this asymmetry, the future operation of the Auger observatory, that will provide good coverage of the southern skies, will allow to check the viability of this model. In particular, a very strong suppression of events above $`E150`$ EeV should be present at latitudes below $`b30^{}`$ for the scenario to survive. Another general feature is that an abrupt kink in the overall spectrum should appear when the secondary images disappear, i.e. when the energy increases beyond the energy of the caustic crossing. Although no particular feature of this kind is apparent in the present data, an increased statistics in the northern skies would be desireable to definitely confront this prediction with observations.
As a final remark, we would like to point out that if a galactic wind is indeed present in the Milky Way, but with a smaller overall strength (so that for instance locally it is below the 2–3 $`\mu `$G amplitude of the spiral field which is inferred from observations, rather than the 7 $`\mu `$G adopted here), it could in any case have interesting observational consequences, especially if EHECRs have a component which is not light. In particular, one can think of a scenario in which an extragalactic source near the north pole, such as M87, produces heavy nuclei with $`E>10^{20}`$ eV (and not necessarily protons at these energies), and their flux is strongly amplified by the galactic wind field. In this case a transition to a heavy composition could result at extremely high energies <sup>2</sup><sup>2</sup>2One would also have to take into account here that for propagation distances beyond $`10`$ Mpc the photodisintegration of nuclei out of the CMB photons starts to be important for $`E>Z\times 10^{19}`$ eV ..
###### Acknowledgments.
Work partially supported by ANPCyT, CONICET and Fundación Antorchas, Argentina.
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# Bias and temperature dependence of the 0.7 conductance anomaly in Quantum Point Contacts
## I Introduction
The quantized conductance through a narrow quantum point contact (QPC), discovered in 1988, is one of the key effects in mesoscopic physics. The quantization of the conductance in units of the spin degenerate conductance quantum, $`G_2=2e^2/h`$, can be explained within a single-particle Fermi-liquid picture in terms of the Landauer-Büttiker formalism as, in the most simple case, adiabatic transport through the constriction. For a review see Ref. .
Since 1995 several experiments on quantum wires and point contacts have revealed deviations from this integer quantization, $`G=nG_2,n=1,2,3,\mathrm{}`$. In particular the 0.7 conductance anomaly, noted for the first time in 1991 but first studied in detail in 1996, poses one of the most intriguing and challenging puzzles in the field both experimentally and theoretically . This anomaly is a narrow plateau, or in some cases just a shoulder-like feature, clearly visible at the low density side of the first conductance plateau in the dependence of the conductance $`G`$ on a gate voltage which tunes the width and the electron density of the QPC. For low bias voltage the conductance value of the anomalous plateau is around $`0.7G_2`$ giving rise to the name of the phenomenon. The 0.7 anomaly has been recorded in many QPC transport experiments involving different materials, geometries and measurement techniques.
In this paper, we present experimental evidence, that the 0.7 conductance anomaly is associated with a density-dependent energy difference separating two transmission channels. We reach this conclusion by measuring both the temperature and the source-drain bias voltage dependence of the differential conductance, $`G=dI/dV_{sd}`$, through shallow-etched QPCs.
The outline of the paper is as follows. In Sec. II we describe the fabrication of the six samples to be investigated. In the following all detailed results on the conductance of the QPCs are shown solely for sample A, and only towards the end of the paper the main results from all samples are shown. In Sec. III we discuss the lateral confinement potential defining the QPC, and we focus in particular on the fact that this potential is controlled by two independent variables: the gate bias and the source-drain bias. Then follows in Sec. IV the results from finite source-drain bias spectroscopy, and the important energy difference $`\mathrm{\Delta }`$ is introduced. We deal with the temperature dependence of the zero-bias conductance in Sec. V and introduce the activation energy $`T_A`$. The main result is obtained in Sec. VI where we show that $`\mathrm{\Delta }=k_BT_A`$ for all six samples. A short conclusion is given in Sec. VII.
## II The shallow etched samples
The quantum point contacts were all fabricated on modulation doped GaAs/GaAlAs heterostructures grown by molecular beam epitaxy (MBE). The layer sequence is: 1 $`\mu `$m $`\mathrm{GaAs}`$ buffer, $`20\mathrm{nm}`$ $`\mathrm{Ga}_{0.7}\mathrm{Al}_{0.3}\mathrm{As}`$ spacer, $`40\mathrm{nm}`$ $`\mathrm{Ga}_{0.7}\mathrm{Al}_{0.3}\mathrm{As}`$ barrier layer with a Si concentration of $`2\times 10^{24}`$ m<sup>-3</sup>, and a 10 nm undoped GaAs cap layer. The carrier density is $`2\times 10^{15}`$ m<sup>-2</sup> and the mobility is 100 m<sup>2</sup>/Vs, measured in the dark at a temperature of 4.2 K.
The samples were processed with a $`20\times 100`$ ($`\mu `$m)<sup>2</sup> mesa, etched 100 nm, and AuGeNi ohmic contacts to the 2DEG were formed by conventional UV-lithography, lift-off and annealing. The narrow QPC constriction was defined using electron beam lithography (EBL) and shallow wet-etching on the mesa. The following procedure was used: The sample was flushed in acetone, methanol and iso-propanol before it was ashed in an oxygen plasma for 20 seconds. The sample was then pre-etched in 18% HCl for 5 minutes, flushed in H<sub>2</sub>O and blown dry in nitrogen. It was then pre-baked for 5 minutes at 185 C before spinning on a 125 nm thick layer of PMMA electron beam resist. The EBL pattern was exposed with an acceleration voltage of 30 kV, and developed in MIBK:iso-propanol (1:3). The sample was post-baked for 5 minutes at 115 C, and ashed 6 seconds before etching 55-60 nm in $`\mathrm{H}_2\mathrm{O}:\mathrm{H}_2\mathrm{O}_2:\mathrm{H}_3\mathrm{PO}_4`$ $`(38:1:1)`$ at an etch rate of 100 nm/min.
Three types of devices were investigated: top-gated (type I), side-gated (type II), and overgrown side-gated (type III). Fig. 1a shows a scanning electron microscope (SEM) picture of a type I QPC constriction. The shallow etched walls of the constriction are shaped as two back-to back parabolas. The picture was taken before the constriction was covered by a 10 $`\mu `$m wide, 100 nm thick Ti/Au top-gate electrode. In type II devices, Fig. 1b, the QPC constriction is formed by etching two semi-circular trenches, $``$250 nm wide and $``$60 nm deep. The etched trenches also define two large areas of 2DEG which are used as side-gates. The same pattern is used in type III devices, but the trenches are etched 90 nm deep to reach the GaAs/GaAlAs interface and then MBE-regrown. In this way the constriction is bounded by heterostructure-interfaces, both vertically and laterally. The e-beam patterning and the MBE-regrowth was made before the Ohmic contacts were deposited. Before the regrowth, the sample was desorbed at $`630^{}`$C for 2 minutes in the MBE-chamber. The sample was then overgrown with 100 nm undoped $`\mathrm{Ga}_{0.9}\mathrm{Al}_{0.1}\mathrm{As}`$ and a 5 nm undoped GaAs cap layer, using a growth temperature of $`590^{}`$C. The sample parameters are tabulated in Table I.
The samples were mounted in a liquid helium refrigerator, and the differential conductance, $`G=dI_{\mathrm{sd}}/dV_{\mathrm{sd}}`$, was measured with a small ac excitation voltage, 5-50 $`\mu `$V rms, using standard lock-in techniques at 33-117 Hz. The effective width of the QPC and the electron density inside it is controlled by a gate voltage, which is applied between the source contact and the top or side gate electrode. Henceforth this gate voltage is denoted $`V_{\mathrm{gs}}`$.
## III The lateral confinement
The shallow etching technique gives rise to a strong lateral confinement in the constriction. We have previously reported observation of quantized conductance at temperatures above 30 K in a 50 nm wide shallow etched QPC with a 1D-subband energy separation $`\mathrm{\Delta }_{01}20`$ meV . In this paper our main example is sample A (type I), but all the measurements reported for this sample have also been performed for the others. Fig. 2 shows the gate-characteristics, i.e. the differential conductance, $`G`$, as function of gate-source voltage, $`V_{\mathrm{gs}}`$, of sample A, measured at different temperatures. The 200 nm wide, etched QPC constriction is depleted at zero gate-voltage, and a positive gate-source voltage is necessary to open it. We estimate the 1D-subband energy separations in the QPC’s from the thermal smearing of the conductance plateaus, and more precisely by finite bias spectroscopy as described below. For the 200 nm wide QPC constriction in device A we find an energy separation between the two lowest 1D-subbands, $`\mathrm{\Delta }_{01}=6.5`$ meV, see also Table I and Sec. IV B.
The confinement potential $`U`$ determines the transmission properties of the device. It is mainly defined by the sample parameters, the geometry, and the gate-source voltage $`V_{gs}`$. However, to some extend, especially near pinch-off where the electron density is low, it does also depend on the bias voltage $`V_{sd}`$ . In short we write $`U=U(V_{gs},V_{sd})`$. This effect of $`V_{sd}`$ influencing $`U`$ we denote ’self-gating’ since it resembles the ordinary gate effect from $`V_{gs}`$ . A sample exhibiting a self-gating can be said to be ’soft’, if not it is ’rigid’.
The current $`I`$ through the QPC can be expressed in terms of the transmission functions $`𝒯_n(\epsilon )`$ and the difference $`\mathrm{\Delta }f(\epsilon )`$ in thermal occupation factors for the source and drain reservoirs as:
$$I=\frac{2e}{h}\underset{n}{}_{\mathrm{}}^{\mathrm{}}𝑑\epsilon 𝒯_n(\epsilon )\mathrm{\Delta }f(\epsilon ),$$
(1)
where
$`𝒯_n(\epsilon )`$ $`=`$ $`𝒯_n[\epsilon ,U(V_{gs},V_{sd})]`$ (2)
$`\mathrm{\Delta }f(\epsilon )`$ $`=`$ $`f[\epsilon \mu \nu eV_{sd}]f[\epsilon \mu +(1\nu )eV_{sd}],`$ (3)
with $`\nu `$ being a number between 0 and 1 describing the ratio of the potential drop on each side of the constriction. Our experimental results are compatible with $`\nu =1/2`$. Writing explicitly the most relevant functional dependencies for the current we obtain:
$$I=I[U(V_{gs},V_{sd}),\mathrm{\Delta }f(V_{sd})].$$
(4)
From this follows to first order in a Taylor expansion the expressions for the differential conductance $`dI/dV_{sd}`$ and the transconductance $`dI/dV_{gs}`$, the quantities measured in the experiments:
$`{\displaystyle \frac{dI}{dV_{sd}}}`$ $``$ $`{\displaystyle \frac{I}{U}}{\displaystyle \frac{U}{V_{sd}}}+{\displaystyle \frac{I}{\mathrm{\Delta }f}}{\displaystyle \frac{\mathrm{\Delta }f}{V_{sd}}},`$ (5)
$`{\displaystyle \frac{dI}{dV_{gs}}}`$ $``$ $`{\displaystyle \frac{I}{U}}{\displaystyle \frac{U}{V_{gs}}}.`$ (6)
We note that any sharp features in the transconductance reminiscent of the characteristic step-like form in the conductance (see Fig. 2) derives from the factor $`I/U`$ in Eq. (6) relating to the opening of new conductance channels. The other factor $`U/V_{gs}`$ is just varying smoothly due to its origin in electrostatics over length scales of the order of at least 100 nm. But $`I/U`$ also appears as a prefactor in the first term of the differential conductance in Eq. (5). Thus the self-gating effect is enhanced when the transconductance is large. Conversely, at low temperatures at the middle of a plateau the current is almost unaffected by changes in $`U`$, at least only very smooth changes are expected. If $`I/U`$ can be neglected, the differential conductance is given by the occupation factor related second term in Eq. (5). As the temperature is enhanced the transconductance becomes more important even at the center of the plateau as is evident for the highest temperatures in Fig. 7.
## IV Bias spectroscopy and the energy difference $`\mathrm{\Delta }`$
An important source of information about the energy subbands in a QPC is finite bias spectroscopy. We use the technique developed by Patel et al. and described theoretically by Glazman and Khaetskii. The differential conductance, $`G=dI/dV_{\mathrm{sd}}`$, at finite dc source-drain bias voltage, $`V_{\mathrm{sd}}`$ is measured by lock-in technique, using a small ac signal, 50 $`\mu `$V rms 117 Hz, superposed on the dc source-drain bias voltage.
### A The differential conductance at finite bias
In Fig. 3 it is shown how at $`T=0.3`$ K the differential conductance of sample $`A`$ depends on the dc source-drain bias. For each trace the gate voltage is fixed, while going from one trace to the next represents an increase in gate voltage of 1 mV. Conductance plateaus appear as dark regions with a high density of traces. Four types of plateaus are observed in the data. (1) The first four integer conductance plateaus are clearly seen at $`G=nG_2`$ around $`V_{\mathrm{sd}}=0`$. (2) The corresponding half-plateaus at approximately $`(n1/2)G_2`$ appears for bias voltages $`2\mathrm{mV}<|V_{sd}|<6\mathrm{mV}`$, when the chemical potential of one reservoir lies above the edge of one subband, while the other potential lies below. (3) We remark that the 0.7 structure is observed observed near $`V_{\mathrm{sd}}=0`$. As the source-drain bias is increased, the $`G`$-value of the conductance anomaly increases, and for $`|V_{\mathrm{sd}}|1`$ mV, the anomaly has evolved into a well-defined plateau with a conductance $`G`$ between 0.8 and 0.9$`G_2`$. (4) Finally, an additional plateau feature is observed at $`G1.4G_2`$ for $`6\mathrm{mV}<|V_{sd}|<8\mathrm{mV}`$.
From the data in Fig. 3 it is seen that the differential conductance depends rather strongly on $`V_{sd}`$. For the lowest conductances a pronounced asymmetry is observed: for negative $`V_{sd}`$ the conductance is higher than for positive $`V_{sd}`$. This effect is always seen when the gate bias is applied relative to the source contact. It persist in all samples even for different grounding points. Furthermore, even at the smallest source-drain bias we observe a strong non-linearity in the conductance at the middle of the integer plateaus, where the chemical potentials lie in the middle of the gap between 1d subband edges: the integer plateaus in Fig. 3 are not flat around $`V_{sd}=0`$. In the following we interpret this non-linearity and the asymmetry in terms of the self-gating effect presented in Sec. III. We subtract this trivial effect from the data to obtain data corresponding to a ’rigid’ QPC not subject to self-gating.
First we treat the asymmetry of the data, which is most strong for the lowest values of $`V_{gs}`$ or equivalently for the lowest electron densities. A simple reason for this can be found in the electrostatics of the QPC. We notice that $`\frac{I}{U}`$ is always antisymmetric with respect to $`V_{sd}`$. However, since the gate voltage is applied relative to the source contact, no special symmetry relations are expected in $`\frac{U}{V_{sd}}`$ as the polarity $`V_{sd}`$ is changed. Especially near pinch-off when the electron density is low in the QPC the effect of a polarity change in $`V_{sd}`$ can be important. Thus we expect on general grounds that regarded as a function of $`V_{sd}`$ the term $`\frac{I}{U}\frac{U}{V_{sd}}`$ from Eq. (5) contains both a symmetric and an antisymmetric part. This conclusion holds true for any value of the ratio $`\nu `$ of the voltage drop in Eq. 3 in contrast to Ref. , where $`\nu 1/2`$ had to be adopted to explain the asymmetry. The antisymmetric part thus attributed to rather trivial electrostatics is subtracted from the data by forming the symmetric combination
$$I(|V_{sd}|)\frac{1}{2}[I(+V_{sd})+I(V_{sd})].$$
(7)
Next we focus on the four $`dI/dV_{sd}`$ traces which for $`V_{sd}=0`$ goes right through the center of each of the first four integer conductance plateaus. As mentioned in Sec. III no appreciable self-gating effect is expected here. Only smooth changes with $`V_{gs}`$ is expected for moderate values of the bias $`V_{sd}`$. Using a second order Taylor expansion of $`dI/dV_{sd}`$ in $`V_{sd}`$ we extend Eq. (5) to the form
$$\frac{dI}{dV_{sd}}(\alpha V_{gs}+\beta )+(\alpha ^{}V_{gs}+\beta ^{})V_{sd},$$
(8)
and fit the four parameters $`\alpha `$, $`\beta `$, $`\alpha ^{}`$, and $`\beta ^{}`$ to the four mid-plateau traces. We then subtract from all the traces the fitted $`V_{sd}`$ dependence. The result of this procedure is shown in Fig. 4. We end up with plots of the integer plateaus in the differential conductance which for moderate values of $`V_{sd}`$ up to 2-3 mV are independent of the finite bias voltage. Note how also the 0.9 anomalous plateau has now become flat. We can thus unambiguously assign constant values for the conductance plateaus in a wide range. The half-plateaus, however, still show a dependence of the bias voltage, although not as strongly as before, indicating the large influence of $`V_{sd}`$ on the potential $`U`$ in the strong non-equilibrium case where one reservoir is injecting electrons above the topmost subband edges and the other not. We note that experimentally we never see $`G=0.5`$ at the first half-plateau but rather a value substantially below and never quite constant but decreasing with increasing bias; in the present case $`G0.3`$. This is probably due to the intricate self-consistent electrostatic effects at pinch-off, but this have to be investigated further. The measured values of the conductance at the plateaus are discussed further in Sec. IV C.
### B The transconductance
To display the features in the conductance traces more clearly we study the transconductance, $`dG/dV_{\mathrm{gs}}`$, which is calculated by numerical differentiation from the measured differential conductance $`G=dI/dV_{sd}`$. The transconductance is zero (or small) on conductance plateaus and shows peaks in the transition regions between plateaus. In Fig. 5 is shown a grayscale plot of the transconductance of sample A, calculated from the data in Fig. 3. The plot covers the range $`10`$ to 10 mV in source-drain bias and 0.25 to 0.50 V in gate voltage corresponding to the first four integer conductance plateaus. Plateau regions (small transconductance) appear as light regions bounded by dark transition regions (high transconductance). The main feature of the plot is the well-known diamond shaped dark transition regions surrounding the integer plateaus $`nG_2`$ and the half-plateaus $`(n1/2)G_2`$, where $`n=1,2,\mathrm{}`$ . The transitions in $`G`$ are due to the crossing of the chemical potentials $`\mu _s`$ and $`\mu _d`$ of the source and drain reservoirs through the subband edges defining the transmitting subbands. The procedure described in Sec. IV A to get rid of the $`V_{sd}`$ dependence of the plateau values allows for an unambiguous assignment of conductance values in each of the diamonds of the transconductance plot. The subband separation $`\mathrm{\Delta }_{01}`$ is extracted from the main diamond structure by reading off the value of $`V_{sd}`$ where the straight black lines surrounding the 1 diamond intersect indicating the appearance of the next subband. The intersection is at $`(V_{gs},V_{sd})=(0.32\mathrm{V},\mathrm{\hspace{0.25em}6.5}\mathrm{mV})`$. Thus $`\mathrm{\Delta }_{01}=6.5`$ mV as listed in Table I.
### C The anomalous subband edge $`\mathrm{\Delta }(V_{gs})`$
In addition to the main feature the anomalous conductance plateaus are seen. The most pronounced is the anomalous $`G=0.9`$ plateau, which appears in the left-hand side of the $`G=1`$ diamond between the leftmost black straight edge and a curved gray anomalous transition line. Note how the anomalous transition line is continued smoothly into the $`G=1.5`$ diamond. Similar, but much weaker, anomalous structures are seen running inside the 2 diamond continuing into the 2.5 diamond, and inside the 3 diamond continuing into 3.5 diamond.
Just as the black straight lines in the grayscale plot of Fig. 5 are due to the crossing of $`\mu _s`$ and $`\mu _d`$ through the subband edges of the transmitting subbands, it is tempting to also associate a subband edge crossing with the anomalous transitions. In particular the strong
transition ridge between the 1.0 and the 0.9 plateau can be analyzed in those terms. In the standard theory changing $`V_{sd}`$ for fixed $`V_{gs}`$ at the first half of the first plateau leads to the sequence $`G=1.0G=0.5`$, since $`\mu _d`$ drops below the lowest lying spin-degenerate subband edge. However, this sequence is not observed in the measurements. To make this point clear we show in Fig. 6 four individual traces at fixed $`V_{gs}`$, denoted A to D, and four traces at fixed $`V_{sd}`$, denoted E to H. In Fig. 6a these traces are drawn as dashed lines in the $`V_{gs}`$-$`V_{sd}`$ plane. In Fig. 6b is shown the differential conductance along trace A to D. The zero-bias point of these four traces corresponds to the following positions on the $`T=0.3`$ K conductance curve of Fig. 2: below the first plateau (A), on the lower half of the first plateau (B), on the upper half of the first plateau (C), and on the lower half of the second plateau (D). First follow trace B. It exhibits the plateau sequence $`G=1.0G=0.85G=0.2`$. Probably due to the ’softness’ of the QPC at low electron densities the value of the ’0.5-plateau’ is around 0.2, where the trace meet with trace A evolving from $`G=0`$ into a plateau at $`G=0.15`$. It is as if the conductance in trace B drops in two steps corresponding to the crossing of two subband edges rather than just one, perhaps as a consequence of lifting of the spin-degeneracy in the QPC.
It seems quite natural to associate the anomalous transition with an anomalous subband edge which lies above the ordinary subband edge and therefore is encountered first as the bias voltage is raised. This would also account for the continuation of the anomalous transition into the 1.5 diamond as seen by studying the behavior of trace C. Increasing $`V_{sd}`$ from 0 this trace exhibits a clear plateau at 1.0 before it rises and develops into a plateau at $`G=1.45`$ as $`\mu _s`$ is raised above the second subband. For slightly larger value of $`V_{sd}`$ $`\mu _d`$ falls below the anomalous subband edge; $`G`$ drops and the trace exhibits a shoulder-like feature around $`G=1.3`$. Only for yet higher values of $`V_{sd}`$ does $`\mu _d`$ drop below the ordinary first subband leading to $`G=1`$ and lower values as in the standard case. Thus as a function of the bias-voltage $`V_{sd}`$ the plateau sequences $`G=1.0G=0.5`$ and $`G=1.0G=1.5G=1.0`$, for the first and second half of the $`G=1`$-plateau, expected from the simple half-plateau model, in experiment are seen rather to be $`G=1.0G=0.85G=0.5`$ and $`G=1.0G=1.5G=1.3G=1.0`$.
The values of the conductance at the plateaus are found after the fitting procedure described in Fig. 4. The most precise way to obtain these values is through Fig. 6c, where the transconductance $`dG/dV_{gs}`$ is plotted versus the differential conductance $`G`$ at four different but fixed bias voltages, traces E, F, G and H. The plateaus appear as minima in the curves, since a minimum in the the transconductance correspond to the point of least slope in plots of $`G`$ versus $`V_{gs}`$. Ideally, if the plateaus are completely flat, the values at the minima are 0. This happens for example at the integer plateaus seen in trace E, and the half-plateaus in trace G. The 0.85-plateau is never completely flat, but in traces F and G it is seen as a well developed minimum.
For comparison with the temperature data presented in Sec.VI we introduce the anomalous gate voltage dependent (and hence density dependent) energy difference $`\mathrm{\Delta }(V_{gs})`$. It is related to that particular gate-voltage dependent value $`V_{sd}^{}`$ of the source-drain bias that maximizes the transconductance along the 0.9-1.0 and 1.35-1.5 ridges in the grayscale plot:
$$\mathrm{\Delta }(V_{gs})=\frac{1}{2}eV_{sd}^{}(V_{gs}),$$
(9)
In terms of an anomalous subband, $`\mathrm{\Delta }`$ is interpreted as the difference between the chemical potential and the anomalous subband edge. In Fig. 5 it is seen that similar ridges appear, progressively weaker, for the higher subbands. The weakening of the effect may be due to less pronounced spin polarization at the higher densities present when more subbands are occupied . Finally, we note that in contrast to the normal plateaus, the anomalous plateaus only appear when both $`\mu _s`$ and $`\mu _d`$ are above a given subband edge: the anomalous plateaus only appear in the left-hand side of the diamonds in the grayscale plots. This is another indication that the anomalous plateaus are related to interaction effects and not simple single-particle subband effects.
## V The activation temperature $`T_a`$
To gain further insight in the conductance anomaly we also study the temperature development of the first conductance plateau, $`G=G_2`$. In Fig. 7 is shown a set of measurements performed on sample A. At the lowest temperature, 0.3 K, the plateau is broad and flat. With a 1D subband energy separation of 6.5 meV, the thermal smearing of the plateau should be negligible at temperatures below 4 K. This is indeed also the case for the upper half of the conductance plateau, $`V_{gs}320340`$ mV, which stays flat as the temperature is raised. On the lower half of the plateau, the conductance is suppressed below the plateau value, $`G_2`$, as the temperature is raised, developing a plateau-like structure around the conductance value $`0.7G_2`$. This is the 0.7 conductance anomaly.
The large 1D-subband energy separation in the shallow etched QPC’s allows us to study the temperature dependence of the $`0.7`$-structure at temperatures up to around 5 K without appreciable thermal smearing of the quantized conductance. In Fig. 8 we present two Arrhenius plots of the conductance suppression shown in Fig. 7 at $`V_{gs}=0.305`$ V and 0.309 V. We plot the relative conductance suppression $`1G(T)/G_0`$ (where $`G_0`$ is the measured conductance value of the plateau) versus $`1/T`$ at the given fixed gate-voltage. The linear behavior in the semilogarithmic Arrhenius plot indicate an activated behavior, $`G(T)/G_0=1C\mathrm{exp}(T_A/T)`$, with the corresponding activation temperatures, $`T_A=0.28`$ K and 1.11 K, extracted from the two slopes, respectively. Fig. 8b shows how the measured activation temperature $`T_A`$ as a function of gate voltage increases from 0 at pinch off to a few kelvin at the middle of the conductance plateau.
In the usual framework of the Landauer-Büttiker formalism the observed activated suppression of the conductance indicates that the $`0.7`$-structure is associated with thermal depopulation of a subband having a gate voltage dependent subband edge. If a subband edge lies $`k_BT_A(V_{gs})`$ below the Fermi level indeed an activated behavior is seen in $`G`$. A phenomenological theory along these lines has been presented by Bruus et al. . Moreover, this picture is in accordance with the discussion presented in Sec. IV C of the crossing of subband edges at finite bias. In the following analysis we connect the measured activation temperature with the energy gap $`\mathrm{\Delta }`$ found by finite-bias spectroscopy.
## VI Comparing $`\mathrm{\Delta }`$ and $`T_a`$
It is possible to ascribe the same origin to the appearance of the plateau at $`0.9G_2`$ at finite bias as to the 0.7 anomaly. The two effects are connected by the energies $`\mathrm{\Delta }(V_{gs})`$ and $`T_A(V_{gs})`$. Consider a fixed gate-voltage on the lower half of the $`G_2`$ plateau. The data are taken at low temperature. At zero bias the excitation energies available for the electrons at the Fermi energy are not sufficient to reach the subband edge lying $`k_BT_A`$ below the Fermi level, and the conductance has the expected, quantized value, $`G_2`$. As the source-drain bias-voltage, $`V_{\mathrm{sd}}=(\mu _s\mu _d)/e`$ is increased we assume that half the potential drop is before and the other half after the QPC, i.e. $`\nu =1/2`$ in Eq. (3). The electrons from the drain reservoir are injected below the subband edge when $`eV_{\mathrm{sd}}/2=k_BT_A`$. This assumption is supported by our experiments. In Fig. 9 we have for all six samples plotted the expected position of the resonance, $`V_{\mathrm{sd}}^{}=2k_BT_A/e`$, versus gate voltage as white circles. The activation temperature used in this plot is obtained from the measured temperature dependence of the $`0.7`$-structure, as the one presented in Figs. 7 and 8. As seen from Fig. 9 the transition from the regular $`G_2`$ plateau to the anomalous $`0.85G_2`$ plateau appears at the expected resonance position.
The quality of the 0.7/0.85 anomalies are varying a lot from sample to sample. The exact reason for this is not known at present. One can think of many reasons such as impurities, geometry related defects, and other sample parameters. But it is noteworthy that for all samples the energy $`\mathrm{\Delta }(V_{gs})`$ characterizing the 0.85 anomaly coincides with the activation energy $`k_BT_A`$ deduced from the 0.7 anomaly.
## VII Conclusion
We have investigated the $`0.7`$ conductance anomaly in six samples of three different types of shallow etched GaAs/GaAlAS QPC’s: top-gated, side-gated and side-gated, overgrown. We note that the QPC confinement potential $`U`$ depends on both $`V_{gs}`$ and $`V_{sd}`$. The influence from $`V_{sd}`$, referred to as self-gating, can explain the distinct asymmetry and non-linearity always observed in differential conductance of QPCs. The QPCs thus appear to be ’soft’, but we have shown how to subtract the self-gating effect from the data. Based on finite bias spectroscopy we have presented experimental evidence, that the 0.7 anomaly is associated with a density-dependent energy difference $`\mathrm{\Delta }`$ of the order of a few kelvin being the distance from the chemical potential to an anomalous subband edge. The shallow etching technique gives rise to a strong lateral confinement with 1D subband energy separations of $`520`$ meV. We have therefore been able study the 0.7 anomaly for higher temperatures than for normal split-gate devices, and this allowed a detailed study of the temperature dependence of the conductance anomaly. We have found an activated behavior of the conductance suppression on the $`0.7`$ anomaly, with a density-dependent activation temperature, $`T_A`$, of a few kelvin. For all six samples the energy difference $`\mathrm{\Delta }`$ is found to be equal to the activation energy $`k_BT_A`$. Our observations supports the idea that the 0.7/0.85 conductance anomaly arises from the existence of an anomalous subband edge in the QPC. The nature of the anomalous subbands is presently unknown. But our observation that the anomalous plateaus only appear when both $`\mu _s`$ and $`\mu _d`$ are above a given subband edge, and the behavior of the 0.7 anomaly as function of magnetic field indicate the importance of interaction effects beyond the simple single-particle subband picture, presumably related to spin polarization .
## Acknowledgements
This research is part of the EU IT-LTR programme Q-SWITCH (No. 20960/30960), and was partly supported by the Danish Technical Research Council (grant no. 9701490) and by the Danish Natural Science Research Council (grants no. 9502937, 9600548 and 9601677). The III-V materials used in this investigation were made at the III-V NANOLAB, operated jointly by the Microelectronics Centre of the Danish Technical University and the Niels Bohr Institute fAPG, University of Copenhagen.
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# Dwarf satellite galaxies in the modified dynamics
## 1 Introduction
The dynamical behavior of dwarf spheroidals and other satellites of the Milky Way holds much information pertinent to the dark-matter problem. Attempts to elicit such knowledge include, on the one hand, measurements of the satellites’ intrinsic properties such as the size, luminosity, and velocity dispersion, which evince mass discrepancies in the satellites (Aaronson $`\&`$ Olszewski (1988); Pryor (1991); Mateo (1998)). This discrepancy is removed in the modified dynamics–MOND (Milgrom (1995); McGaugh $`\&`$ de Blok (1998); Mateo (1998)). On the other hand, the satellites can be used to probe the gravitational field of their mother galaxy (specifically, the Milky Way) by using them as test particles to probe the galaxy’s potential field (e.g., Little $`\&`$ Tremaine (1987)), or by studying tidal effects of the galaxy on the structure of the satellite taken as a finite body (e.g., Faber $`\&`$ Lin (1983)). In Newtonian dynamics, the history of the center-of-mass motion may influence the internal workings of the satellite via tidal effects. Tidal disruption may have culled from the satellite population those that are internally weakly bound and/or move on elongated orbits, thus affecting the distribution of galactic orbits seen today (see, e.g., Lynden-Bell, Cannon, $`\&`$ Godwin (1983)).
In MOND, the interaction between the internal and center-of-mass motions, brought about by the theory’s nonlinearity, goes beyond the Newtonian effects. For small systems (smaller than the scale over which the external field varies) the effect goes in one direction: while the center-of-mass motion is not affected by the internal motions, it may strongly affect them as explained in Bekenstein $`\&`$ Milgrom (1984) and Milgrom (1986). This occurs when the accelerations inside the satellite are of the order of or smaller than its center-of-mass acceleration; it is also required that the internal accelerations be small compared with the acceleration constant of MOND, $`a_0`$, as is always the case for the Milky Way’s dwarf satellites. Due to this external-field effect (EFE) a satellite that plunges into the galaxy on an eccentric orbit increases in size, making itself an easier victim for tidal disruption. An additional destructive effect results when the changes in the external field become resonant with the internal motions.
The purpose of the paper is to describe, and demonstrate the pertinence, of these processes, which are peculiar to MOND.
In the next section we briefly recapitulate the external field effect. Next, in section 3, we consider a dwarf on an elongated orbit, delineating the different regimes of application of the MOND effects, and give an analytic description of dynamics in the adiabatic regime. In section 4 we describe the MOND N-body simulations, the results of which are described in section 5. Section 6 lists our conclusions and briefly comments on the Milky Way’s dwarf satellites.
## 2 The External-Field Effect
We work with the formulation of MOND as modified gravity described in Bekenstein $`\&`$ Milgrom (1984) whereby the Poisson equation for the nonrelativistic gravitational potential is replaced by
$$[\mu (|\varphi |/a_0)\varphi ]=4\pi G\rho $$
(1)
(and the gravitational acceleration is given by $`\varphi `$). For a small system, freely falling in an external field that dominates its own, equation(1) can be linearized in the internal field, expanding about the value of the external acceleration $`g_{ex}`$ (approximately constant over the extent of the small system). As shown in Milgrom (1986) one gets a quasi-Newtonian internal dynamics with an effective gravitational constant $`G_e=G/\mu (g_{ex}/a_0)`$. We use the term “quasi-Newtonian” because, in addition to the increased effective gravitational constant, the dynamics is anisotropic with some dilation along the direction of the external field. If the external field is in the $`z`$ direction then the linearized equation reduces to the Poisson equation in the coordinates $`x,y,z(1+L)^{1/2}`$, where $`Ldln[\mu (s)]/dln(s)`$ at $`s=g_{ex}/a_0`$ takes a value between 0 and 1. In the deep-MOND limit \[$`s1`$, where $`\mu (s)s`$\], assumed all along in this paper, we then have
$$G_e=Ga_0/g_{ex}G,$$
(2)
and $`L1`$.
## 3 A satellite on an elongated orbit
MOND’s basic premise is that our galaxy, like others, does not contain dynamically important dark matter. Thus, as long as the orbit of the satellite under consideration does not take it within a few scale lengths of the mother galaxy, the latter may be treated as a point mass. Even within galaxies the mean accelerations never much exceed $`a_0`$; at large galactocentric distances the acceleration is always smaller than $`a_0`$, as we assume all along. The MOND-limit acceleration at a distance $`R`$ from a point mass $`M`$ is
$$g(R)=V_{\mathrm{}}^2/R,$$
(3)
where $`V_{\mathrm{}}(GMa_0)^{1/4}`$ is the asymptotic, circular-orbit speed. This $`g(R)`$ is the external acceleration field $`g_{ex}`$ that enters the quasi-Newtonian internal dynamics of the dwarf on the sections of its orbit where $`g(R)`$ outweighs the internal accelerations. If along some portion of the dwarf’s orbit the change in the external field is slow–i.e., occurs on time scales long compared with the internal dynamical time–the quantity $`vr`$ is expected to remain constant as an adiabatic invariant. Here $`v`$ is some mean internal velocity, and $`r`$ is the mean radius of the system. Also, in the quasi-Newtonian regime an effective, Newtonian virial relation should hold: $`v^2G_eMr^1`$. As $`R`$ varies along this section of the orbit, and with it $`G_e`$, we expect $`v`$ and $`r`$ to follow according to $`vG_eR`$, and $`rG_e^1R^1`$. As the dwarf plunges in on an eccentric orbit it puffs up–an effect that does not appear in Newtonian dynamics with dark matter–thus rendering itself more susceptible to tidal breakup than it would be due to the increasing external-field gradients alone.
To consider more quantitatively the interplay between adiabaticity, external-field dominance, and tidal breakup, consider a satellite described by its gross properties: the (baryonic) mass $`m`$, the root-mean-square velocity of the constituents with respect to the center-of-mass $`v`$, and the size $`r`$ (say the rms distance of constituents from the center). It moves on an orbit $`R(t)`$ with velocity $`V(t)`$ in the field of the point-like mother galaxy of mass $`M`$. (We neglect the secondary effects of anisotropy and so we only consider the magnitude of the position vector $`\stackrel{}{R}`$.)
The parameter
$$\beta v^2/rg(R)=v^2R/V_{\mathrm{}}^2r$$
(4)
measures the importance of the internal acceleration vis-a-vis the external one. The EFE is pertinent when $`\beta 1`$. In terms of the radii and masses we can write
$$\beta \{\begin{array}{cc}(R/r)(m/M)^{1/2}\hfill & \text{if }\beta 1\text{;}\hfill \\ (R/r)^2(m/M)\hfill & \text{if }\beta 1\text{.}\hfill \end{array}$$
(5)
Here we used Newtonian expressions with $`G_e`$ from eq.(2) when $`\beta 1`$. The parameter
$$\gamma (R/V_{\mathrm{}})/(r/v)=(R/r)^{1/2}\beta ^{1/2}$$
(6)
is useful for measuring the degree of adiabaticity (achieved when $`\gamma 1`$) when the orbit is mildly eccentric, because then the orbital changes occur on a time scale $`R/V`$, and $`VV_{\mathrm{}}`$. (The MOND potential far from a central mass is logarithmic, for which the virial relation reads $`V^2=V_{\mathrm{}}^2`$. The velocities at perigalacticon, $`V_p`$, and apogalacticon, $`V_a`$, are related by $`V_p^2V_a^2=V_{\mathrm{}}^2ln(R_a/R_p)`$, where $`R_p`$ and $`R_a`$ are the respective distances.) We can write
$$\gamma \{\begin{array}{cc}(R/r)(m/M)^{1/4}\hfill & \text{if }\beta 1\text{;}\hfill \\ (R/r)^{3/2}(m/M)^{1/2}\hfill & \text{if }\beta 1\text{.}\hfill \end{array}$$
(7)
Because only $`Rr`$ is of interest, we see from eq.(6) that $`\gamma 1`$ in the whole region $`\beta >1`$.
Tidal effects in the bulk are important when the mean internal acceleration, $`g_{in}`$, is smaller than the increment of the external acceleration over the extent $`r`$; i.e., when $`g_{in}V_{\mathrm{}}^2r/R^2=g(R)r/R`$. We take as the criterion for the importance of tidal effects
$$\alpha [g_{in}/(V_{\mathrm{}}^2r/R^2)]^{1/3}=(vR/V_{\mathrm{}}r)^{2/3}1.$$
(8)
Again, since only cases for which $`rR`$ are of interest we see that tidal effects need concern us only when $`g_{in}g`$ ($`\beta 1`$). In this regime we have $`\alpha (R/r)(m/M)^{1/3}`$. Note in general that $`\alpha =\gamma ^{2/3}`$. This means that non-adiabaticity and tidal effects enter at about the same place on the orbit, as is indeed verified in our numerical calculations. Clearly, the inflation of the satellite due to the EFE continues in the tidal phase.
We can now qualitatively see what happens to a dwarf on an elongated orbit. If the entire orbit has $`\beta 1`$ the satellite is unaffected by the Galaxy. If its orbit takes it to a small enough galactocentric distance $`R_0`$ where $`\beta =1`$ the EFE enters into action there. At this point the situation is adiabatic with $`\gamma _0=\gamma (R_0)(R_0/r_0)^{1/2}`$, which is $`10`$ for the typical value of $`R_0/r_0100`$. As $`R`$ decreases further we are, at first, in the adiabatic regime with $`rr_0R_0/R`$, $`vv_0R/R_0`$, and $`\beta `$ decreasing still below 1: $`\beta (R/R_0)^4`$. In this region $`\gamma \gamma _0(R/R_0)^3`$ from eq.(7), so roughly at $`R=R_0\gamma _0^{1/3}`$ adiabaticity is lost, and at the same time tidal effects set in, in which case we have to resort to MOND, N-body calculations, as described below.
The comparison of the MOND predictions on the onset of tidal effects with those of Newtonian dynamics (ND) (with dark matter) depends on what exactly is measured, and on what is assumed in ND (e.g., on the dark-matter distribution in the dwarf). But, in any event, the puffing up of a dwarf in the $`\beta 1`$ region, which has no analogue in ND, makes dwarfs more vulnerable to tidal disruption. To take a specific example, suppose a satellite is observed at $`R=R_1`$ with measured size, internal, and center-of-mass velocities. Its future orbit can then be deduced, and also the Newtonian, dynamical mass it contains, $`m_N`$. Suppose it is already in the $`\beta 1`$ regime. It is easy to see that the value of its tidal parameter $`\alpha `$ as deduced in ND is the same as that given by MOND, since $`m_N=mG_e(R_1)/G`$ and the galactic mass within $`R_1`$ is $`M_N(R_1)=MG_e(R_1)/G`$. As we saw, MOND predicts that $`\alpha R^2`$, while Newtonian dynamics predicts $`\alpha R^{2/3}`$, since $`r`$ is then assumed to remain constant while $`M_N(R)R`$. So tidal effects will clearly enter at larger radii in MOND.
## 4 N-body simulations
The numerical simulations involve a model dwarf comprising N identical particles that starts with some equilibrium distribution function in compliance with MOND dynamics. The model is then subjected to different types of variable external influences that mimic aspects of the influence of the mother galaxy. The underlying potential field equation is eq.(1), or its linearized, approximate form. This nonlinear potential equation is solved numerically using multi-grid methods as detailed in Brada (1996) and adumbrated in Brada $`\&`$ Milgrom (1999). The particles are then propagated in the derived potential. It is only interesting to study the dwarf when it is in the external-field-dominated region. To isolate the different effects discussed above we proceed in three steps. First, to pinpoint the effects of non-adiabaticity, and verify our analytic deductions for the adiabatic regime, we start with a quasi-Newtonian King model for the dwarf, assume quasi-Newtonian dynamics, and simply vary $`G_e`$ periodically and see how the model reacts for different frequencies of the perturbation. In the second step we still consider a quasi-Newtonian behavior but the applied variations in $`G_e`$ and the direction of the external field correspond to actual orbits of a dwarf. The third stage, which is more costly, is to simulate the complete system of dwarf plus a point-mass galaxy. Since the construction of initial models for the dwarf are peculiar to MOND we describe them briefly now, referring the reader for more details to Brada (1996).
### 4.1 Constructing steady-state galactic models
We use as initial states King models (For details see, e.g., Binney $`\&`$ Tremaine (1987)) properly modified to constitute, as the case may require, quasi-Newtonian or deep-MOND steady states. The distribution function for the Newtonian models is
$$f_K=\{\begin{array}{cc}\rho _1(2\pi \sigma ^2)^{3/2}(e^{\epsilon /\sigma ^2}1)\hfill & \text{if }\epsilon >0\hfill \\ 0\hfill & \text{if }\epsilon 0\text{,}\hfill \end{array}$$
(9)
where $`\epsilon E+\varphi _0`$, $`E=v^2/2+\varphi `$, and the parameter $`\varphi _0`$ is the upper energy cutoff. Equation(9) is integrated over velocities to obtain the density $`\rho _K(\mathrm{\Psi })`$ as a function of the relative potential $`\mathrm{\Psi }\varphi +\varphi _0`$. Instead of the Poisson equation we solve here the spherically symmetric MOND equation (in the deep-MOND limit assumed all along):
$$[a_0^1r^2(\mathrm{\Psi }^{})^2]^{}=4\pi G\rho _K(\mathrm{\Psi })$$
(10)
(the apostrophe signifies derivative with respect to $`r`$), which provides an ordinary differential equation for $`\mathrm{\Psi }(r)`$ that can be integrated numerically with the boundary condition $`\mathrm{\Psi }^{}(0)=0`$. The second boundary condition is $`\mathrm{\Psi }(0)`$, which together with $`\varphi _0`$ determines the model. The model can also be specified in terms of other parameters from among the tidal radius $`r_t`$, the total mass, the central density $`\rho (0)`$, and the King radius $`r_0(9\sigma ^2/4\pi G\rho _0)^{1/2}`$. Note that for MOND King models $`r_0`$ defined in this way is not some characteristic radius; it is just a convenient representation of $`\rho _0`$ (so we can have $`r_0>r_t`$, for example).
In constructing a quasi-Newtonian model we remember that the transformed potential $`\varphi ^{}(x^{},y^{},z^{})\varphi (x,y,z)`$, with $`x^{}=x,y^{}=y,z^{}=z(1+L)^{1/2}`$, with $`L=dln[\mu (s)]/dln(s)`$, satisfies the usual Poisson equation with $`G_e`$ as gravitational constant, and $`\rho ^{}(x^{},y^{},z^{})\rho (x,y,z)`$ as density ($`z`$ is taken in the direction of the external field). We thus begin by constructing a Newtonian model in the auxiliary coordinates $`x^{},x^{},z^{}`$ remembering that the conserved quantity on which the distribution function depends by the Jeans theorem is $`[v_x^{}_{}{}^{}2+v_y^{}_{}{}^{}2+(1+L)v_z^{}_{}{}^{}2]/2+\varphi ^{}(x^{},y^{},z^{})`$. This model has a spherical mass distribution, but a $`v_z^{}`$ dispersion that is smaller by a factor $`(1+L)^{1/2}`$ than those in the other directions. We draw positions and velocities for the auxiliary coordinates of the N particles. Then we multiply all $`z^{}`$ and $`v_z^{}`$ values by $`(1+L)^{1/2}`$. The total mass of the model is multiplied by the same factor because $`\rho ^{}(\stackrel{}{r}^{})d^3r^{}=(1+L)^{1/2}\rho (\stackrel{}{r})d^3r`$. The resulting model is elongated in the $`z`$ direction and has an isotropic global velocity distribution.
## 5 Results of the simulations
### 5.1 Testing for the external-field effect
We first want to establish the consequences of the external-field effect and learn what is the time scale necessary for a change to be adiabatic. As was discussed in section 3, in the adiabatic regime we expect the average velocity in a system, $`v`$, to be proportional to $`G_e=G(a_0/g_{ex})`$, and the average size of the system to be inversely proportional to $`G_e`$. We start with a Newtonian King model having $`10^5`$ particles with $`\sigma ^2=1`$, $`r_0=1`$, and $`\mathrm{\Psi }(0)/\sigma ^2=1`$. We also take $`G=1`$ \[this fixes $`\rho (0)`$; the total mass of the Newtonian model is then $`m=0.72`$, and its tidal radius $`r_t=1.975`$\]. We then perform the stretching by $`2^{1/2}`$ (we take $`L=1`$ for the deep-MOND limit) to get a model with $`m=1.02`$ and root-mean-square values of the coordinate and velocity components (designated by capital letters) $`X=Y=0.44`$, $`Z=0.62`$, $`V_x=V_y=V_z=0.39`$. The natural dynamical time $`r/v`$ is of order unity. We ran simulations on a cubical grids with $`129^3`$ grid points using the quasi-Newtonian field equation with a fixed time step $`dt=0.04`$ for $`10^3`$ time steps. We varied $`G_e`$ periodically with time taking $`G_e(t)=(0.7+0.3\mathrm{cos}\omega t)^1`$ for three values of $`\omega =2\pi (1/40,1/20,1/10)`$. ( The physical grid spacing is changed in proportion to $`G_e^1`$ during the simulation, since we expect $`r`$ to scale as $`G_e^1`$.) In analogy with the adiabaticity parameter, $`\gamma `$, defined locally on an orbit, we can define here as some global measure of adiabaticity $`\widehat{\gamma }(T/4)/[X(0)/V_x(0)]`$, where $`T=2\pi /\omega `$ is the period. For the three models presented $`\widehat{\gamma }=10,5,2.5`$ respectively. The results of these simulations for the $`x`$ components are shown in figure (1) describing how the extent and velocity dispersion in the $`x`$ direction vary with time. The results for the $`z`$ direction are the same within a few percents. The results for $`\omega =2\pi /40`$ show the expected external-field effect with strict adiabaticity; small departures from adiabaticity appear in the time dependence of the size and mean velocity for $`\omega =2\pi /20`$; while for $`\omega =2\pi /10`$ clear departure from adiabaticity is evident.
We also learn that departure from adiabaticity brings about a secular increase in the radius and decrease in the internal velocities.
### 5.2 The evolution of a quasi-Newtonian model along a realistic orbit
Before going to the more complete models that utilize the full MOND field equation and include the tidal forces, we follow the evolution of a quasi-Newtonian model varying the value of $`G_e`$ and the direction of the external field according to the location of the dwarf on an actual orbit in the logarithmic potential of the point-mass mother galaxy. Tidal forces are then not taken into account. The purpose of these experiments is to test the degree to which the changes in dwarf characteristics are adiabatic for sample models with realistic parameters. In particular, to see what lasting effects non-adiabaticity near perigalacticon has on the dwarf in disjunction from tidal effects. We thus took orbits with strong adiabaticity at apogalacticon, where we start, but a breakdown of adiabaticity near perigalacticon. Note that $`G_e`$ goes back to the same value at subsequent apogalacticons; so, apart from some remaining oscillations on the dynamical time scale, the virial relation is reestablished in the mean, and, thus, $`v^2r`$ must come back to the same value.
We take the mass of the mother galaxy as unit, $`M=1`$. Since we also use units in which $`G=1`$ and $`a_0=1`$, the unit of velocities becomes $`V_{\mathrm{}}(MGa_0)^{1/4}=1`$, which in cgs units is about $`220kms^1`$ for the MW. Length is then measured in units of $`V_{\mathrm{}}^2/a_0`$, about 10 kpc for the MW. In these units the MW satellites are typically at distances between 5 and 20, of size $`3\times 10^2`$ to $`10^1`$, velocity dispersion $`2\times 10^2`$ to $`5\times 10^2`$, and of mass $`10^5`$ to $`10^4`$. Accordingly, we construct our dwarf model as a quasi-Newtonian King model having the following properties: $`\sigma ^2=8\times 10^4`$, $`r_0=1/16`$, $`\mathrm{\Psi }(0)/\sigma ^2=1`$, $`r_t=0.1234`$ (before stretching), and $`m=5.09\times 10^5`$. We simulated its evolution for two orbits with pericenter and apocenter distances of $`R_{min}=6`$, $`R_{max}=9`$, and $`R_{min}=6`$, $`R_{max}=12`$. At galactocentric distance $`R`$ the external acceleration is $`R^1`$ (in our units of $`a_0`$), and $`G_e=R`$.
The results of the simulations are summarized in figures (2)(3). Adiabaticity is more severely violated in the second model, which has $`R_{max}=12`$. We see that, as a result of violating adiabaticity near $`R=R_{min}`$, the dwarf still oscillates on the dynamical time scale when it next enters the adiabatic regime around apocenter. More importantly, it attains a larger radius (averaged over the fluctuations). The velocity dispersion is correspondingly smaller ($`v^2r`$ is preserved and $`vr`$ is larger). On the next close passage the dwarf will be even less adiabatic and becomes more vulnerable and will continue to increase in size.
### 5.3 The evolution of a full MOND model along a realistic orbit.
We then followed the full evolution of a dwarf obeying MOND orbiting a point mother galaxy. We started by producing an isolated MOND King model with the following parameters: $`\sigma ^2=8.53\times 10^3`$, $`r_0=0.3`$ (remember that $`r_0`$ is just a proxy for $`\rho _0`$, not a characteristic radius), and $`\mathrm{\Psi }(0)/\sigma ^2=1`$. The resulting model has a total mass $`m=3.59\times 10^5`$ and a tidal radius $`r_t=0.1346`$. Note that $`\sigma `$ is not the mean velocity dispersion of the MOND model. This can be gotten from the deep-MOND virial relation (Milgrom (1994)), which in our units reads $`v^2=2m^{1/2}/3`$, where $`v^2`$ is the three-dimensional rms velocity. So we get for the one-dimensional rms velocity $`v_x^2=1.33\times 10^3`$.
These global parameters are similar to the ones of the quasi-Newtonian model we have used in the previous subsection, but the density profiles of the two models differ: the MOND model is less concentrated than the quasi-Newtonian model. Our model dwarf is put on an eccentric orbit starting at an apogalacticon distance of $`R_{max}=12`$ and reaching a perigalacticon distance of $`R_{min}=6`$. From the above model we construct a family of five models by scaling up the mass of the model ending up with masses $`m`$, $`2m`$, $`3m`$, $`4m`$, and $`16m`$, and scaling the velocity dispersions up accordingly. These are all models for isolated MOND dwarfs. However, at apocenter already we have to start with models under some external-field influence. So, before we let the models evolve along the orbit we need to switch on adiabatically the external field. This is done in a preliminary simulation where we gradually increase the mass of the mother galaxy from zero to one.
The presence of the mother galaxy enters these simulations through the boundary conditions used by the potential solver. The $`16m`$ model hardly changes when the external field is switched on, while for the $`m`$ model the rms radius, $`r`$, increases by as much as $`50\%`$ when the external field is switched on. The values for the parameter $`\beta `$, which measures the ratio of the internal field to the external field, calculated at $`R=12`$ are $`0.233`$, $`0.7`$, $`1.03`$, $`1.29`$, and $`3.09`$, respectively. Thus, the $`m`$ model is dominated by the external field while the $`3m`$ is the borderline case. We then integrate the internal dynamics and the center-of-mass motion of each of the models on the orbit. The orbits lie in the $`X`$,$`Y`$ plane and we start at $`Y=12`$, $`X=0`$. The simulations lasted for 40 time units and consisted of $`10^3`$ individual time steps. The rms value of $`r`$ and $`v`$ as functions of time are given in figure (4). Also shown there are the adiabaticity ($`\gamma `$) and the tidal ($`\alpha `$) parameters along the orbit. Projections of the dwarf structure for the four smaller-mass models, on the $`XY`$ plane, are given in figure (5).
The value of the parameter $`\beta `$ at pericenter for the five models was: $`0.0315`$, $`0.056`$, $`0.168`$, $`0.2106`$, and $`1.36`$, respectively. The three models with masses $`m`$, $`2m`$ and $`3m`$ show clear signs of tidal disruption. The $`m`$ model seems to have been totally destroyed by the tidal forces and there is no clear core that remained after the passage near the galaxy. The $`2m`$ model was strongly influenced by the tidal forces and lost about $`25\%`$ of its mass. The $`3m`$ model lost only a few percent of its mass through tidal interaction. We can attribute these mass losses to the combined action of tidal forces and the extra non-adiabatic expansion of the models near $`R=R_{min}`$.
In comparison, Newtonian dynamics applied to dwarfs observed with the same initial positions, center-of-mass velocities, sizes, and velocity dispersions would predict much less tidal disruption. Consider, in particular, the two most vulnerable models with masses $`m`$ and $`2m`$. They start at $`R=12`$ with $`\beta <1`$, so, from the discussion at the end of section 3 we see that their Newtonian $`\alpha `$ values there are the same as the MOND values. (Of course, a Newtonist will assume that they contain more mass for the same sizes and velocity dispersions.) From figure 4 we see that the two models start with $`\alpha 4,7`$ respectively. The MOND scaling ($`\alpha R^2`$ for $`\beta 1`$) implies that the two models should have $`\alpha 1,1.7`$ at perigalacticon, $`R=6`$, as they approximately do. In Newtonian dynamics, where $`\alpha R^{2/3}`$, we would get at perigalacticon $`\alpha 2.5,4.4`$ for the two models, making these initial model dwarfs much safer from later tidal disruption. It need perhaps be clarified that the Newtonist will continue to get the same $`\alpha `$ values as in MOND if he uses at every point the observed properties, but this will lead him to conclude that the mass of the dwarf varies. Here we speak of what the Newtonist’s predictions will be given only the initial data, and assuming that the dwarf mass is constant.
## 6 Summary and conclusions
We have studied the existence, the nature, and the influence on dwarf satellites of the external field effect in MOND. For dwarf parameters in the EFE regime two situation are grossly distinguished: a) the adiabatic regime, in which tidal effects are not so important and b) the impulsive region, which also roughly coincides with the region where tidal forces become important. Due to the EFE the radius of a dwarf in the adiabatic regime increases as it approaches the mother galaxy. If the whole orbit is in the adiabatic regime, the structure of the dwarf simply changes periodically with the orbital period. If, however, some segment of the orbit is in the impulsive-tidal regime near pericenter, then the dwarf might lose much of its mass there. Even if it does not, it can emerge from this region having a larger radius and smaller velocity dispersion (hence, a longer intrinsic dynamical time). In its next approach to perigalacticon it will thus enter the impulsive-tidal regime at a larger distance from the center.
Clearly, all the above is highly germane to the dwarf system of the MW. The distribution of intrinsic and orbital parameters of presently observed dwarfs must have been greatly affected by interaction with the MW. And, one expects, MOND would give a different answer than Newtonian dynamics with dark matter. To actually deduce the present-day properties of the dwarfs would, however, require knowledge of the initial distribution of the orbital and intrinsic parameters of the dwarf-satellite population. Nothing is really known about this, so we refrain from speculating on the subject. We only estimate where our dwarf satellites stand as regards external-field dominance, adiabaticity, and the importance of tidal effects.
We consider the 10 dwarf spheroidal satellites with known parameters (Mateo (1998)): Sculptor, LSG 3, Fornax, Carina, Leo I, Sextans, Leo II, Ursa Minor, Draco, and Sagittarius. We take for the MW $`V_{\mathrm{}}=220`$ kms<sup>-1</sup>. Since only core radii, $`r_c`$, are given we write for the mean radius $`r=\eta r_c`$ to get for the adiabaticity parameter of those dwarfs $`\gamma \eta ^1(22,150,14,15,47,8,39,14,20,2)`$, respectively. So, with the exception of Sagittarius–which is known to be in the throes of disruption–and perhaps Sextans, these dwarfs are in the adiabatic regime within reasonable margins for $`\eta `$, and even considering the approximate nature of the $`\gamma `$ criterion. According to our analysis they are also only weakly affected by tidal forces at their present positions. As has been pointed out (Milgrom (1995), McGaugh $`\&`$ de Blok (1998)), most of the above dwarfs (with the exception of LSG 3, Leo I, and Leo II) are materially affected by the EFE: with the above choices of system parameters we get $`\beta \eta ^1(0.7,4.4,0.7,0.5,1.9,0.2,1.2,0.6,0.9,0.1)`$.
If we apply the MOND scaling $`\alpha R^2`$, which is valid in the $`\beta 1`$ regime, to the dwarfs with $`\eta \beta <1`$ (except for Sagittarius) we can estimate the minimum galactocentric distance above which the bulk of the dwarf is immune to tidal effects. This is given by $`R_t^MR_0\alpha _0^{1/2}=R_0\gamma _0^{1/3}`$, where here a subscript 0 marks present-day values. (If a dwarf in now on an outgoing section of its orbit it will return to the same $`R`$, as it goes in, in the same state.) For Sculptor, Fornax, Carina, Sextans, UMi, and Draco we get, respectively $`R_t^M\eta ^{1/3}(28,57,41,43,27,32)`$ kpc. The corresponding Newtonian values ($`R_t^NR_0\gamma _0^1`$) are $`R_t^N\eta (4,10,7,11,5,4)`$ kpc. They are smaller than the corresponding MOND values if $`\eta `$ is not so large that $`\alpha _0<1`$. (For some dwarfs these Newtonian radii may fall within the stellar MW where our approximation of a spherical, logarithmic potential is not valid.)
Our results imply that for a given dwarf in the adiabatic regime on an elongated orbit under a strong EFE the size and velocity dispersion would be strongly dependent on the distance from the mother galaxy. One might then try to look for such correlations in the time-frozen population as seen today. This seems to us quite hopeless at present because the effects will be swamped by other factors of which we know very little; in particular, the unknown distribution of initial (intrinsic and orbital) parameters for the dwarfs. This is aggravated by the small sample size.
We leave for a future publication some other interesting effects predicted by MOND that result from the EFE. For example, in a dwarf in the EFE regime the total angular momentum is not conserved. We alluded to the fact that the direction of the external field if felt by the “internal” dynamics of the dwarf. In a static or adiabatic situation only the angular momentum along the external-field direction is conserved.
We thank the referee, Tad Pryor, for many useful comments and suggestions
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# The Algebraic Model for scattering in three-s-cluster systems. I. Theoretical Background
## I Introduction
Since 1980, the so-called Algebraic Model (AM) of the Resonating Group Method has been used in the investigation of bound and continuum states of nuclear systems. Initially the AM was applied to binary cluster configurations . Later on it was extended to describe binary clusters coupled to collective (quadrupole and monopole) channels . Quite recently, three-cluster configurations were considered in the AM framework . Such configurations play a significant role in light nuclei, in particular for reactions of astrophysical interest.
The Algebraic Model represents the nuclear many-particle wave functions through their expansions in harmonic oscillator eigenstates. The use of a basis of square integrable states reduces the Schrödinger equation to a matrix equation. This procedure is well-known for bound states, but is also applicable to continuum states when the appropriate boundary conditions are imposed on the expansion coefficients . Thus the Algebraic Model provides a unified approach to bound and continuous spectra based on familiar matrix techniques.
Solving an AM scattering problem requires two major steps. The first is to define the basis states relevant to the scattering channels being investigated and to compute the Hamiltonian matrix elements in that basis. The second is to consider the description of the asymptotic region and boundary conditions in the basis and to solve the matrix equation subject to those boundary conditions.
We elaborate on both steps in this paper for a fully antisymmetrized three-cluster system, and introduce the proper three-cluster continuum boundary conditions. Such a description is still lacking in the literature. Three-cluster systems have been considered in some calculations, but there the continuum was approximated through sets of two-cluster configurations .
In our calculation the Pauli principle will be treated rigorously by taking full antisymmetrization into account for all nucleons in the system. The basis states will be defined by a composition of three s-clusters with “frozen” internal structure. The relative motions of the clusters are the relevant degrees of freedom. They are described by the Hyperspherical Harmonics method (HHM). It makes a natural formulation of the proper boundary conditions possible. The matrix elements of kinetic and potential energy in the three-cluster basis states are calculated using the Generating Function technique , and methods for an explicit derivation will be presented.
The asymptotic behavior of relative motion of clusters is obtained by considering large inter-cluster distances where the interaction between clusters approaches zero. We use the folding approximation to determine the asymptotic wave-functions. The folding approximation amounts to solving the scattering equations without nucleon-nucleon interaction and without inter-cluster antisymmetrization.
In the case of neutral clusters, i.e. when inter-cluster Coulomb interactions are absent, the asymptotic equations are uncoupled in the hyperangular momentum $`K`$ that is typical in the HHM. We will therefore introduce individual scattering channels characterized by quantum number $`K`$. In configurations with charged clusters the Coulomb interaction also couples $`K`$-channels in the asymptotic region. Again the folding approximation will be used to arrive at the appropriate asymptotic wave equations. A multiple channel approach which couples all $`K`$-channels has to be used to solve the full AM system of equations. This extends the theoretical foundation of the AM given in for two-cluster systems.
In part II of this paper, the techniques discussed here will be applied to the three cluster configurations $`\alpha +n+n`$ for <sup>6</sup>He, and $`\alpha +p+p`$ for <sup>6</sup>Be.
## II Three-s-cluster systems
### A Wave Functions
The many-particle wave function for a three-cluster system of $`A`$ nucleons ($`A=A_1+A_2+A_3`$) can be written, using the anti-symmetrization operator $`𝒜`$, as follows
$$\mathrm{\Psi }(𝐪_1,..,𝐪_{A1})=𝒜\left[\mathrm{\Psi }_1\left(A_1\right)\mathrm{\Psi }_2\left(A_2\right)\mathrm{\Psi }_3\left(A_3\right)\mathrm{\Psi }_R\left(R\right)\right]$$
(1)
where the centre of mass of the $`A`$-nucleon system has been eliminated by the use of Jacobi coordinates $`𝐪_i`$. The cluster wave functions $`\mathrm{\Psi }_i\left(A_i\right)`$
$$\mathrm{\Psi }_i\left(A_i\right)=\mathrm{\Psi }_i(𝐪_1^{\left(i\right)},..,𝐪_{A_i1}^{\left(i\right)})(i=1,2,3)$$
(2)
represent the internal structure of the $`i`$-th cluster, centered around its centre of mass $`𝐑_i`$. In our AM study these cluster functions are fixed and they are Slater determinants of harmonic oscillator (0s)-states, corresponding to the groundstate configuration of the cluster ($`A_i4`$ for all $`i`$). The $`\mathrm{\Psi }_R\left(R\right)`$ wave function
$$\mathrm{\Psi }_R\left(R\right)=\mathrm{\Psi }_R(𝐪_1^{\left(R\right)},𝐪_2^{\left(R\right)})=\mathrm{\Psi }_R(𝐪_1,𝐪_2)$$
(3)
represents the relative motion of the three clusters with respect to one another, and $`𝐪_1`$ and $`𝐪_2`$ represent Jacobi coordinates. In Fig. 1 we indicate an enumeration of possible Jacobi coordinates and their relation to the component clusters.
The state (3) is not limited to any particular type of orbital; on the contrary we will use a complete basis of harmonic oscillator states for the relative motion degrees of freedom. Thus the full $`A`$-particle state cannot be expressed as a single Slater determinant of single particle orbitals.
An important approximation is obtained by breaking the Pauli principle between the individual clusters, but retaining a proper quantum-mechanical description of the clusters, which is described by the wave function
$$\mathrm{\Psi }_F(𝐪_1,..,𝐪_{A1})=\mathrm{\Psi }_1\left(A_1\right)\mathrm{\Psi }_2\left(A_2\right)\mathrm{\Psi }_3\left(A_3\right)\mathrm{\Psi }_R\left(R\right)$$
(4)
Because each cluster wave function is antisymmetric (they are Slater determinants) one is indeed neglecting the inter-cluster anti-symmetrization only. This results in what is known as the “Folding” model. It has the advantage of preserving the identities of the clusters and, if the intra-cluster structure is kept “frozen”, it reduces the many-particle problem to that of the relative motion of the clusters.
The folding approximation will be the natural choice for calculating the asymptotic behavior of the three cluster-system, i.e. the disintegration of the system in the three non-interacting individual clusters. This amounts to the situation that all three clusters are a sufficient distance apart and intra-cluster interactions are no longer in force.
The folding model is however also an acceptable approximation in the interaction region and can serve as comparison to the fully antisymmetrized calculations. In the folding model, the clusters interact through a local, inter-cluster potential called the folding potential. As in the current paper (0s)-determinants $`\mathrm{\Psi }_i\left(A_i\right)`$ are used to describe the internal state of the clusters, the folding potential will be easily calculated and is a sum of three terms
$$V^{(F)}=V^{(F)}(𝐑_{12})+V^{(F)}(𝐑_{23})+V^{(F)}(𝐑_{31})$$
(5)
where each term is simply the integral
$$V^{(F)}(𝐑_{\tau \upsilon })=\underset{iA_\tau }{}\underset{jA_\upsilon }{}𝑑\tau _\tau 𝑑\tau _\upsilon |\mathrm{\Psi }_\tau (A_\tau )|^2V(𝐫_i𝐫_j+𝐑_{\tau \upsilon })|\mathrm{\Psi }_\upsilon (A_\upsilon )|^2$$
(6)
The coordinates $`𝐑_{\tau \upsilon }`$ are associated with the relative position of the clusters
$$𝐑_{\tau \upsilon }=\frac{1}{A_\tau }\underset{iA_\tau }{}𝐫_i\frac{1}{A_\upsilon }\underset{jA_\upsilon }{}𝐫_j$$
(7)
and sum to zero; they are equivalent to the $`𝐪_1,𝐪_2`$ Jacobi coordinates introduced earlier. In this way the folding approximation turns the three-cluster problem into an effective three-particle problem for the relative motion coordinates.
Because the cluster states are fixed and built up of (0s)-orbitals, the problem of labeling the basis states with quantum numbers relates to the inter-cluster wave function only. This holds true whether one uses the full anti-symmetrization or the folding approximation. In a two-cluster case, the set of quantum numbers describing inter-cluster motion is unambiguously defined. In a three-cluster case, several schemes can be used to classify the inter-cluster wave function in the oscillator representation. In three distinct but equivalent schemes were considered. One of these used the quantum numbers provided by the Hyperspherical Harmonics (HH) method (see for instance ). This is the classification that we will adopt. Even within this particular scheme there are several ways to classify the basis states. We shall restrict ourselves to the so-called Zernike-Brinkman basis . This corresponds to the following reduction of the unitary group $`U(6)`$, the symmetry group of the three-particle oscillator Hamiltonian,
$$U(6)O(6)O(3)O(3)O(3)$$
(8)
This reduction provides the quantum numbers $`K`$, the hypermomentum, $`n`$, the hyperradial excitation, $`l_1`$, the angular momentum connected with the first Jacobi vector, $`l_2`$, the angular momentum connected with the second Jacobi vector, and $`L`$ and $`M`$ the total angular momentum obtained from coupling the partial angular momenta $`l_1`$, $`l_2`$. Collectively these quantum numbers will be denoted by $`\nu `$, i.e. $`\nu =\{n,K,(l_1l_2)LM\}`$ in the remainder of the text.
There are a number of relations and constraints on these quantum numbers:
* the total angular momentum is the vector sum of the partial angular momenta $`𝐥_1`$and $`𝐥_2`$, i.e. $`𝐋=`$ $`𝐥_1+`$ $`𝐥_2`$ or $`\left|l_1l_2\right|Ll_1+l_2`$.
* by fixing the values of $`l_1`$ and $`l_2`$, we impose restrictions on the hypermomentum $`K=l_1+l_2,l_1+l_2+2,l_1+l_2+4,\mathrm{}`$ This condition implies that for certain values of hypermomentum $`K`$ the sum of partial angular momenta $`l_1+l_2`$ cannot exceed $`K`$.
* the partial angular momenta $`l_1`$ and $`l_2`$ define the parity of the three-cluster state by the relation $`\pi =\left(1\right)^{l_1+l_2}`$.
* for the “normal” parity states $`\pi =\left(1\right)^L`$ the minimal value of hypermomentum is $`K_{\mathrm{min}}=L`$, whereas $`K_{\mathrm{min}}=L+1`$ for the so-called “abnormal” parity states $`\pi =\left(1\right)^{L+1}`$.
* oscillator shells with $`N`$ quanta are characterized by the constraint $`N=2n+K`$.
Thus for a given hyperangular and rotational configuration the quantumnumber $`n`$ ladders the oscillator shells of increasing oscillator energy.
## III The Algebraic Model
### A Asymptotic solutions in coordinate representation
The Algebraic Model implements a method to solve the Schrödinger equation for quantum scattering systems, in particular for nuclear cluster systems. It is based on a matrix representation of the Schrödinger equation in terms of a square integrable basis, usually Harmonic Oscillator states, and boundary conditions in terms of the asymptotic behavior of the expansion coefficients of the wave function. In this paper we restrict ourselves to a presentation tailored to the treatment of three-cluster systems.
In the case of three-cluster calculations, one needs to determine a proper approximation for the wave function (3). Consider an expansion of the relative wave function
$$\mathrm{\Psi }_R(𝐪_1,𝐪_2)=\underset{\nu }{}c_\nu \mathrm{\Psi }_\nu (𝐪_1,𝐪_2)$$
(9)
with $`\nu =\{n,K,(l_1l_2)LM\}`$ and $`\left\{\mathrm{\Psi }_\nu \right\}`$ a complete basis of six-dimensional oscillator states. It covers all possible types of relative motion between the three clusters.
To obtain the asymptotic behavior of the three-cluster system, we consider the folding approximation. The assumption that antisymmetrization effects between clusters are absent in the asymtptotic region is a natural one. The relative motion problem of the three clusters in the absence of a potential can then be explicitly solved in the Hyperspherical Harmonics (HH) method (see for instance ). It involves the transformation of the Jacobi coordinates $`𝐪_1`$ and $`𝐪_2`$ to the hyperradius $`\rho `$ and a set of hyperangles $`\mathrm{\Omega }`$. The inter-cluster wave function in coordinate representation is expanded in hyperspherical harmonics $`H_K^{\nu _0}\left(\mathrm{\Omega }\right)`$ where $`\nu _0`$ has been chosen as a shorthand for $`(l_1l_2)LM`$.
In the absence of the Coulomb interaction this leads to a set of equations for the hyperradial asymptotic solutions, with the kinetic energy operator as reference Hamiltonian
$$\left\{\frac{\mathrm{}^2}{2m}\left[\frac{d^2}{d\rho ^2}+\frac{5}{\rho }\frac{d}{d\rho }\frac{K\left(K+4\right)}{\rho ^2}\right]E\right\}R_{K,\nu _0}\left(\rho \right)=0$$
(10)
The solutions can be obtained analytically and are represented by a pair of Hänkel functions for the ingoing and outgoing solutions:
$$R_{K,\nu _0}^{(\pm )}\left(\rho \right)=\left\{\begin{array}{c}H_{K+2}^{\left(1\right)}\left(k\rho \right)/\rho ^2\hfill \\ H_{K+2}^{\left(2\right)}\left(k\rho \right)/\rho ^2\hfill \end{array}\right\}$$
(11)
where
$`k=\sqrt{{\displaystyle \frac{2mE}{\mathrm{}^2}}}`$
One notices that these asymptotic solutions are independent of all quantum numbers $`\nu _0`$, and are determined by the value of hypermomentum $`K`$ only.
When charged clusters are considered the asymptotic reference Hamiltonian consists of the kinetic energy and the Coulomb interaction:
$$\left\{\frac{\mathrm{}^2}{2m}\left[\frac{d^2}{d\rho ^2}+\frac{5}{\rho }\frac{d}{d\rho }\frac{𝒦}{\rho ^2}\right]+\frac{Z_{eff}}{\rho }E\right\}\left(\rho \right)=0$$
(12)
The matrix $`𝒦`$ is diagonal with matrix elements $`K\left(K+4\right)`$, and $`Z_{eff}`$, the “effective charge”, is off-diagonal in $`K`$ and $`(l_1l_2)`$. Different $`K`$-channels are now coupled. A standard approximation for solving these equations is to decouple them by assuming that the off-diagonal matrix-elements of $`Z_{eff}`$ are sufficiently small:
$$\left\{\frac{\mathrm{}^2}{2m}\left[\frac{d^2}{d\rho ^2}+\frac{5}{\rho }\frac{d}{d\rho }\frac{K\left(K+4\right)}{\rho ^2}\right]+\frac{Z_{eff}}{\rho }E\right\}R_{K,\nu _0}\left(\rho \right)=0$$
(13)
The constants $`Z_{eff}`$ depends on $`K`$ and $`\nu _0`$ and all parameters of the many-body system under consideration. We will restrict ourselves to this decoupling approximation, but it is to be understood that its validity has to be checked for any specific three-cluster system.
The asymptotic solutions then become
$$R_{K,\nu _0}^{(\pm )}\left(\rho \right)=\left\{\begin{array}{c}W_{i\eta ,\mu }\left(2ik\rho \right)/\rho ^{\frac{5}{2}}\hfill \\ W_{i\eta ,\mu }\left(2ik\rho \right)/\rho ^{\frac{5}{2}}\hfill \end{array}\right\}$$
(14)
where $`W`$ is the Whittaker function, $`\mu =K+2`$ and $`\eta `$ is the well-known Sommerfeld parameter
$$\eta =\frac{m}{\mathrm{}^2}\frac{Z_{eff}}{k}$$
(15)
As $`\eta `$ is a function of $`K,`$ $`l_1`$ and $`l_2`$ through the parameter $`Z_{eff}`$, the asymptotic solutions will now be dependent on $`K`$ and $`\nu _0`$.
### B Asymptotic solutions in oscillator representation
The Algebraic Model relies on an expansion in terms of oscillator functions, and the asymptotic behavior of the corresponding expansion coefficients $`c_\nu `$. It was conjectured (see for instance ) that for very large values of the oscillator quantum number $`n`$ the expansion coefficients for physically relevant wave-functions behave like
$$c_n=n|\psi \sqrt{2}\rho _n^2\psi (b\rho _n)$$
(16)
where $`\rho _n=\sqrt{4n+2K+6}`$ corresponds to the classical turning point, $`b`$ is the oscillator parameter, and $`\psi `$ is the hyperradial wave function.
In the case of neutral clusters this leads after substitution of the hyperradial asymptotic solutions to the following expansion coefficients $`c_n^{(\pm )}`$
$$c_n^{(\pm )K}\sqrt{2}\left\{\begin{array}{c}H_{K+2}^{\left(1\right)}\left(kb\rho _n\right)\hfill \\ H_{K+2}^{\left(2\right)}\left(kb\rho _n\right)\hfill \end{array}\right\}$$
(17)
This result can be obtained in an alternative way by representing the Schrödinger equation, with the kinetic energy operator $`\widehat{T}`$ as the Hamiltonian to describe the asymptotic situation, in a (hyperradial) oscillator representation
$$\underset{m=0}{\overset{\mathrm{}}{}}n,(K,\nu _0)\left|\widehat{T}E\right|m,(K,\nu _0)c_m^{K,\nu _0}=0$$
(18)
This matrix equation is of a three-diagonal form because of the properties of $`\widehat{T}`$ and the oscillator basis. Solving for the expansion coefficients $`c_n^{K,\nu _0}`$ leads to a three-term recurrence relation
$$T_{n,n1}^{K,\nu _0}c_{n1}^{K,\nu _0}+\left(T_{n,n}^{K,\nu _0}E\right)c_n^{K,\nu _0}+T_{n,n+1}^{K,\nu _0}c_{n+1}^{K,\nu _0}=0$$
(19)
where
$$T_{n,m}^{K,\nu _0}=n,(K,\nu _0)\left|\widehat{T}\right|m,(K,\nu _0)$$
(20)
The asymptotic solutions (i.e. for high $`n`$) of this recurrence relation are then precisely given by (17).
When the Coulomb interaction is present we again apply (16) to obtain
$$c_n^{(\pm )K}\sqrt{2}\left\{\begin{array}{c}W_{i\eta ,\mu }\left(2ikb\rho _n\right)/\sqrt{\rho _n}\hfill \\ W_{i\eta ,\mu }\left(2ikb\rho _n\right)/\sqrt{\rho _n}\hfill \end{array}\right\}$$
(21)
In this case the oscillator representation of the Schrödinger equation is no longer of a tridiagonal form, and cannot be solved analytically for the asymptotic solutions to corroborate this result.
It should be noted, that the above elaborations are valid for relatively small values of momentum $`k`$ and sufficiently large values of discrete hyperradius $`\rho _n`$, or when
$$\frac{k^2}{\rho _n^2}1$$
(22)
which defines the “asymptotic regime”; it shows that for any value of $`k`$ one can find values for $`n`$ where the relation is satisfied.
As we consider an asymptotic decoupling in the $`(K,\nu _0)`$ quantumnumbers, one will deal with asymptotic channels characterized by the $`(K,\nu _0)`$ values. So only in the internal (or interaction) region will states with different $`K`$ and $`\nu _0`$ be coupled by the short-range nuclear potential and the Coulomb potential. The three-cluster system can therefore be described by a coupled-channels approach, where the individual channels are characterized by a single $`K`$-value, and we will henceforth refer to these channels as “$`K`$-channels”.
### C Multi-channel AM equations
In the current many-channel description of the Algebraic Model for three-cluster systems, the channels will be characterized by the a specific value of the set of quantumnumbers $`K,\nu _0`$, whereas the relative motion of clusters within the channel is connected to the oscillator index $`n`$. We will use $`K`$ henceforth as a corporate index for individual channels, and assume it represents all $`K,\nu _0`$ quantum numbers.
The Schrödinger equation can be cast in a matrix equation of the form
$$\underset{K^{}}{}\underset{m}{}n,K\left|\widehat{H}E\right|m,K^{}c_m^K^{}=0$$
(23)
We will now use a representation of the dynamical equations presented in . As we will consider an $`S`$-matrix formulation of the problem, the expansion coefficients are rewritten as
$$c_n^K=c_n^{(0)K}+\delta _{K_iK}c_n^{()K}S_{K_iK}c_n^{(+)K}$$
(24)
where, for the current channel $`K`$, the $`c_n^{(0)K}`$ are the so-called residual coefficients, the $`c_n^{(\pm )K}`$ are the incoming and outgoing asymptotic coefficients (valid for all $`n`$). The matrix element $`S_{K_iK}`$ describes the coupling between the current channel $`K`$ and the entrance channel $`K_i`$.
As shown in , the $`c_n^{(\pm )K}`$ satisfy the following system of equations for a given channel $`K`$
$$\underset{m=0}{\overset{\mathrm{}}{}}n,K\left|\widehat{H}_0E\right|m,Kc_m^{(\pm )K}=\beta _0^{(\pm )K}\delta _{n,0}$$
(25)
$`\widehat{H}_0`$ being the asymptotic reference Hamiltonian, which equals the kinetic energy operator for uncharged clusters, and the kinetic energy operator plus Coulomb interaction for charged clusters. The right-hand side features $`\beta _0^{(\pm )K}`$ which is a regularization factor to account for the irregular behavior of the $`c_0^{(\pm )K}`$. This factor allows one to solve (25) for all values of $`n`$. The value of $`\beta _0^{(\pm )K}`$ can be obtained for both reference Hamiltonians (i.e. with or without Coulomb). The set of equations (25) for the asymptotic coefficients can then be solved numerically to different degrees of approximation depending on the requested precision. The $`c_n^{K(\pm )\text{ }}`$have the desired asymptotic behavior (cfr eqs 17 and 21)
Substitution of (24) in the equations (23) then leads to the following system of dynamical equations for the many-channel system:
$`{\displaystyle \underset{K^{}}{}}{\displaystyle \underset{m}{}}n,K\left|\widehat{H}E\right|m,K^{}c_m^{K^{}\left(0\right)}`$ $`{\displaystyle \underset{K^{}}{}}S_{K_iK^{}}\left[\beta _0^{(+)K^{}}\delta _{n,0}\delta _{K^{}K}+V_n^{KK^{}(+)}\right]`$ (26)
$`=`$ $`\beta _0^{()K}\delta _{n,0}\delta _{K_iK}V_n^{KK_i()}`$ (27)
where the dynamical coefficients $`V_n^{KK^{}(\pm )}`$, defined in , are given by
$$V_n^{(\pm )KK^{}}=\underset{m=0}{\overset{\mathrm{}}{}}n,K\left|\widehat{V}\right|m,K^{}c_m^{(\pm )K^{}}$$
(28)
This system of equations should be solved for both the residual coefficients $`c_n^{K\left(0\right)}`$ and the $`S`$-matrix elements $`S_{K^{}K}`$.
To obtain an appropriate approximation to the exact solution of (27), we consider an internal region corresponding to $`n<N`$ and an asymptotic region with $`nN`$. The choice of $`N`$ is such that one can expect the residual expansion coefficients $`\left\{c_n^{(0)K}\right\}`$ to be sufficiently small in the asymptotic region. Under these assumptions (27) reduces to the following set of $`N+1`$ equations ($`n=0..N`$):
$`{\displaystyle \underset{K^{}}{}}{\displaystyle \underset{m<N}{}}n,K\left|\widehat{H}E\right|m,K^{}c_m^{(0)K^{}}`$ $`{\displaystyle \underset{K^{}}{}}S_{K_iK^{}}\left[\beta _0^{(+)K^{}}\delta _{n,0}\delta _{K^{}K}+V_n^{KK^{}(+)}\right]`$ (29)
$`=\beta _0^{()K}\delta _{n,0}\delta _{K_iK}V_n^{KK_i()}`$ (30)
The total number of equations for a given entrance channel $`K_i`$ amounts to $`N_{ch}\left(N+1\right)`$, and solving the set of equations by traditional numerical linear algebra leads to the $`N_{ch}N`$ residual coefficients $`\{c_n^{K\left(0\right)}\text{}K=K_{\mathrm{min}}..K_{\mathrm{max}}\text{}n=0..N1\}`$ and $`N_{ch}`$ $`S`$-matrix elements $`\{S_{K_iK}\text{}K=K_{\mathrm{min}}..K_{\mathrm{max}}\}`$. The set of equations has to be solved for all $`N_{ch}`$ entrance channels labelled by $`K_i`$.
### D Numerical solution and convergence
The numerical solution of the AM equations crucially depends on a proper choice of $`N`$, distinguishing the internal from the external region. The determining factor in this is the form of the potential energy matrix elements which, contrary to the short-range coordinate character of the potential, can be of a slowly descending nature in $`n`$. If this is the case, a sufficiently large value of $`N`$ has to be chosen.
In the case of three-cluster systems it is known from literature that the potential asymptotically behaves as $`1/\rho ^3`$ in the hyperradius, with a corresponding effect on the matrix elements. It will be shown later on that the asymptotic form of the effective potential in the current case follows this behavior. It is well known that potentials with an asymptotic tail $`1/\rho ^3`$ dramatically change the phase shift behavior in the low energy region and that special care should be taken to get convergent results. This cannot always be obtained by merely choosing a sufficiently large value of $`N`$.
In reference a numerical strategy to account for long tails in the potential in the AM was developed. It dramatically improves the convergence of the results with significantly smaller values for $`N`$. We refer to for further details.
## IV The Generating Function method
### A General principle
In this section, the general principles for calculating matrix elements in a three-cluster basis will be explained. The two main quantities of interest are: the overlap matrix, and the hamiltonian matrix. The former is of importance because of the proper normalization of the basis states. The latter is decomposed into the kinetic energy operator, the matrix elements of which are obtained mainly by group-theoretical considerations, the potential energy operator, which in our case will be chosen to be a semi-realistic two-body interaction based on a superposition of Gaussians, and the Coulomb contribution.
In this work matrix elements for two-body Gaussian interactions will be derived. From these matrix elements of other functional forms of two-body interactions can be obtained using Gaussian transforms. This latter procedure will be followed to calculate the matrix elements of the Coulomb interaction.
The basic principle of generating functions is well-known from mathematical physics. A generating function or generator state depends on a parameter, referred to as the generating coordinate, in such a way that an expansion with respect to that parameter yields basis states as expansion terms. A familiar example are the single-particle translated Gaussian wave functions
$$\varphi (𝐫|𝐑)=\mathrm{exp}\left\{\frac{1}{2}𝐫^2+\sqrt{2}𝐑𝐫\frac{1}{2}𝐑^2\right\}$$
(31)
with the translation parameter $`𝐑`$ acting as generator coordinate. The choice of parametrization of the generator coordinate influences the quantum numbers of the individual basis states that are generated. In a Cartesian parametrization $`𝐑=(R_x,R_y,R_z)`$ one generates the familiar Cartesian $`\varphi _{n_x}(R_x)\varphi _{n_y}(R_y)\varphi _{n_z}(R_z)`$ oscillator states. With a radial parametrization $`𝐑=R\stackrel{ˇ}{𝐑}`$ (where the inverted hat stands for a unit vector) the expansion yields
$$\varphi (𝐫|𝐑)=\underset{n,l,m}{}𝒩_{nl}R^{2n+l}Y_{lm}(\stackrel{ˇ}{𝐑})\varphi _{nlm}(𝐫)$$
(32)
An underlying mathematical connection exists between such expansions, group representation theory and coherent state analysis . In the present work we exploit the generating function principle to facilitate the computation of matrix elements. The matrix element of any operator between generating states is a function of the generating coordinates on the left and right
$$X(𝐑,𝐑^{})=\varphi \left(𝐫|𝐑\right)\left|\widehat{𝐗}\right|\varphi \left(𝐫|𝐑^{}\right)$$
(33)
Expansion of this function will yield the matrix elements between the basis states. They can be identified in the expansion by the appropriate dependence on the generator coordinates
$$X(𝐑,𝐑^{})=\underset{nlm}{}\underset{n^{}l^{}m^{}}{}𝒩_{nl}𝒩_{n^{}l^{}}R^{2n+l}R^{(2n^{}+l^{})}Y_{lm}^{}(\stackrel{ˇ}{𝐑})Y_{l^{}m^{}}(\stackrel{ˇ}{𝐑}^{^{}})\varphi _{nlm}(𝐫)\left|\widehat{𝐗}\right|\varphi _{n^{}l^{}m^{}}(𝐫)$$
(34)
Of course, one is not required to expand with respect to all parameters at once. Elimination of the angular dependence first, yields a partial generating function for the radial matrix elements:
$`X(𝐑,𝐑^{})`$ $`=`$ $`{\displaystyle \underset{lm}{}}{\displaystyle \underset{l^{}m^{}}{}}Y_{lm}^{}(\stackrel{ˇ}{𝐑})Y_{l^{}m^{}}(\stackrel{ˇ}{𝐑}^{^{}})X_{lm;l^{}m^{}}(R,R^{})`$ (35)
$`X_{lm;l^{}m^{}}(R,R^{})`$ $`=`$ $`{\displaystyle \underset{n}{}}{\displaystyle \underset{n^{}}{}}𝒩_{nl}𝒩_{n^{}l^{}}R^{2n+l}R^{(2n^{}+l^{})}\varphi _{nlm}(𝐫)\left|\widehat{𝐗}\right|\varphi _{n^{}l^{}m^{}}(𝐫)`$ (36)
Such partial generating functions will prove to be particularly useful when we apply the generating function method to the three-cluster problem with six generator coordinates (corresponding to six degrees of freedom) and an extensive set $`\nu =\{n,K,(l_1l_2)LM\}`$ of quantum numbers.
The calculation of the matrix elements with the generating function method is a two-step process. The first step is the calculation of the generating function for the operator involved. Usually this is accomplished with analytical techniques. The second step is the expansion of the generating function w.r.t. the generator coordinates. Several approaches have been used in this respect. Explicit differentiation is one of them. Using recurrence relations for the expansion terms is another one . In any case, the work involved here is straightforward but extremely tedious; both approaches are best implemented using algebraic manipulation software such as Mathematica or Maple.
In this paper we introduce a representation of the generating functions in a manageable form to obtain explicit matrix elements and their connecting recurrence relations.
### B Three-cluster generator state
The customary generator state for the inter-cluster basis functions is given by (in what follows we shall use small $`𝐪`$ for the Jacobi vectors and capital $`𝐐`$ for the corresponding generating coordinates)
$$\mathrm{\Psi }(𝐪_\mathrm{𝟏},𝐪_\mathrm{𝟐}|𝐐_\mathrm{𝟏},𝐐_\mathrm{𝟐})=\mathrm{exp}\left\{\frac{1}{2}\left(𝐪_1^2+𝐪_\mathrm{𝟐}^\mathrm{𝟐}\right)+\sqrt{2}\left(𝐐_1𝐪_1+𝐐_2𝐪_2\right)\frac{1}{2}\left(𝐐_1^2+𝐐_2^2\right)\right\}$$
(37)
The choice of parametrization is linked to the basis states one intends to generate. Associated with our choice of basis (Zernike-Brinkman ), we introduce hyperspherical coordinates. The hyperradius and hyperangles, both for spatial coordinates and for generating parameters, are defined by:
$`\rho =\sqrt{𝐪_1^2+𝐪_2^2},`$ $`q_1=\rho \mathrm{cos}\theta ,`$ $`q_2=\rho \mathrm{sin}\theta ;`$ (38)
$`R=\sqrt{𝐐_1^2+𝐐_2^2},`$ $`Q_1=R\mathrm{cos}\mathrm{\Theta },`$ $`Q_2=R\mathrm{sin}\mathrm{\Theta },`$ (39)
Using these, one expands the generating function (37) in hyperspherical harmonic functions:
$$\mathrm{\Psi }(𝐪_\mathrm{𝟏},𝐪_\mathrm{𝟐}|𝐐_\mathrm{𝟏},𝐐_\mathrm{𝟐})=\underset{\nu }{}\mathrm{\Psi }_\nu (\rho ,\theta ,\stackrel{ˇ}{𝐪}_1,\stackrel{ˇ}{𝐪}_2)\mathrm{\Xi }_\nu ^{}(R,\mathrm{\Theta },\stackrel{ˇ}{𝐐}_1,\stackrel{ˇ}{𝐐}_2)$$
(40)
where the full set of quantum numbers $`\nu `$ (introduced previously) is involved in the summation. The oscillator basis functions are
$$\mathrm{\Psi }_\nu (\rho ,\theta ,\stackrel{ˇ}{𝐪}_1,\stackrel{ˇ}{𝐪}_2)=𝒩_{n,K}\rho ^K\mathrm{exp}\{\rho ^2/2\}L_n^{K+2}(\rho ^2)H_K^{(l_1l_2)LM}(\theta ,\stackrel{ˇ}{𝐪}_1,\stackrel{ˇ}{𝐪}_2)$$
(41)
and the generator coordinate functions are
$$\mathrm{\Xi }_\nu (R,\mathrm{\Theta },\stackrel{ˇ}{𝐐}_1,\stackrel{ˇ}{𝐐}_2)=𝒩_{n,K}R^{K+2n}H_K^{(l_1l_2)LM}(\mathrm{\Theta },\stackrel{ˇ}{𝐐}_1,\stackrel{ˇ}{𝐐}_2)$$
(42)
Here $`H`$ denotes the hyperspherical harmonic function
$`H_K^{(l_1l_2)LM}(\mathrm{\Theta },\stackrel{ˇ}{𝐐}_1,\stackrel{ˇ}{𝐐}_2)`$ $`=`$ $`𝒩_K^{(l_1l_2)LM}\mathrm{\Phi }_K^{(l_1l_2)}\left(\mathrm{\Theta }\right)\left\{Y_{l_1}(\stackrel{ˇ}{𝐐}_1)\times Y_{l_2}(\stackrel{ˇ}{𝐐}_2)\right\}_{LM}`$ (43)
$`\mathrm{\Phi }_K^{(l_1l_2)}\left(\mathrm{\Theta }\right)`$ $`=`$ $`\left(\mathrm{cos}\mathrm{\Theta }\right)^{l_1}\left(\mathrm{sin}\mathrm{\Theta }\right)^{l_2}P_{\frac{Kl_1l_2}{2}}^{l_2+\frac{1}{2},l_1+\frac{1}{2}}(\mathrm{cos}2\mathrm{\Theta })`$ (44)
From (42), one easily deduces the procedure for selecting basis functions with fixed quantum numbers $`\nu =\{n,K,\left(l_1l_2\right)LM\}`$. One has to differentiate the generating function $`\left(K+2n\right)`$-times with respect to $`R`$ and then to set $`R=0`$. After that one has to integrate over $`\mathrm{\Theta }`$ with the weight $`\mathrm{\Phi }_K^{(l_1l_2)}`$ to project onto the hypermomentum $`K`$; one has to integrate over unit vectors $`\stackrel{ˇ}{𝐐}_1`$ and $`\stackrel{ˇ}{𝐐}_2`$ with weights $`Y_{l_1m_1}(\stackrel{ˇ}{𝐐}_1)`$ and $`Y_{l_2m_2}(\stackrel{ˇ}{𝐐}_2)`$ to project onto partial angular momenta. The order of these operations is not important and is a matter of convenience for each specific case. However these calculations, in particular those connected with integrating over hyperangle $`\mathrm{\Theta }`$, are extremely extensive and cumbersome. For this reason, we introduce a new generating function appropriate for three-cluster calculations.
We start from the function below which depends on seven generating coordinates namely:
$`\mathrm{\Psi }(𝐪_\mathrm{𝟏},𝐪_\mathrm{𝟐}|ϵ,𝐐_\mathrm{𝟏},𝐐_\mathrm{𝟐})`$ (45)
$`=`$ $`(1+ϵ)^3\mathrm{exp}\left\{{\displaystyle \frac{1}{2}}{\displaystyle \frac{1ϵ}{1+ϵ}}\left(𝐪_1^2+𝐪_\mathrm{𝟐}^\mathrm{𝟐}\right)+{\displaystyle \frac{\sqrt{2}}{1+ϵ}}\left(𝐐_1𝐪_1+𝐐_2𝐪_2\right){\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{1+ϵ}}\left(𝐐_1^2+𝐐_2^2\right)\right\}`$ (46)
It was used previously in a different context to describe the coupling between monopole and two-cluster degrees of freedom. In those cases, the parameter $`𝐐`$ generates basisfunctions of inter-cluster motion while parameter $`ϵ`$ generates collective monopole excitations of the $`A`$-nucleon system. Here, we will modify the function somewhat and use it only for the inter-cluster motion. We exploit the redundancy in the set of generating parameters (seven parameters vs. six degrees of freedom) and the fact that all expressions up to now are valid for complex generator coordinates also. We restrict the moduli of $`𝐐_1`$ and $`𝐐_2`$ and set
$$Q_1=S,Q_2=iS$$
(47)
When the complex conjugate version of (45) is used, e.g. in the calculation of matrix elements, (47) is also complex conjugated.
We now consider the new set of generator coordinates $`ϵ,S,\stackrel{ˇ}{𝐐}_1,\stackrel{ˇ}{𝐐}_\mathrm{𝟐}`$, and substitute (47) in (45), where the inverted hats on $`𝐐_1`$ and $`𝐐_2`$ again indicate the angular components of both variables. This leads to
$`\mathrm{\Psi }\left(𝐪|ϵ,S,\stackrel{ˇ}{𝐐}_1,\stackrel{ˇ}{𝐐}_2\right)`$ $`=`$ $`(1+ϵ)^3\mathrm{exp}\left\{{\displaystyle \frac{1}{2}}{\displaystyle \frac{1ϵ}{1+ϵ}}\left(𝐪_1^2+𝐪_\mathrm{𝟐}^\mathrm{𝟐}\right)+{\displaystyle \frac{\sqrt{2}S}{1+ϵ}}\left(\stackrel{ˇ}{𝐐}_1𝐪_1i\stackrel{ˇ}{𝐐}_2𝐪_2\right)\right\}`$ (48)
$`=`$ $`{\displaystyle \underset{\nu }{}}\mathrm{\Psi }_\nu (\rho ,\theta ,\stackrel{ˇ}{𝐪}_1,\stackrel{ˇ}{𝐪}_2)\mathrm{\Phi }_\nu (ϵ,S,\stackrel{ˇ}{𝐐}_1,\stackrel{ˇ}{𝐐}_2)`$ (49)
where the weights associated with each basis function are given by
$$\mathrm{\Phi }_\nu (ϵ,S,\stackrel{ˇ}{𝐐}_1,\stackrel{ˇ}{𝐐}_2)=𝒩_K^{(l_1l_2)LM}𝒩_{n,K}\left(i\right)^{l_2}ϵ^nS^K\left\{Y_{l_1}(\stackrel{ˇ}{𝐐}_1)\times Y_{l_2}(\stackrel{ˇ}{𝐐}_2)\right\}_{LM}$$
(50)
These are of a simpler structure and easier to use than (42) because through (47) the dependence on the hyperangular coordinate has been eliminated.
The full generating function for the matrix elements $`X_{\nu ,\nu ^{}}`$ of operator $`\widehat{𝐗}`$ now has the following general structure:
$`X(ϵ,S,\stackrel{ˇ}{𝐐}_1,\stackrel{ˇ}{𝐐}_\mathrm{𝟐};ϵ^{},S^{},\stackrel{ˇ}{𝐐}_1^{},\stackrel{ˇ}{𝐐}_\mathrm{𝟐}^{})`$ (51)
$`=`$ $`{\displaystyle \underset{\nu ,\nu ^{}}{}}\mathrm{\Psi }_\nu (\rho ,\theta ,\stackrel{ˇ}{𝐪}_1,\stackrel{ˇ}{𝐪}_2)\left|\widehat{𝐗}\right|\mathrm{\Psi }_\nu ^{}(\rho ,\theta ,\stackrel{ˇ}{𝐪}_1,\stackrel{ˇ}{𝐪}_2)\mathrm{\Phi }_\nu ^{}(ϵ,S,\stackrel{ˇ}{𝐐}_1,\stackrel{ˇ}{𝐐}_2)\mathrm{\Phi }_\nu ^{}(ϵ^{},S^{},\stackrel{ˇ}{𝐐}_1^{},\stackrel{ˇ}{𝐐}_\mathrm{𝟐}^{})`$ (52)
$`=`$ $`{\displaystyle \underset{\nu ,\nu ^{}}{}}X_{\nu ,\nu ^{}}𝒩_K^{(l_1l_2)LM}𝒩_K^{}^{(l_1^{}l_2^{})L^{}M^{}}𝒩_{n,K}𝒩_{n^{},K^{}}\left(i\right)^{l_2}\left(i\right)^{l_2^{}}ϵ^nϵ^n^{}S^KS^K^{}`$ (54)
$`\left\{Y_{l_1}(\stackrel{ˇ}{𝐐}_1)\times Y_{l_2}(\stackrel{ˇ}{𝐐}_2)\right\}_{LM}\left\{Y_{l_1^{}}(\stackrel{ˇ}{𝐐}_1^{})\times Y_{l_2^{}}(\stackrel{ˇ}{𝐐}_2^{})\right\}_{L^{}M^{}}`$
where again the shorthand notation for the quantum numbers $`\nu =(n,K,(l_1l_2)LM)`$ and $`\nu ^{}=(n^{},K^{},(l_1^{}l_2^{})L^{}M^{})`$ is used. As explained before, we will also consider partial generating functions which have been reduced with respect to a subset of generating coordinates and their corresponding quantum numbers. Most often, we will use a reduction with respect to the angular momentum dependence i.e.
$`X(ϵ,S,\stackrel{ˇ}{𝐐}_1,\stackrel{ˇ}{𝐐}_\mathrm{𝟐};ϵ^{},S^{},\stackrel{ˇ}{𝐐}_1^{},\stackrel{ˇ}{𝐐}_\mathrm{𝟐}^{})`$ (55)
$`=`$ $`{\displaystyle \underset{\left(l_1l_2\right)LM,(l_1^{}l_2^{})L^{}M^{}}{}}𝒳_{\left(l_1l_2\right)LM;(l_1^{}l_2^{})L^{}M^{}}(ϵ,S,ϵ^{},S^{})`$ (57)
$`\left\{Y_{l_1}(\stackrel{ˇ}{𝐐}_1)\times Y_{l_2}(\stackrel{ˇ}{𝐐}_2)\right\}_{LM}\left\{Y_{l_1^{}}(\stackrel{ˇ}{𝐐}_1^{})\times Y_{l_2^{}}(\stackrel{ˇ}{𝐐}_2^{})\right\}_{L^{}M^{}}`$
with the partial generating function
$`𝒳_{\left(l_1l_2\right)LM;(l_1^{}l_2^{})L^{}M^{}}(ϵ,S,ϵ^{},S^{})=`$ (58)
$`{\displaystyle \underset{n,K,n^{},K^{}}{}}X_{K,n,\left(l_1l_2\right)LM;K,n,(l_1^{}l_2^{})L^{}M^{}}𝒩_K^{(l_1l_2)LM}𝒩_K^{}^{(l_1^{}l_2^{})L^{}M^{}}𝒩_{n,K}𝒩_{n^{},K^{}}\left(i\right)^{l_2}\left(i\right)^{l_2^{}}ϵ^nϵ^n^{}S^KS^K^{}`$ (59)
generating matrix elements for specified $`\left(l_1l_2\right)LM`$ and $`(l_1^{}l_2^{})L^{}M^{}`$ only.
The asymmetry in the treatment of the different quantum numbers is motivated by the methodology in which the matrix elements will be used. Indeed, all quantities in this paper are calculated in the context of the Algebraic Model. As explained earlier the spatial asymptotic behavior is mapped onto the asymptotic behavior of the expansion coefficients in the oscillator basis. As for fixed $`K`$ the $`n`$ quantum numbers ladder through the oscillator shells, they will be needed for sufficiently high values in order to properly describe the asymptotic region.
## V Matrix elements in the folding approximation
In this section we derive the generating functions for the overlap and hamiltonian in the folding approximation, for several reasons. The folding approximation is indeed the natural representation for discussing the asymptotic behavior of the three cluster system, as the antisymmetrization between clusters vanishes at large inter-cluster distances. The calculation of generating functions in this approximation is also illuminating for the subsequent derivation of generating functions in a fully antisymmetrized setting, as the principles are identical, but the implementation is more complex. Finally, the folding approximation provides an interesting model to discuss the importance of antisymmetrization in the interaction region.
### A Matrix elements for the overlap
The overlap of two generating functions of the form (45) is easily obtained, and can be written as
$$I(ϵ,𝐐_1,𝐐_2;ϵ^{},𝐐_1^{},𝐐_2^{})=\mathrm{\Delta }^3\mathrm{exp}\left\{\frac{1}{\mathrm{\Delta }}\underset{i=1}{\overset{2}{}}\left[𝐐_i𝐐_i^{}+\frac{1}{2}\left(ϵ^{}𝐐_i^2+ϵ𝐐_i^2\right)\right]\right\}$$
(60)
where
$`\mathrm{\Delta }=1ϵϵ^{}`$
After substitution of (47) in (60), one obtains
$$I(ϵ,S,𝐐_1,𝐐_2;ϵ^{},S^{},𝐐_1^{},𝐐_2^{})=\mathrm{\Delta }^3\mathrm{exp}\left\{\frac{SS^{}}{\mathrm{\Delta }}\stackrel{ˇ}{𝐐}_1\stackrel{ˇ}{𝐐}_1^{}\right\}\mathrm{exp}\left\{\frac{SS^{}}{\mathrm{\Delta }}\stackrel{ˇ}{𝐐}_2\stackrel{ˇ}{𝐐}_2^{}\right\}$$
(61)
It is interesting to note that the arguments of the exponential factors are diagonal in the generator coordinates $`S`$ and $`S^{}`$. To obtain an expansion in terms of angular momenta, the well known relation
$$\mathrm{exp}\{𝐚𝐛\}=4\pi \underset{lm}{}i_l(ab)Y_{lm}^{}(\widehat{𝐚})Y_{lm}^{}(\widehat{𝐛})$$
(62)
can be applied, where $`i_l(x)`$ is the Modified Spherical Bessel function of the first kind. Substitution of (62) in (61), and applying traditional angular momentum coupling techniques leads to
$`I(ϵ,S,𝐐_1,𝐐_2;ϵ^{},S^{},𝐐_1^{},𝐐_2^{})`$ (63)
$`=`$ $`\left(4\pi \right)^2{\displaystyle \underset{l_1,m_1,l_2,m_2}{}}i_{l_1}\left({\displaystyle \frac{SS^{}}{\mathrm{\Delta }}}\right)i_{l_2}\left({\displaystyle \frac{SS^{}}{\mathrm{\Delta }}}\right)Y_{l_1m_1}^{}(\stackrel{ˇ}{𝐐}_1)Y_{l_2m_2}^{}(\stackrel{ˇ}{𝐐}_2)Y_{l_1m_1}^{}(\stackrel{ˇ}{𝐐}_1^{})Y_{l_2m_2}^{}(\stackrel{ˇ}{𝐐}_2^{})`$ (64)
$`=`$ $`\left(4\pi \right)^2{\displaystyle \underset{\left(l_1l_2\right)LM}{}}i_{l_1}\left({\displaystyle \frac{SS^{}}{\mathrm{\Delta }}}\right)i_{l_2}\left({\displaystyle \frac{SS^{}}{\mathrm{\Delta }}}\right)\left\{Y_{l_1}(\stackrel{ˇ}{𝐐}_1)\times Y_{l_2}(\stackrel{ˇ}{𝐐}_2)\right\}_{LM}\left\{Y_{l_1}^{}(\stackrel{ˇ}{𝐐}_1^{})Y_{l_2}^{}(\stackrel{ˇ}{𝐐}_2^{})\right\}_{LM}`$ (65)
The reduced generating function (cfr. (59)) then becomes
$$_{\left(l_1l_2\right)LM;(l_1l_2)LM}(ϵ,S;ϵ^{},S^{})=\left(4\pi \right)^2\mathrm{\Delta }^3i_{l_1}\left(\frac{SS^{}}{\mathrm{\Delta }}\right)i_{l_2}\left(\frac{SS^{}}{\mathrm{\Delta }}\right)$$
(66)
This reduced generating function (66) is diagonal in the partial angular momenta $`l_1`$ and $`l_2`$, and independent of total angular momentum $`L`$, thus valid for all angular momenta $`L`$ compatible with $`l_1`$ and $`l_2`$.
The matrix elements with quantumnumbers $`K`$ (also $`n`$) can now be obtained through a standard procedure, e.g. by differentiating (66) with respect to $`S`$ and $`S^{}`$ ($`ϵ`$ and $`ϵ^{}`$ for $`n`$). In particular, as (66) depends on the product of $`S`$ and $`S^{}`$ only, the overlap is diagonal in $`K`$. Likewise, the dependence on $`ϵ`$ and $`ϵ^{}`$ appears as a product in the factor $`\mathrm{\Delta }`$, so that the overlap is again diagonal in $`n`$. The fact that the generated overlap matrix elements are diagonal is a confirmation of the fact that these matrix elements were generated in an orthogonal basis. The calculation of the matrix is however not unimportant, as it provides a straightforward way to obtain the norm of the generator coordinate basisfunction (50).
### B Matrix elements for the kinetic energy
In order to calculate the kinetic energy one can use the properties of the oscillator basis. The matrix elements of the kinetic energy of relative motion of the clusters are related to those of the oscillator potential by the virial theorem
$`N\left|\widehat{T}_R\right|N`$ $`=`$ $`N\left|\widehat{V}_O\right|N={\displaystyle \frac{1}{2}}E_N`$ (67)
$`N\pm 2\left|\widehat{T}_R\right|N`$ $`=`$ $`N\pm 2\left|\widehat{V}_O\right|N`$ (68)
where
$$\widehat{V}_O=\frac{1}{2}\mathrm{}\omega \underset{i=1}{}𝐪_i^2$$
(69)
and $`E_N=\mathrm{}\omega \left[N+3\right]=\mathrm{}\omega \left[2n+K+3\right]`$ is the oscillator energy of relative motion.
One easily obtains a representation of the oscillator potential in the manifold spanned by the generating function (49):
$`\left(𝐪_1^2+𝐪_2^2\right)\mathrm{\Psi }(𝐪_1,𝐪_2|ϵ,S;\stackrel{ˇ}{𝐐}_1,\stackrel{ˇ}{𝐐}_2)`$ (70)
$`=`$ $`\left[(1+ϵ)^2{\displaystyle \frac{d}{dϵ}}+(1+ϵ)S{\displaystyle \frac{d}{dS}}+3(1+ϵ)\right]\mathrm{\Psi }(𝐪_1,𝐪_2|ϵ,S;\stackrel{ˇ}{𝐐}_1,\stackrel{ˇ}{𝐐}_2)`$ (71)
so that the kinetic energy operator $`\widehat{T}_R`$ can be represented by
$$𝒯_R=\frac{1}{2}\mathrm{}\omega \left[(1ϵ)^2\frac{d}{dϵ}+(1ϵ)S\frac{d}{dS}+3(1ϵ)\right]$$
(72)
In a wider context, one can associate this operator with the generators $`^{(+)}`$, $`^{()}`$ and $`^{\left(0\right)}`$
$`^{(+)}`$ $`=`$ $`\left[ϵ^2{\displaystyle \frac{d}{dϵ}}+ϵS{\displaystyle \frac{d}{dS}}+3ϵ\right],`$ (73)
$`^{()}`$ $`=`$ $`\left[{\displaystyle \frac{d}{dϵ}}\right]`$ (74)
$`^{\left(0\right)}`$ $`=`$ $`\left[2ϵ{\displaystyle \frac{d}{dϵ}}+S{\displaystyle \frac{d}{dS}}+3\right]`$ (75)
of the $`Sp(2,R)`$ group classifying the space of relative motion of three-particle states, and whose irreducible representations are labelled by hypermomentum $`K`$ so that
$`𝒱_R`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{}\omega \left[^{\left(0\right)}+^{(+)}+^{()}\right]`$ (76)
$`𝒯_R`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{}\omega \left[^{\left(0\right)}^{(+)}^{()}\right]`$ (77)
Using again the shorthand notation $`\nu _0=(l_1l_2)LM`$, matrix elements are then readily found to be
$`n,K,\nu _0\left|\widehat{T}_R\right|n,K,\nu _0`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{}\omega \left[2n+K+3\right]`$ (78)
$`n+1,K,\nu _0\left|\widehat{T}_R\right|n,K,\nu _0`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{}\omega \sqrt{\left(n+1\right)\left(n+K+3\right)}`$ (79)
$`n1,K,\nu _0\left|\widehat{T}_R\right|n,K,\nu _0`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{}\omega \sqrt{n\left(n+K+2\right)}`$ (80)
The full kinetic energy of the three cluster system must include the internal kinetic energy of the clusters: $`\widehat{T}=\widehat{T}_R+\widehat{T}_{cl}`$. As we consider frozen $`s`$-clusters only, this contribution is purely diagonal and equal to
$$T_{cl}=\frac{3}{4}\mathrm{}\omega \underset{i=1}{\overset{3}{}}\left(A_i1\right)=\frac{3}{4}\mathrm{}\omega \left(A3\right)$$
(81)
### C Matrix elements for a Gaussian potential
For a Gaussian two-body interaction with strength $`V_0`$ and range $`a`$
$$V(𝐫_i,𝐫_j)=V_0\mathrm{exp}\left\{\frac{(𝐫_i𝐫_j)^2}{a^2}\right\}$$
(82)
the folding potential (5) for (0s)-clusters can be calculated analytically. The result obtained from (6) is again of a Gaussian form but with a modified strength and interaction length now also depending on the oscillator parameter $`b`$
$$V^{(F)}(𝐑_{\tau \upsilon })=V_0z^{3/2}\mathrm{exp}\left\{\frac{z}{a^2}𝐑_{\tau \upsilon }^2\right\}$$
(83)
$`z=\left(1+{\displaystyle \frac{b^2}{a^2}}[2\mu _{\tau \upsilon }^1]\right)^1,\mu _{\tau \upsilon }={\displaystyle \frac{A_\tau A_\upsilon }{A_\tau +A_\upsilon }}`$
The matrix element of a Gaussian potential between two generating functions of the form (45) is not necessarily diagonal in terms of the chosen Jacobi coordinates. In the folding approximation one can however easily find a set Jacobi coordinates in which the potential matrix element is diagonal. Both sets of Jacobi coordinates will then be related by an orthogonal transformation.
In what follows, we distinguish the two types of Jacobi coordinates as follows: the original coordinates are denoted by $`𝐪`$ and were introduced in section II A, and shown in Fig. 1; the diagonalizing coordinates will be denoted by $`𝐱`$. An explicit instance of the latter are easily obtained. Indeed, for any choice of two clusters $`i`$ and $`j`$ ($`ij)`$, $`k`$ being the third particle, one obtains a system of coordinates $`𝐱`$, uniquely defined by index $`k`$, as follows,
$`𝐱_1`$ $`=`$ $`\sqrt{c_{1,k}}(𝐫_i𝐫_j)`$ (84)
$`𝐱_2`$ $`=`$ $`\sqrt{c_{2,k}}(𝐫_k{\displaystyle \frac{A_i𝐫_i+A_j𝐫_j}{A_i+A_j}})`$ (85)
$$c_{1,k}=\mu _{ij}=\frac{A_iA_j}{A_i+A_j},c_{2,k}=\frac{A_k(A_i+A_j)}{A_i+A_j+A_k}$$
(86)
Each such Jacobi coordinate system $`k`$ leads to a diagonal representation for the potential energy between clusters $`i`$ and $`j`$. As both Jacobi systems $`𝐪`$ and $`𝐱`$ are related through an orthogonal transformation, we can invoke the Raynal-Revai theorem . The latter states that any orthogonal transformation of Jacobi coordinates leads to an orthogonal transformation of the Hyperspherical Harmonics (44) in the wave function (41), preserving the hypermomentum quantum number $`K`$
$$H_K^{\left(l_1l_2\right)LM}(\theta _𝐪,\stackrel{ˇ}{𝐪}_1,\stackrel{ˇ}{𝐪}_2)=\underset{\lambda _1\lambda _2}{}O_{\lambda _1\lambda _2}^{l_1l_2}H_K^{\left(\lambda _1\lambda _2\right)LM}(\theta _𝐱,\stackrel{ˇ}{𝐱}_1,\stackrel{ˇ}{𝐱}_2)$$
(87)
The $`O_{\lambda _1\lambda _2}^{l_1l_2}`$ are known as the Raynal-Revai coefficients .
This transformation can then be used to obtain the matrix elements of the potential in the original set of coordinates $`𝐪`$ through
$`n,K,\left(l_1l_2\right)LM\left|\widehat{V}\right|n^{},K^{},\left(l_1^{}l_2^{}\right)LM_{(𝐪)}`$ (88)
$`=`$ $`{\displaystyle \underset{\lambda _1\lambda _2}{}}{\displaystyle \underset{\lambda _1^{}\lambda _2^{}}{}}O_{\lambda _1\lambda _2}^{l_1l_2}O_{\lambda _1^{}\lambda _2^{}}^{l_1^{}l_2^{}}n,K,\left(\lambda _1\lambda _2\right)LM\left|\widehat{V}\right|n^{},K^{},\left(\lambda _1^{}\lambda _2^{}\right)LM_{(𝐱)}`$ (89)
where the matrix element on the right hand side, which is diagonal in the $`𝐱`$ representation, can be calculated in a straightforward way.
To obtain the latter matrix element we consider a generating function of an identical structure as (45), in which we replace the $`𝐪`$ coordinates and $`𝐐`$ generator coordinates by $`𝐱`$ and $`𝐗`$. A generating function for the two-body matrix elements is then easily obtained as a product of two integrals over $`𝐱_1`$ and $`𝐱_2`$ leading to
$$V(ϵ,S,\stackrel{ˇ}{𝐗}_1,\stackrel{ˇ}{𝐗}_\mathrm{𝟐};ϵ^{},S^{},\stackrel{ˇ}{𝐗}_1^{},\stackrel{ˇ}{𝐗}_\mathrm{𝟐}^{})=V_{pre}\mathrm{exp}\left\{\frac{SS^{}}{\mathrm{\Lambda }}\stackrel{ˇ}{𝐗}_1\stackrel{ˇ}{𝐗}_1^{}\right\}\mathrm{exp}\left\{\frac{SS^{}}{\mathrm{\Delta }}\stackrel{ˇ}{𝐗}_\mathrm{𝟐}\stackrel{ˇ}{𝐗}_\mathrm{𝟐}^{}\right\}$$
(90)
$`V_{pre}`$ $`=`$ $`V_0\left(\mathrm{\Delta }\mathrm{\Lambda }\right)^{3/2}\mathrm{exp}\left\{{\displaystyle \frac{\gamma }{2\mathrm{\Delta }\mathrm{\Lambda }}}\left[\left(\xi ^{}S\right)^2+\left(\xi S^{}\right)^2\right]\right\}`$
$`\xi `$ $`=`$ $`1+ϵ,\xi ^{}=1+ϵ^{}`$
$`\mathrm{\Delta }`$ $`=`$ $`1ϵϵ^{}=\xi +\xi ^{}\xi \xi ^{}`$
$`\mathrm{\Lambda }`$ $`=`$ $`\mathrm{\Delta }+\gamma \xi \xi ^{}`$
$`\gamma `$ $`=`$ $`{\displaystyle \frac{b^2}{a^2}}{\displaystyle \frac{1}{c_{1,\alpha }}}`$
By using (62), one again eliminates factorized terms and sums to obtain the following reduced generating function (cfr. (59)), in complete analogy to the procedure taken for (66):
$$𝒱_{\left(l_1l_2\right)LM;(l_1l_2)LM}(ϵ,S;ϵ^{},S^{})=V_{pre}i_{l_1}\left(\frac{SS^{}}{\mathrm{\Lambda }}\right)i_{l_2}\left(\frac{SS^{}}{\mathrm{\Delta }}\right)$$
(91)
As was the case with the overlap, the matrix elements of the potential in the $`𝐱`$ coordinate system are diagonal with respect to the partial angular momenta $`l_1`$ and $`l_2`$ and do not depend directly upon the total angular momentum $`L`$ ; this is a direct consequence of the characteristics of the operator. Expression (91) is again valid for all values of the total angular momenta $`L`$ that are compatible with the given partial angular momenta $`l_1`$ and $`l_2`$, and thus represents a generating function for matrix elements with specific total angular momentum $`L`$.
One obtains matrix elements with specific $`K`$ and $`n`$ quantumnumbers through the standard procedures (differentiation, recurrence relations, …). So for example one obtains
$`n,K,\left(l_1l_2\right)LM\left|\widehat{V}\right|n^{},K^{},\left(l_1l_2\right)LM_{(𝐱)}`$ (92)
$`=`$ $`{\displaystyle \frac{\left({\displaystyle \frac{d}{dϵ}}\right)^n\left({\displaystyle \frac{d}{dϵ^{}}}\right)^n^{}\left({\displaystyle \frac{d}{dS}}\right)^K\left({\displaystyle \frac{d}{dS^{}}}\right)^K^{}𝒱_{\left(l_1l_2\right)LM;(l_1l_2)LM}(ϵ,S;ϵ^{},S^{})|}{S=S^{}=0}}`$ (93)
$`ϵ=ϵ^{}=0`$ (94)
### D Matrix elements for the Coulomb potential
The Coulomb interaction in the folding approximation between two clusters with $`Z_\tau `$ and $`Z_\upsilon `$ number of protons is given by (6) as
$$V_C^{(F)}(𝐑_{\tau \upsilon })=Z_\tau Z_\upsilon e^2\underset{iA_\tau }{}\underset{jA_\upsilon }{}𝑑\tau _\tau 𝑑\tau _\upsilon |\mathrm{\Psi }_\tau (A_\tau )|^2\frac{1}{\left|𝐫_i𝐫_j+𝐑_{\tau \upsilon }\right|}|\mathrm{\Psi }_\upsilon (A_\upsilon )|^2$$
(95)
A straightforward calculation of its matrix elements is impractical and very tedious. By however using the following Gauss-transform
$$\frac{1}{r}=\frac{2}{\sqrt{\pi }}_0^{\mathrm{}}𝑑x\mathrm{exp}\{r^2x^2\}$$
(96)
one rewrites (95) as
$`V_C^{(F)}(\tau \upsilon )`$ $`=`$ $`{\displaystyle \frac{2Z_\tau Z_\upsilon e^2}{\sqrt{\pi }}}{\displaystyle \underset{iA_\tau }{}}{\displaystyle \underset{jA_\upsilon }{}}{\displaystyle 𝑑\tau _\tau 𝑑\tau _\upsilon |\mathrm{\Psi }_\tau (A_\tau )|^2_0^{\mathrm{}}𝑑x\mathrm{exp}\{\left(𝐫_i𝐫_j+𝐑_{\tau \upsilon }\right)^2x^2\}|\mathrm{\Psi }_\upsilon (A_\upsilon )|^2}`$ (97)
$`=`$ $`{\displaystyle \frac{2Z_\tau Z_\upsilon e^2}{\sqrt{\pi }b}}{\displaystyle _0^{\mathrm{}}}𝑑\gamma z^{3/2}\mathrm{exp}\left\{{\displaystyle \frac{R_{\tau \upsilon }^2}{b^2}}\gamma ^2z\right\}`$ (98)
where
$$z=\left(1+t\gamma ^2\right)^1,t=2\mu _{\tau \upsilon }^1,\gamma =bx$$
(99)
and its matrix elements can be obtained by integrating matrix elements (depending on $`z`$) of the Gaussian potential.
Introducing the integration variable $`s`$
$$s=\frac{t\gamma ^2}{1+t\gamma ^2},$$
(100)
transforms (98) to
$$V_C^{(F)}(\tau \upsilon )=\frac{2Z_\tau Z_\upsilon e^2}{\sqrt{\pi }b}\frac{1}{2\sqrt{t}}_0^1𝑑ss^{\frac{1}{2}}\mathrm{exp}\left\{\frac{R_{\tau \upsilon }^2}{tb^2}s\right\}=\frac{2Z_\tau Z_\upsilon e^2}{\sqrt{\pi }b}erf(\frac{R_{\tau \upsilon }^2}{tb^2})$$
(101)
This form shows that the integration can be reduced to a finite interval. It also shows that the Coulomb interaction between clusters does not behave as $`R_{\tau \upsilon }^1`$ as could be expected. For very large value of $`R_{\tau \upsilon }`$ however (101) properly reduces to $`\frac{Z_\tau Z_\upsilon e^2}{R_{\tau \upsilon }}`$, due to the asymptotic form of the error function.
### E Asymptotic behavior of the potential contributions
As the folding model is used for defining the asymptotic channels, it is clear that the asymptotic behavior of the potential energy matrix elements in this model will be of vital importance for the rate of convergence of the AM solutions. The effective potential in terms of the hyperradius $`\rho `$ is defined by integrating the folding potential over all hyperangles as formally indicated by
$`W(\rho )`$ $`=`$ $`{\displaystyle \underset{\tau ,\upsilon }{}}W_{\tau \upsilon }(\rho )W_{K,l_1,l_2}(\rho )`$ (102)
$`=`$ $`K,l_1,l_2\left|{\displaystyle \underset{\tau ,\upsilon }{}}\widehat{V}(𝐑_{\tau \upsilon })\right|K,l_1,l_2`$ (103)
Its asymptotic behavior is then obtained for large values of $`\rho `$.
In the diagonal representation with Jacobi coordinates $`𝐱`$ the calculation of $`W(\rho )`$ between one pair of clusters $`\tau `$ and $`\upsilon `$ amounts to a straightforward integration over the hyperangles.
The Gaussian interaction then exhibits the following asymptotic behavior
$$W_{\tau \upsilon }(\rho )(1)^{Kl_1l_2}V_0N_K^{l_1l_2}N_K^{l_1l_2}\left(\frac{1}{2}\right)^{K+3}\mathrm{\Gamma }(l_1+\frac{3}{2})\left(\frac{\sqrt{\mu _{\tau \upsilon }}a}{\rho }\right)^{2l_1+3}\left(\begin{array}{c}\frac{K+l_1l_2+1}{2}\\ Kl_1l_2\end{array}\right)\left(\begin{array}{c}\frac{K+l_1l_2+1}{2}\\ Kl_1l_2\end{array}\right)$$
(104)
which indeed shows a worst-case behavior (for $`l_1=0`$) of the form $`1/\rho ^3`$ as predicted in section III D.
The analogous evaluation for the Coulomb interaction leads to the following (exact) expression
$$W_{\tau \upsilon }(\rho )=\frac{Z_{K,l_1,l_2}^{eff}}{\rho }$$
(105)
with
$`Z_{K,l_1,l_2}^{eff}`$ $`=`$ $`{\displaystyle \frac{1}{2}}Z_\tau Z_\upsilon e^2\sqrt{\mu _{\tau \upsilon }}`$ (116)
$`{\displaystyle \underset{n,m=0}{\overset{\frac{Kl_1l_2}{2}}{}}}()^{Kl_1l_2nm}\left(\begin{array}{c}\frac{Kl_1+l_2+1}{2}\\ n\end{array}\right)\left(\begin{array}{c}\frac{Kl_1+l_2+1}{2}\\ m\end{array}\right)\left(\begin{array}{c}\frac{K+l_1l_2+1}{2}\\ \frac{Kl_1l_2}{2}n\end{array}\right)\left(\begin{array}{c}\frac{K+l_1l_2+1}{2}\\ \frac{Kl_1l_2}{2}m\end{array}\right)`$
$`B(Kl_1nm+{\displaystyle \frac{3}{2}},n+m+l_1+1)`$
and $`B`$ stands for the beta-function.
These results corroborate the fact that special care should be taken to get properly convergent results, even more critically when a Coulomb contribution between the clusters is present.
## VI Matrix elements with Full Antisymmetrization
When considering full antisymmetrization between all particles of the three-cluster wave function the normalization of the basis states becomes a non-trivial problem. The overlap, i.e. the matrix representation of the antisymmetrization operator $`𝒜`$, needs to be explicitly calculated.
As was carried out in the previous section one would normally start from the generating state (45) to obtain generating matrix elements in terms of $`ϵ`$ and $`ϵ^{}`$, facilitating the treatment of the hyperangular coordinates. We propose an alternative representation for the scaled generating state (45) which is more suited to our calculations, and hereto introduce the following integral transformation
$$\mathrm{\Psi }\left(𝐪|ϵ,𝐐\right)=𝑑𝐤\mathrm{exp}\{k^2\}\mathrm{\Psi }(𝐪;𝐐+\sqrt{2ϵ}𝐤)$$
(117)
This allows to obtain (45) by scaling on the generator coordinate only. In other words, generating matrix elements can be obtained with the simpler generating state (31) and scaled later on, reducing effectively the calculational burden.
### A Matrix elements for the overlap
Because the individual cluster states are Slater determinants one can use the familiar determinantal formulae to calculate the generating function. Starting then from the generating state (31) for the single-particle orbitals one obtains
$$I(𝐐_1,𝐐_2;𝐐_1^{},𝐐_2^{})=\underset{\nu }{}D^{(\nu )}\mathrm{exp}\left\{\underset{i,j=1}{\overset{2}{}}B_{ij}^{(\nu )}𝐐_i𝐐_j^{}\right\}$$
(118)
The coefficients $`D^{(\nu )}`$ and $`B_{ij}^{(\nu )}`$ as well as the number of the terms depend on the specific type of three-cluster configuration, viz. the number of nucleons per cluster and their spin-isospin quantum numbers.
In order to reduce (118) with respect to the angular quantum numbers it is profitable to diagonalize the forms in the exponentials of (118). This can easily be achieved by diagonalizing its 2 by 2 coefficient matrix $`B_{ij}^{(\nu )}`$. This again amounts to making an orthogonal transformation from the original Jacobi coordinates $`𝐪`$ to new Jacobi coordinates $`𝐱`$. This orthogonal transformation will induce a corresponding transformation of the basis functions which can be handled by the Raynal-Revai theorem as discussed in the previous section. The block-diagonal form of an exponential term can be written as
$$\mathrm{exp}\left\{\underset{i=1}{\overset{2}{}}\lambda _i𝐗_i𝐗_i^{}\right\}$$
(119)
We now introduce the scaling on $`ϵ`$ and $`ϵ^{}`$ by carrying out the transformation (117) on both generator coordinates $`𝐗`$ and $`𝐗^{}`$ for every (block-diagonal) term in (118), leading to
$$I(ϵ,𝐗_1,𝐗_2;ϵ^{},𝐗_1^{},𝐗_2^{})=\left(\mathrm{\Delta }_1\mathrm{\Delta }_2\right)^{3/2}\mathrm{exp}\left\{\underset{i=1}{\overset{2}{}}\frac{\lambda _i}{\mathrm{\Delta }_i}\left[𝐗_i𝐗_i^{}+\frac{\lambda _i}{2}\left(ϵ^{}𝐗_i^2+ϵ𝐗_i^2\right)\right]\right\}$$
(120)
where
$`\mathrm{\Delta }_i=1\lambda _i^2ϵϵ^{}`$
In order to further reduce the generating function in the partial angular momenta $`l_1`$ and $`l_2`$ we substitute (47) in (120)
$`I(ϵ,S,\stackrel{ˇ}{𝐗}_1,\stackrel{ˇ}{𝐗}_2;ϵ^{},S^{},\stackrel{ˇ}{𝐗}_1^{},\stackrel{ˇ}{𝐗}_2^{})`$ (121)
$`=`$ $`\left(\mathrm{\Delta }_1\mathrm{\Delta }_2\right)^{3/2}\mathrm{exp}\left\{{\displaystyle \frac{\lambda _1^2\lambda _2^2}{2\mathrm{\Delta }_1\mathrm{\Delta }_2}}\left[S^2ϵ^{}+S^2ϵ\right]\right\}\mathrm{exp}\left\{{\displaystyle \frac{\lambda _1}{\mathrm{\Delta }_1}}SS^{}\stackrel{ˇ}{𝐗}_1\stackrel{ˇ}{𝐗}_1^{}\right\}\mathrm{exp}\left\{{\displaystyle \frac{\lambda _2}{\mathrm{\Delta }_2}}SS^{}\stackrel{ˇ}{𝐗}_2\stackrel{ˇ}{𝐗}_2^{}\right\}`$ (122)
By using (62) this produces in complete analogy to (66) a generating function for all total angular momenta $`L`$ compatible with $`l_1`$ and $`l_2`$:
$`_{(l_1,l_2)LM;(l_1l_2)LM}(ϵ,S;ϵ^{},S^{})`$ (123)
$`=`$ $`\left(4\pi \right)^2\left(\mathrm{\Delta }_1\mathrm{\Delta }_2\right)^{3/2}\mathrm{exp}\left\{{\displaystyle \frac{\lambda _1^2\lambda _2^2}{2\mathrm{\Delta }_1\mathrm{\Delta }_2}}\left[S^2ϵ^{}+S^2ϵ\right]\right\}i_{l_1}\left({\displaystyle \frac{\lambda _1}{\mathrm{\Delta }_1}}SS^{}\right)i_{l_2}\left({\displaystyle \frac{\lambda _2}{\mathrm{\Delta }_2}}SS^{}\right)`$ (124)
If $`\lambda _1=\lambda _2`$, the overlap depends on the factor $`SS^{}`$ only and is therefore diagonal with respect to hypermomentum $`K`$, as was the case in the folding approximation. In the general case one will have $`\lambda _1\lambda _2`$, and (124) generates nondiagonal matrix elements in $`K`$. Hypermomentum is thus no longer a good quantum number for three-cluster systems, contrary to the folding approximation.
The generating function for the overlap with $`K=0`$, and consequently for angular momenta $`l_1=l_2=L=0`$, is obtained immediately by putting $`S=S^{}=0`$. This simplicity is an indication of the very suitable form of our generating function.
The matrix elements $`I_{n,n^{}}^{K,\nu _0;K^{},\nu _0}`$ can now again be generated by a standard procedure, such as differentiation or recurrence relations.
### B Matrix elements for the kinetic energy
The matrix elements of the kinetic energy operator can be derived without use of a generating function. One of the effects of the antisymmetrization operator is to mix basis states within a fixed oscillator shell. The “diagonal” (i.e. within a shell) and “off-diagonal” (within neighboring shells) kinetic energy matrix elements are easily found to be connected to the matrix elements of the overlap by
$`T_{n,n^{}}^{K,\nu _0;K^{},\nu _0}={\displaystyle \frac{1}{2}}\mathrm{}\omega \left[2n+K+3+{\displaystyle \frac{3}{2}}\left(A3\right)\right]I_{n,n^{}}^{K,\nu _0;K^{},\nu _0}\delta _{2n+K,2n^{}+K^{}}`$
$`T_{n+1,n^{}}^{K,\nu _0;K^{},\nu _0}={\displaystyle \frac{1}{2}}\mathrm{}\omega \sqrt{\left(n+1\right)\left(n+K+3\right)}I_{n,n^{}}^{K,\nu _0;K^{},\nu _0}\delta _{2n+K,2n^{}+K^{}}`$
where the restriction on the quantum numbers to remain on the same oscillator shell ($`2n+K=2n^{}+K^{}`$) has been accounted for.
### C Matrix elements for a Gaussian potential
If one considers a Gaussian form for the nucleon-nucleon potential, and calculates a generating matrix element for the interaction using Slater determinants with individual orbitals of the form (31) one obtains terms of the form
$$V(𝐐_1,𝐐_\mathrm{𝟐};𝐐_1^{},𝐐_\mathrm{𝟐}^{})=V_0\left(1\zeta \right)^{3/2}\mathrm{exp}\left\{\zeta \left[\underset{i=1}{\overset{2}{}}\left(C_i𝐐_i+C_i^{}𝐐_i^{}\right)\right]^2+\underset{i,j=1}{\overset{2}{}}B_{ij}𝐐_i𝐐_j^{}\right\}$$
(125)
where
$$\zeta =\frac{2b^2}{2b^2+a^2}$$
(126)
and $`V_0`$ stands for any of the even ($`V_{13}`$ and $`V_{31}`$) and odd ($`V_{11}`$ and $`V_{33}`$) components of the $`NN`$–interaction. Again $`b`$ is the oscillator radius and $`a`$ the range of the potential well. The first term in the exponent of (125) contains the factor $`\zeta `$ and the vectors $`C_1𝐐_1+C_2𝐐_2`$ and $`C_1^{}𝐐_1^{}+C_2^{}𝐐_2^{}`$. The latter have a simple meaning: they define the distance between of the two clusters which respectively contain one of the nucleons in the interacting pair.
In the diagonal representation (125) becomes
$$V(𝐗_1,𝐗_2;𝐗_1^{},𝐗_\mathrm{𝟐}^{})=V_0\mathrm{exp}\{\zeta \left[\underset{i=1}{\overset{2}{}}(G_i𝐗_i+G_i^{}𝐗_i^{})\right]^2+\underset{i=1}{\overset{2}{}}\lambda _i𝐗_i𝐗_i^{})\}$$
(127)
and the coefficients $`C_i`$, $`C_i^{}`$ and $`G_i`$, $`G_i^{}`$ are trivially related by the orthogonal diagonalizing transformation.
In order to eliminate crossterms in $`𝐗_1𝐗_2`$ and $`𝐗_1^{}𝐗_2^{}`$ we introduce an additional transformation through the integral identity:
$$e^{𝐚^2}=\frac{1}{\pi ^{3/2}}_{\mathrm{}}^{\mathrm{}}𝑑𝐙\mathrm{exp}\left\{Z^2+2i(𝐙𝐚)\right\}$$
(128)
leading to the following integral form for the block (127)
$`V(𝐗_1,𝐗_2;𝐗_1^{},𝐗_\mathrm{𝟐}^{})`$ (129)
$`=`$ $`V_0{\displaystyle \frac{1}{\pi ^{3/2}}}\mathrm{exp}\left\{{\displaystyle \underset{i=1}{\overset{2}{}}}\lambda _i𝐗_i𝐗_i^{}\right\}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑𝐙\mathrm{exp}\left\{Z^2+2i\zeta {\displaystyle \underset{i=1}{\overset{2}{}}}\left(G_i𝐙𝐗_i+G_i^{}𝐙𝐗_i^{}\right)\right\}`$ (130)
Using (117) on both $`𝐗_i`$ and $`𝐗_i^{}`$, one obtains after integrating over the corresponding $`𝐤_i`$ and $`𝐤_i^{}`$ :
$`V(ϵ;𝐗_1,𝐗;ϵ^{};𝐗_1^{},𝐗_\mathrm{𝟐}^{})=\left(\mathrm{\Delta }_1\mathrm{\Delta }_2\right)^{3/2}\mathrm{exp}\left\{{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \frac{\lambda _i}{2\mathrm{\Delta }_i}}\left(ϵ^{}\lambda _i𝐗_i^2+ϵ\lambda _i𝐗_i^2+2𝐗_i𝐗_i^{}\right)\right\}`$ (131)
$`{\displaystyle \frac{1}{\pi ^{3/2}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑𝐙\mathrm{exp}\left\{\mathrm{\Delta }𝐙^2+2i\zeta {\displaystyle \underset{i=1}{\overset{2}{}}}\left(\xi _i𝐙𝐗_i+\xi _i^{}𝐙𝐗_i^{}\right)\right\}`$ (132)
where again $`\mathrm{\Delta }_i=1\lambda _i^2ϵϵ^{}`$ and
$`\xi _i`$ $`=`$ $`{\displaystyle \frac{G_i+ϵ^{}G_i^{}\lambda _i}{\mathrm{\Delta }_i}},\xi _i^{}={\displaystyle \frac{G_i^{}+ϵG_i\lambda _i}{\mathrm{\Delta }_i}}`$
$`\mathrm{\Delta }`$ $`=`$ $`\left\{1+{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \frac{1}{\mathrm{\Delta }_i}}\left[ϵG_i^2+ϵ^{}G_i^2+2ϵϵ^{}\lambda _iG_iG_i^{}\right]\right\}`$
To obtain a reduced generating function for specific angular momentum quantum numbers, we again use expansion formulae of an exponential in terms of spherical harmonics. For exponential terms with a real scalar product we use (62), whereas for exponential terms with an imaginary scalar product we use
$$\mathrm{exp}\{i𝐚𝐛\}=4\pi \underset{lm}{}i^lj_l(ab)Y_{lm}^{}(\stackrel{ˇ}{𝐚})Y_{lm}^{}(\stackrel{ˇ}{𝐛})$$
(133)
where $`j_l(ab)`$ is the well-known spherical Bessel function.
Applying both (62) and (133) to (132) leads to
$`V(ϵ,𝐗_1,𝐗;ϵ^{},𝐗_1^{},𝐗_\mathrm{𝟐}^{})=(4\pi )^6\left(\mathrm{\Delta }_1\mathrm{\Delta }_2\right)^{3/2}\mathrm{exp}{\displaystyle \underset{i=1}{\overset{2}{}}}\left\{{\displaystyle \frac{\lambda _i^2}{2\mathrm{\Delta }_i}}\left(ϵ^{}𝐗_i^2+ϵ𝐗_i^2\right)\right\}`$ (134)
$`{\displaystyle \underset{k_1,k_2,k_1^{},k_2^{},k_1^{\prime \prime },k_2^{\prime \prime }}{}}i^{k_1^{}+k_2^{}+k_1^{\prime \prime }+k_2^{\prime \prime }}{\displaystyle \underset{i=1}{\overset{2}{}}}i_{k_i}\left({\displaystyle \frac{\lambda _i}{\mathrm{\Delta }_i}}X_iX_i^{}\right)`$ (135)
$`{\displaystyle \frac{1}{\pi ^{3/2}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑𝐙\mathrm{exp}\{\mathrm{\Delta }Z^2\}P_{k_1,k_2,k_1^{},k_2^{},k_1^{\prime \prime },k_2^{\prime \prime }}{\displaystyle \underset{i=1}{\overset{2}{}}}j_{k_i^{}}\left(2\xi _iX_iZ\right)j_{k_i^{\prime \prime }}\left(2\xi _i^{}X_i^{}Z\right)`$ (136)
Before performing the integration over $`𝐙`$, we again make the substitution (47) in (136) to obtain
$`V(ϵ,S,\stackrel{ˇ}{𝐗}_1,\stackrel{ˇ}{𝐗};ϵ^{},S^{},\stackrel{ˇ}{𝐗}_1^{},\stackrel{ˇ}{𝐗}_\mathrm{𝟐}^{})`$ (137)
$`=`$ $`(4\pi )^6\left(\mathrm{\Delta }_1\mathrm{\Delta }_2\right)^{3/2}\mathrm{exp}\left\{{\displaystyle \frac{\lambda _1^2\lambda _2^2}{2\mathrm{\Delta }_1\mathrm{\Delta }_2}}\left(ϵ^{}S^2+ϵS^2\right)\right\}`$ (140)
$`{\displaystyle \underset{k_1,k_2,k_1^{},k_2^{},k_1^{\prime \prime },k_2^{\prime \prime }}{}}i^{k_1^{}+k_2^{}+k_1^{\prime \prime }+k_2^{\prime \prime }}i_{k_1}\left({\displaystyle \frac{\lambda _1}{\mathrm{\Delta }_1}}SS^{}\right)i_{k_2}\left({\displaystyle \frac{\lambda _2}{\mathrm{\Delta }_2}}SS^{}\right)`$
$`{\displaystyle \frac{1}{\pi ^{3/2}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑𝐙\mathrm{exp}\{\mathrm{\Delta }Z^2\}j_{k_1^{}}\left(2\xi _1SZ\right)j_{k_1^{\prime \prime }}\left(2\xi _1^{}S^{}Z\right)j_{k_2^{}}\left(i2\xi _2SZ\right)j_{k_2^{\prime \prime }}\left(i2\xi _2^{}S^{}Z\right)P_{k_1,k_2,k_1^{},k_2^{},k_1^{\prime \prime },k_2^{\prime \prime }}`$
where:
$`P_{k_1,k_2,k_1^{},k_2^{},k_1^{\prime \prime },k_2^{\prime \prime }}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{2}{}}}\left(Y_{k_i}(\stackrel{ˇ}{𝐗}_i)Y_{k_i}(\stackrel{ˇ}{𝐗}_i^{})\right)\left(Y_{k_i^{}}(\stackrel{ˇ}{𝐗}_i)Y_{k_i^{}}(\stackrel{ˇ}{𝐙})\right)\left(Y_{k_i^{\prime \prime }}(\stackrel{ˇ}{𝐙})Y_{k_i^{\prime \prime }}(\stackrel{ˇ}{𝐗}_i^{})\right)`$ (141)
$`=`$ $`{\displaystyle \underset{(l_1,l_2)L;(l_1^{},l_2^{})L^{};L^{\prime \prime }}{}}𝒫_{\{k_1,k_2,k_1^{},k_2^{},k_1^{\prime \prime },k_2^{\prime \prime }\};\{(l_1,l_2)L;(l_1^{},l_2^{})L^{};L^{\prime \prime }\}}`$ (143)
$`(\{Y_{l_1}(\stackrel{ˇ}{𝐗}_1)\times Y_{l_2}(\stackrel{ˇ}{𝐗}_2)\}_L\{\{Y_{l_1^{}}(\stackrel{ˇ}{𝐗}_1^{})\times Y_{l_2^{}}(\stackrel{ˇ}{𝐗}_2^{})\}_L^{}\times Y{}_{L^{\prime \prime }}{}^{}(\stackrel{ˇ}{𝐙})\}_L)`$
The intermediate (but redundant) index $`L^{\prime \prime }`$ is connected to the integration variable $`𝐙`$ of the integral transformation (128). Because of the orthogonality between spherical harmonics and of the scalar character of the potential operator $`L`$ and $`L^{}`$ will be equal, and $`L^{\prime \prime }=0`$ after integration over $`𝐙`$: we therefore anticipate by simplifying (143) to
$`P_{k_1,k_2,k_1^{},k_2^{},k_1^{\prime \prime },k_2^{\prime \prime }}`$ $`=`$ $`{\displaystyle \underset{(l_1,l_2)L;(l_1^{},l_2^{})L}{}}𝒫_{\{k_1,k_2,k_1^{},k_2^{},k_1^{\prime \prime },k_2^{\prime \prime }\};\{(l_1,l_2)L;(l_1^{},l_2^{})L;0\}}{\displaystyle \frac{1}{\sqrt{4\pi }}}`$ (145)
$`\left(\left\{Y_{l_1}(\stackrel{ˇ}{𝐗}_1)\times Y_{l_2}(\stackrel{ˇ}{𝐗}_2)\right\}_L\times \left\{Y_{l_1^{}}(\stackrel{ˇ}{𝐗}_1^{})\times Y_{l_2^{}}(\stackrel{ˇ}{𝐗}_2^{})\right\}_L\right)`$
The expansion coefficients are then easily shown to be
$`𝒫_{\{k_1,k_2,k_1^{},k_2^{},k_1^{\prime \prime },k_2^{\prime \prime }\};\{(l_1,l_2)L;(l_1^{},l_2^{})L;0\}}={\displaystyle \underset{k}{}}(1)^{l_1+l_1^{}+k+L}\left\{\begin{array}{ccc}l_2^{}& l_2& k\\ l_1& l_1^{}& L\end{array}\right\}C_{k0k0}^{00}`$ (148)
$`{\displaystyle \underset{i=1}{\overset{2}{}}}\left(1\right)^{k_1}\left(2k_i+1\right)\left(2k_i^{}+1\right)\left(2k_i^{\prime \prime }+1\right)\left\{\begin{array}{ccc}k_i^{\prime \prime }& l_i^{}& k_i\\ l_i& k_i^{}& k\end{array}\right\}C_{k_i^{}0k_i0}^{l_i0}C_{k_i^{\prime \prime }0k_i0}^{l_i^{}0}C_{k_i^{}0k_i^{\prime \prime }0}^{k0}`$ (151)
Under the anticipative assumptions $`L=L^{}`$ and $`L^{\prime \prime }=0`$, the integration over the angles $`\stackrel{ˇ}{𝐙}`$ in (140) is now trivial. The remaining integration over $`Z`$ is easily done after substituting the power expansions for the Bessel functions $`j_l(x)`$ and $`i_l(x)`$. The final result provides a tractable though very bulky result for $`𝒱_{\left(l_1l_2\right)LM;(l_1l_2)LM}(ϵ,S;ϵ^{},S^{})`$ that is not reproduced here, as it carries no further additional information.
Again, by using the standard procedures such as differentiation or recurrence relations one obtains the effective matrix elements. In particular the form of (92) remains valid, though the analytic differentiation preferably should be performed within an algebraic package such as Mathematica or Maple due to the bulkiness of the formulae.
### D Matrix elements for the Coulomb potential
The matrix elements of the Coulomb potential can now be most easily obtained from the Gaussian results. We consider again the Gaussian integral representation
$$\frac{1}{\left|𝐫_i𝐫_j\right|}=\frac{2}{\sqrt{\pi }}_0^{\mathrm{}}𝑑x\mathrm{exp}\{x^2(𝐫_i𝐫_j)^2\}$$
(152)
We use (140), replace $`a^2`$ by $`1/x^2`$, and introduce a new integration variable $`t=\gamma /(1\gamma )`$ where $`\gamma =x^2b^2`$, to obtain symbolically
$`V_C(ϵ,S,\stackrel{ˇ}{𝐗}_1,\stackrel{ˇ}{𝐗}_\mathrm{𝟐};ϵ^{},S^{},\stackrel{ˇ}{𝐗}_1^{},\stackrel{ˇ}{𝐗}_\mathrm{𝟐}^{})`$ (153)
$`=`$ $`{\displaystyle \frac{2}{\sqrt{\pi }}}{\displaystyle _0^{\mathrm{}}}𝑑xV(ϵ,S,\stackrel{ˇ}{𝐗}_1,\stackrel{ˇ}{𝐗}_\mathrm{𝟐};ϵ^{},S^{},\stackrel{ˇ}{𝐗}_1^{},\stackrel{ˇ}{𝐗}_\mathrm{𝟐}^{})`$ (154)
$`=`$ $`{\displaystyle \frac{2}{b\sqrt{\pi }}}{\displaystyle _0^1}{\displaystyle \frac{dt}{(1t)^{\frac{3}{2}}t^{\frac{1}{2}}}}V(ϵ,S,\stackrel{ˇ}{𝐗}_1,\stackrel{ˇ}{𝐗}_\mathrm{𝟐};ϵ^{},S^{},\stackrel{ˇ}{𝐗}_1^{},\stackrel{ˇ}{𝐗}_\mathrm{𝟐}^{})`$ (155)
Mutatis mutandis we apply the integral transformation directly to the generating matrix elements $`𝒱_{\left(l_1l_2\right)LM;(l_1^{}l_2^{})LM}(ϵ,S;ϵ^{},S^{})`$
$`𝒱_{\left(l_1l_2\right)LM;(l_1^{}l_2^{})LM}^C(ϵ,S;ϵ^{},S^{})`$ (156)
$`=`$ $`{\displaystyle \frac{2}{\sqrt{\pi }}}{\displaystyle _0^{\mathrm{}}}𝑑x𝒱_{\left(l_1l_2\right)LM;(l_1^{}l_2^{})LM}(ϵ,S;ϵ^{},S^{})`$ (157)
$`=`$ $`{\displaystyle \frac{2}{b\sqrt{\pi }}}{\displaystyle _0^1}{\displaystyle \frac{dt}{(1t)^{\frac{3}{2}}t^{\frac{1}{2}}}}𝒱_{\left(l_1l_2\right)LM;(l_1^{}l_2^{})LM}(ϵ,S;ϵ^{},S^{})`$ (158)
This however leads to a very intricate evaluation of the integral, and thus for the reduced generating matrix element of the Coulomb potential.
A better procedure is to generate the quantum numbers $`K`$ and $`K^{}`$ first, by differentiation of $`𝒱_{\left(l_1l_2\right)LM;(l_1^{}l_2^{})LM}(ϵ,S;ϵ^{},S^{})`$ on $`S`$ and $`S^{}`$, then setting $`S=S^{}=0`$. This leads to reduced generating functions for each $`K`$ and $`K^{}`$ now only dependent on $`ϵ`$ and $`ϵ^{}`$, which are of a much simpler form and allow for an analytic integration in (158). The further derivation on matrix elements for $`n`$ and $`n^{}`$ is then straightforward.
## VII Conclusion
In this paper we presented a framework for a microscopic three-cluster model within the Algebraic Model for scattering. It was shown that it is possible to obtain matrix elements for fully antisymmetrized three-cluster configurations, as well as a proper description for the the three-cluster continuum in terms of a hyperspherical description. The corresponding AM equations in a multichannel description were also introduced. In the current work the individual clusters were limited to contain only s-orbitals, thus reducing the mass of clussters to that of a four-nucleon system. The latter restriction is however not a fundamental one, and was taken to restrict the analytical and calculational burden.
In order to prove the validity and feasibility of the current model we will apply it to two specific three-cluster configurations believed to be of importance to astrophysical physiscs, $`\alpha +n+n`$ for <sup>6</sup>He, and $`\alpha +p+p`$ for <sup>6</sup>Be. These results appear in part II of this paper.
## VIII Acknowledgments
This work was partially supported by an INTAS Grant-(“INTAS-93-755-extension). One of the authors (V. V.) gratefully acknowledges the University of Antwerp (RUCA) for a “RAFO-gastprofessoraat 1998-1999” and the kind hospitality of the members of the research group “Computational Quantum Physics” of the Department of “Mathematics and Computer Sciences”, University of Antwerp, RUCA, Belgium.
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# Untitled Document
Proof of the Marginal Stability
Bound for the Swift-Hohenberg Equation
and Related Equations
P. Collet<sup>1</sup> and J.-P. Eckmann<sup>2,3</sup>
<sup>1</sup>Centre de Physique Théorique, Laboratoire CNRS UMR 7644, Ecole Polytechnique, F-91128 Palaiseau Cedex, France
<sup>2</sup>Dépt. de Physique Théorique, Université de Genève, CH-1211 Genève 4, Switzerland
<sup>3</sup>Section de Mathématiques, Université de Genève, CH-1211 Genève 4, Switzerland
Abstract. We prove that if the initial condition of the Swift-Hohenberg equation
$$_tu(x,t)=\left(\epsilon ^2\left(1+_x^2\right)^2\right)u(x,t)u^3(x,t)$$
is bounded in modulus by $`Ce^{\beta x}`$ as $`x+\mathrm{}`$, the solution cannot propagate to the right with a speed greater than
$$\underset{0<\gamma \beta }{sup}\gamma ^1\left(\epsilon ^2+4\gamma ^2+8\gamma ^4\right).$$
This settles a long-standing conjecture about the possible asymptotic propagation speed of the Swift-Hohenberg equation. The proof does not use the maximum principle and is simple enough to generalize easily to other equations. We illustrate this with an example of a modified Ginzburg-Landau equation, where the minimal speed is not determined by the linearization alone.
1. Introduction The marginal stability conjecture deals with the possible propagation speed of solutions of dissipative partial differential equations. It was formulated in the late 1970’s by several authors. Its clearest form is obtained for the Ginzburg-Landau equation
$$_tu(x,t)=_x^2u(x,t)+u(x,t)u^3(x,t),$$
$`(1.1)`$
where $`u:𝐑\times 𝐑^+𝐑`$. When the initial data have compact support, then the solution cannot propagate with a speed faster than some critical speed $`c`$, which happens to be 2 for this example. The number 2 can be understood as follows. One writes $`u(x,t)=v(xct)`$, and looks for a solution of (1.1) expressed for $`v`$:
$$0=_x^2v+c_xv+vv^3.$$
$`(1.2)`$
If one makes the assumption that $`v(\xi )=C_1e^{\beta \xi }`$ as $`\xi +\mathrm{}`$, one finds obviously that $`\beta `$ and $`c`$ should be related through the equation
$$0=\beta ^2\beta c+1,$$
$`(1.3)`$
since the non-linear term is irrelevant at $`\xi =\mathrm{}`$ in this case. For fixed $`\beta `$ we clearly find $`c=(\beta ^2+1)/\beta `$, and since functions which are (in absolute value) bounded by $`C\mathrm{exp}(\beta x)`$ are also bounded by $`C^{}\mathrm{exp}(\beta ^{}x)`$ for $`0<\beta ^{}<\beta `$ one finds in this case an upper bound
$$c_\beta ^{GL}=\underset{0<\gamma <\beta }{inf}\frac{\gamma ^2+1}{\gamma },$$
$`(1.4)`$
and this is equal to 2 for $`\beta 1`$. Using the maximum principle for parabolic PDE’s, Aronson and Weinberger were able to show \[AW\] that no positive solution starting from initial conditions with compact support can move faster than the speed $`c_2^{GL}=2`$. Using essentially the same argument, it was also shown in \[CE\] that if the initial condition decays like $`e^{\beta x}`$ with $`\beta <1`$, then the solution cannot move faster than $`c_\beta ^{GL}`$. However, in cases where the maximum principle does not apply, such as in $`\mathrm{}`$SH, the maximum possible speed was only conjectured, and tested numerically, but no rigorous result was obtained, see e.g., \[LMK, DL, BBDKL\].
In a somewhat different direction, there is the important, and difficult, issue on whether there is actually a solution moving with the maximal allowed velocity. In general, its realization depends on the details of the nonlinearity, and this question has been extensively discussed in the literature \[AW, B, DL, BBDKL, vS\]. It will not be treated here.
The main result of our paper is an upper bound on the speed of propagation of solutions to the Swift-Hohenberg equation
$$_tu=\left(\epsilon (1+_x^2)^2\right)uu^3.$$
$`(1.5)`$
The polynomial equation analogous to (1.3) turns out to be
$$0=\epsilon ^2+4\beta ^2+8\beta ^4c\beta ,$$
$`(1.6)`$
and we define in this case
$$c_\beta =c_\beta ^{SH}=\underset{0<\gamma \beta }{inf}\frac{\epsilon ^2+4\gamma ^2+8\gamma ^4}{\gamma }.$$
$`(1.7)`$
The polynomial is an absolute maximum of the real part of $`P`$, as we explain at the end of the introduction and in the Appendix. This will be the minimal speed.While it looks different from the standard discussion in \[BBDKL\], we explain in the Appendix that the two definitions coincide. The current formulation has the advantage of being expressed in terms of real variables, although the traveling wave in this case is actually modulated \[CE\]. Our result can be expressed informally as follows: If the initial data for the problem are bounded in absolute value by $`Ce^{\beta x}`$ as $`x+\mathrm{}`$ then the solution cannot advance faster to the right than $`c_\beta `$ in the sense that
$$\underset{t\mathrm{}}{lim}u(x+ct,t)=\mathrm{\hspace{0.17em}0},$$
for all $`c>c_\beta `$. In particular, if the initial condition has compact support, the above hypotheses are satisfied for any $`\beta >0`$ and we find an upper bound on the speed which is $`c_{}=inf_\beta c_\beta `$: This is the absolute minimum of $`(\epsilon ^2+4\beta ^2+8\beta ^4)/\beta `$.
Remark. The precise formulation is given in $`\mathrm{}`$main.
Before explaining the main steps of the proof we note a well-known result, namely that if the initial condition $`u_0`$ is bounded in $`𝒞^3`$, i.e.,
$$\underset{j=0,\mathrm{},3}{\mathrm{max}}\underset{x𝐑}{sup}|_x^ju_0(x)|K,$$
$`(1.8)`$
then there is a constant $`L=L(K)`$ such that for all $`t>0`$ one has
$$\underset{j=0,\mathrm{},3}{\mathrm{max}}\underset{x𝐑}{sup}|_x^ju(x,t)|L(K).$$
$`(1.9)`$
The proof of the main result is really quite easy and consists of 3 steps:
i) An a priori bound on the Green’s function of the semigroup generated by the linear part $`\epsilon ^2(1+_x^2)^2`$ of the Swift-Hohenberg equation.
ii) The observation that if the initial condition satisfies $`lim_x\mathrm{}e^{\beta x}_x^ju_0(x)=0`$, for $`j=0,\mathrm{},3`$, then the same holds for $`u(x,t)`$. This is needed later on to ensure that integration by parts does not produce boundary terms at infinity.
iii) An energy-like estimate which shows that
$$\underset{t\mathrm{}}{lim}_{ct}^{\mathrm{}}𝑑x|u(x,t)|^2e^{2\beta (xct)}=\mathrm{\hspace{0.17em}0},$$
when $`c>c_\beta `$ (if it is finite at $`t=0`$, see below for details). Thus, the solution is outrun by a frame moving with speed $`c>c_\beta `$. While this is similar to what was observed in the proofs where the maximum principle could be used, it has here a quite different origin of dynamical nature.
In $`\mathrm{}`$nonlin , we consider the case of the Ginzburg-Landau equation when the nonlinearity $`uu^3`$ is replaced by a general function $`f(u)`$ with the properties $`f(0)=0`$, $`0<f^{}(0)<\mathrm{}`$ and $`lim\; sup_z\mathrm{}f(z)/z<0`$. In such a case, the bound (1.4) is replaced by
$$c_\beta ^{GL^{}}=\underset{0<\gamma <\beta }{inf}\frac{\gamma ^2+\underset{u}{sup}\frac{f(u)}{u}}{\gamma }.$$
In the case of the Swift-Hohenberg equation the bound generalizes as follows: Assume the equation is
$$_tu=(\epsilon ^2(1+_x^2)^2))u+f(u).$$
Then we get for the maximal possible speed:
$$c_\beta ^{SH^{}}=\underset{0<\gamma \beta }{inf}\frac{\epsilon ^2+4\gamma ^2+8\gamma ^4+\underset{u}{sup}\frac{f(u)}{u}}{\gamma }.$$
In an appendix, we show that the expression (1.6) is nothing but
$$\underset{k_\beta ^{}}{sup}ReP(z)|_{z=\beta +ik_\beta ^{}},$$
where the sup is over the solutions $`k_\beta ^{}`$ of
$$\frac{dReP(\beta +ik)}{dk}=\mathrm{\hspace{0.17em}0}.$$
We also show that these conditions are the same as those found in \[BBDKL\].
Finally, it should be noted that the method is not restricted to 1-dimensional problems, and can also be applied to questions of grows of ‘‘bubbles’’ in the 2-dimensional Swift-Hohenberg equation.
2. A pointwise bound on the Green’s function Here we bound the Green’s function of the operator $`\epsilon ^2(1+_x^2)^2`$ by a method which generalizes immediately to other problems of similar type. Let $`P`$ be a polynomial in $`k`$ which is of the form
$$P(ik)=a_nk^n+\underset{m=0}{\overset{n1}{}}a_mk^ma_nk^n+R(k),$$
and assume $`n`$ even and $`a_n>0`$. (For the Swift-Hohenberg equation, $`P(z)=\epsilon ^2(1+z^2)^2`$.) Then the Green’s function
$$G_t(x)=𝑑ke^{ikx}e^{P(ik)t},$$
satisfies:
Lemma 2.1. Given $`0<\beta <\mathrm{}`$, there is a constant $`C(\beta )`$ such that for all $`t(0,1]`$ one has the bound
$$t^{1/n}|G_t(x)|e^{(\beta ^{}+2t^{1/n})|x|}C(\beta ),$$
$`(2.1)`$
for all $`\beta ^{}[0,\beta ]`$.
Remark. This clearly also implies, for all $`t(0,1]`$ and all $`\beta ^{}[0,\beta ]`$:
$$𝑑x|G_t(x)|e^{\beta ^{}|x|}C(\beta ),$$
$`(2.2)`$
since $`𝑑x|G_t(x)|e^{\beta ^{}|x|}C(\beta )𝑑xt^{1/n}e^{2|x|t^{1/n}}C(\beta )`$. Proof. We will show the bound in the form
$$t^{1/n}|G_t(zt^{1/n})|e^{\gamma t^{1/n}z}C(\beta ),$$
$`(2.3)`$
with $`\gamma =\beta +2t^{1/n}`$, and it clearly suffices to consider $`z>0`$. Proving (2.3) is a straightforward calculation which is probably well-known. Indeed, the l.h.s. of (2.3) equals (without the absolute values)
$$\begin{array}{cc}& 𝑑kt^{1/n}\mathrm{exp}\left(\gamma t^{1/n}z+ikt^{1/n}za_nk^nt+R(k)t\right)\hfill \\ & =𝑑\mathrm{}\mathrm{exp}\left(\gamma t^{1/n}z+i\mathrm{}za_n\mathrm{}^n+R(\mathrm{}t^{1/n})t\right).\hfill \end{array}$$
Since the integrand is an entire function in $`\mathrm{}`$ we can shift the contour from $`\mathrm{}`$ to $`\mathrm{}^{}=\mathrm{}i\gamma t^{1/n}`$ and the last expression is seen to be equal to
$$𝑑\mathrm{}^{}\mathrm{exp}\left(i\mathrm{}^{}za_n(\mathrm{}^{}+i\gamma t^{1/n})^n+R(\mathrm{}^{}t^{1/n}+i\gamma )t\right).$$
Note now that
$$\begin{array}{ccc}\hfill |& \mathrm{exp}(i\mathrm{}^{}za_n(\mathrm{}^{}+i\gamma t^{1/n})^n+R(\mathrm{}^{}t^{1/n}+i\gamma )t)|\hfill & \\ & =\left|\mathrm{exp}\left(a_n(\mathrm{}^{}+i\gamma t^{1/n})^n+R(\mathrm{}^{}t^{1/n}+i\gamma )t\right)\right|,\hfill & (2.4)\hfill \end{array}$$
and for bounded $`\beta `$ and $`t(0,1]`$ we find that $`\gamma t^{1/n}=(\beta +2t^{1/n})t^{1/n}\beta +2`$, and hence (2.4) is uniformly integrable in $`\mathrm{}^{}`$, since $`a_n>0`$. The proof of Lemma 2.1 is complete.
3. Exponential decay of solutions In this section, we prove a bound in the laboratory frame, showing that if the initial condition goes exponentially to 0 then the solution at time $`t`$ goes to zero as well, with the same rate.
Theorem 3.1. Assume that $`u_0`$ is bounded in $`𝒞^3`$ and that
$$\underset{x\mathrm{}}{lim}e^{\beta x}_x^ju_0(x)=\mathrm{\hspace{0.17em}0},$$
$`(3.1)`$
for $`j=0,\mathrm{},3`$ and some $`\beta >0`$. Then the solution $`u(x,t)`$ of (1.5) with initial data $`u_0`$ satisfies for all $`t>0`$:
$$\underset{x\mathrm{}}{lim}e^{\beta x}_x^ju(x,t)=\mathrm{\hspace{0.17em}0},$$
$`(3.2)`$
for $`j=0,\mathrm{},3`$.
Proof. The proof is in steps of some (fixed) time $`\tau _{}`$. We define first
$$g_\xi (x)=\mathrm{\hspace{0.17em}1}+e^{\beta (x\xi )}.$$
The assumption means that $`u_0`$ satisfies (1.8) for some $`K`$. From (3.1), and because $`L(K)K`$, we conclude that there is a $`\xi >0`$ for which
$$\underset{x𝐑}{sup}g_\xi (x)|_x^ju_0(x)|\mathrm{\hspace{0.17em}2}L(K),$$
$`(3.3)`$
for $`j=0,\mathrm{},3`$. Note that we do not have any control on the size of $`\xi `$, but such a control is not needed.
From (1.8) we also conclude (see (1.9)) that
$$\underset{t0}{sup}\underset{x𝐑}{sup}|_x^ju(x,t)|L(K),$$
$`(3.4)`$
for $`j=0,\mathrm{},3`$.
The crucial step in the proof of Theorem 3.1 is
Lemma 3.2. There are a $`\tau _{}>0`$ and a $`\rho `$, independent of $`\xi `$, such that for $`t[0,\tau _{}]`$ one has
$$\underset{j=0,\mathrm{},3}{sup}\underset{x𝐑}{sup}g_\xi (x)|_x^ju(x,t)|\rho .$$
$`(3.5)`$
Proof. We use the estimates on the convolution kernel $`G_t`$ associated with the semigroup $`t\mathrm{exp}(t(\epsilon ^2(1+_x^2)^2)`$ which were proven in Section 2. One has
$$u_t=G_tu_0_0^t𝑑sG_{ts}u_s^3,$$
where $`u_s(x)=u(x,s)`$. We define $`_\xi `$ as the space of uniformly continuous functions $`f`$ for which
$$f_\xi =\underset{x𝐑}{sup}g_\xi (x)|f(x)|<\mathrm{}.$$
Using this quantity as a norm makes $`_\xi `$ a Banach space. Consider next the space $`𝒦=𝒦_{\xi ,\tau _{}}=𝒞^0([0,\tau _{}],_\xi )`$ of functions $`h:(x,t)h(x,t)`$, with the norm
$$h_{\xi ,\tau _{}}=\underset{t[0,\tau _{}]}{sup}h(,t)_\xi .$$
This is again a Banach space. For $`v𝒦`$ we define the map $`v𝒬v`$ by
$$\left(𝒬v\right)(x,t)=\left(G_tu_0\right)(x)_0^t𝑑s\left(G_{ts}v_s^3\right)(x).$$
$`(3.6)`$
Note that if $`𝒬v=v`$, then $`v`$ is a solution to (1.5) with initial condition $`u_0`$. To find $`v`$, we will show that for sufficiently small $`\tau _{}>0`$ the operator $`𝒬`$ contracts a small ball of $`𝒦_{\xi ,\tau _{}}`$ to itself. The center of this ball is the function $`(x,t)0`$.
First we bound $`G_tu_0`$. Note that from the definition of $`g_\xi `$ we find
$$\frac{g_\xi (x)}{g_\xi (y)}e^{\beta |xy|},$$
since for $`x<y`$ the quotient is bounded by 1 and for $`x>y`$ we have the (very rough) bound $`e^{\beta (xy)}`$. From Lemma 2.1, we have for all $`t(0,1]`$ and all $`x𝐑`$:
$$|G_t(x)|e^{2\beta |x|}C(\beta )t^{1/4}e^{2|x|t^{1/4}},$$
$`(3.7)`$
and, clearly, $`C(\beta )`$ can be chosen the same value for all smaller $`\beta `$. Using this, we find
$$\begin{array}{cc}\hfill |\left(G_tu_0\right)(x)g_\xi (x)|& 𝑑y|G_t(xy)u_0(y)|g_\xi (y)\frac{g_\xi (x)}{g_\xi (y)}\hfill \\ & 𝑑y|G_t(xy)u_0(y)|g_\xi (y)e^{\beta |xy|}\hfill \\ & 𝑑z|G_t(z)e^{\beta |z|}|\underset{z^{}𝐑}{sup}|u_0(z^{})|g_\xi (z^{})\hfill \\ & C(\beta )\underset{z^{}𝐑}{sup}|u_0(z^{})|g_\xi (z^{}).\hfill \end{array}$$
$`(3.8)`$
Combining these bounds with (3.3) we get
$$|\left(G_tu_0\right)(x)g_\xi (x)|C_2L(K).$$
In fact, we can do a little better in (3.8) by extracting a factor of $`e^{\beta |xy|}`$. The last two lines in (3.8) are replaced by
$$\begin{array}{cc}\hfill |(& G_tu_0)(x)g_\xi (x)|\hfill \\ & 𝑑z|G_t(z)e^{2\beta |z|}|\underset{y𝐑}{sup}|u_0(y)|g_\xi (y)e^{\beta |xy|}\hfill \\ & C(2\beta )\underset{y𝐑}{sup}|u_0(y)|g_\xi (y)e^{\beta |xy|}.\hfill \end{array}$$
$`(3.9)`$
Since $`|u_0(y)|g_\xi (y)`$ is bounded and converges to 0 as $`y+\mathrm{}`$, we conclude that the quantity in (3.9) tends to 0 as $`x+\mathrm{}`$. Thus, we also have
$$\underset{x\mathrm{}}{lim}|\left(G_tu_0\right)(x)g_\xi (x)|=\mathrm{\hspace{0.17em}0}.$$
$`(3.10)`$
We next bound the non-linear term. Assume $`v𝒦_{\xi ,\tau _{}}`$ and $`v_{\xi ,\tau _{}}<\rho `$. Then any power ($`1`$) of $`v`$ is also in $`𝒦_{\xi ,\tau _{}}`$ and one has a bound of the form
$$v^3_{\xi ,\tau _{}}C_3\rho ^3.$$
Therefore, the method leading to (3.8) now yields
$$\left|_0^t𝑑s\left(G_{ts}v_s^3\right)(x)g_\xi (x)\right|C_4\rho ^3t,$$
and if also $`w_{\xi ,\tau _{}}<\rho `$, then a variant of that method gives:
$$\begin{array}{cc}\hfill |_0^tds(G_{ts}& v_s^3)(x)g_\xi (x)_0^tds(G_{ts}w_s^3)(x)g_\xi (x)|\hfill \\ & C_5\rho ^2t\underset{s[0,t]}{sup}\underset{x𝐑}{sup}|v_s(x)w_s(x)|g_\xi (x).\hfill \end{array}$$
Taking the center of the ball at $`(x,t)0`$ and the radius $`\rho =2C_2K`$ and then $`\tau _{}<\mathrm{min}\{(4C_4\rho ^3)^1,(4C_5\rho ^2)^1\}`$, we have a contraction and hence a unique fixed point $`v`$ for $`𝒬`$. For $`j=1,2,3`$, we use the same methods since we can push all derivatives from the operator $`G_t`$ to the function $`v`$, because $`G_t`$ is a convolution. The details are left to the reader. The existence of this fixed point clearly shows Lemma 3.2.
We come back to the proof of Theorem 3.1. We define
$$\mathrm{\Gamma }(t)=\underset{x\mathrm{}}{lim\; sup}|u(x,t)g_\xi (x)|.$$
By assumption, we have $`\mathrm{\Gamma }(0)=0`$ and by Lemma 3.2 we have
$$|u(x,t)|\rho /g_\xi (x),$$
so that $`\mathrm{\Gamma }(t)\rho `$ for $`t\tau _{}`$. We now show it is actually 0 for those $`t`$. Consider $`𝒬`$ as in (3.6). Note that
$$\begin{array}{cc}\hfill \mathrm{\Gamma }(t)& =\underset{x\mathrm{}}{lim\; sup}|u(x,t)g_\xi (x)|\hfill \\ & =\underset{x\mathrm{}}{lim\; sup}g_\xi (x)|\left(G_tu_0\right)(x)|+\underset{x\mathrm{}}{lim\; sup}g_\xi (x)|_0^t𝑑s\left(G_{ts}u_s^3\right)(x)|.\hfill \end{array}$$
The first term vanishes by (3.9). Thus, $`\mathrm{\Gamma }`$ only depends on the nonlinear part. Using (3.7), that part can be bounded as
$$\begin{array}{cc}\hfill g_\xi (x)\left|_0^t𝑑sG_{ts}u_s^3(x)\right|& _0^t𝑑s𝑑y\frac{g_\xi (x)}{g_\xi ^3(y)}\left|G_{ts}(xy)\right|\left|g_\xi (y)u_s(y)\right|^3\hfill \\ & _0^t𝑑s𝑑y\left|G_{ts}(xy)\right|e^{\beta |xy|}\left|g_\xi (y)u_s(y)\right|^3\hfill \\ & C(\beta )_0^t𝑑s𝑑z(ts)^{1/4}e^{2|z|(ts)^{1/4}}\hfill \\ & \left|g_\xi (xz)u_s(xz)\right|^3.\hfill \end{array}$$
$`(3.11)`$
We need an upper bound for the $`lim\; sup_x\mathrm{}`$ of this expression. Fix an $`\epsilon >0`$. For $`s[0,t]`$, we can find an $`\eta (s,\epsilon )>0`$ such that
$$\underset{y\eta (s,\epsilon )}{sup}|g_\xi (y)u_s(y)|\mathrm{\Gamma }(s)+\epsilon .$$
There is also a number $`\zeta (\epsilon )>0`$ such that for any $`s[0,t]`$:
$$_{|z|>\zeta (\epsilon )}𝑑z(ts)^{1/4}e^{2|z|(ts)^{1/4}}\epsilon .$$
If $`x>\zeta (\epsilon )+\eta (s,\epsilon )`$, we have
$$𝑑z(ts)^{1/4}e^{2|z|(ts)^{1/4}}|g_\xi (xz)u_s(xz)|^3(\mathrm{\Gamma }(s)+\epsilon )^3+\rho ^3\epsilon ,$$
by Lemma 3.2. We cannot conclude directly by integration over $`s`$ because $`\eta `$ depends on $`s`$. However, $`\eta (s,\epsilon )`$ is finite for almost every $`s`$ (in reality for every $`s`$). Therefore, we can find a finite number $`\mathrm{\Theta }(\epsilon )`$ such that the set
$$E(\epsilon )=\left\{s[0,t]\right|\eta (s,\epsilon )>\mathrm{\Theta }(\epsilon )\}$$
has Lebesgue measure at most $`\epsilon ^2`$ (note that $`E(\epsilon )`$ is measurable). Therefore, if $`x>\mathrm{\Theta }(\epsilon )+\zeta (\epsilon )`$ we have
$$\begin{array}{cc}\hfill _0^t& ds𝑑z(ts)^{1/4}e^{2|z|(ts)^{1/4}}|g_\xi (xz)u_s(xz)|^3\hfill \\ & =_{([0,t]E(\epsilon ))E(\epsilon )}𝑑s𝑑z(ts)^{1/4}e^{2|z|(ts)^{1/4}}|g_\xi (xz)u_s(xz)|^3\hfill \\ & C_6_0^t𝑑s\left((\mathrm{\Gamma }(s)+\epsilon )^3+\rho ^3\epsilon \right)+C_7\rho ^3_{E(\epsilon )}𝑑s(ts)^{1/4}.\hfill \end{array}$$
The last integral is of order $`\epsilon ^{1/2}`$ by the Schwarz inequality. Since $`\epsilon >0`$ is arbitrary, we get
$$\mathrm{\Gamma }(t)C_8_0^t𝑑s\mathrm{\Gamma }(s)^3.$$
Since $`\mathrm{\Gamma }`$ is bounded by what we said above and $`\mathrm{\Gamma }(0)=0`$, it follows from Gronwall’s lemma that $`\mathrm{\Gamma }(t)=0`$ for $`t\tau _{}`$. One then repeats the argument for all consecutive intervals of length $`\tau _{}`$. The proof of the corresponding bounds on the derivatives is similar and is left to the reader.
4. Bound on the speed We define $`J_\xi `$ by
$$J_\xi (t)=_\xi ^{\mathrm{}}𝑑x|u(x,t)|^2e^{2\beta (x\xi )},$$
$`(4.1)`$
where $`u(x,t)`$ is the solution of the Swift-Hohenberg equation. The main result of this paper is
Theorem 4.1. Let $`u(x,t)`$ be a solution of the Swift-Hohenberg equation (1.5) for an initial condition $`u_0(x)=u(x,0)`$ which is in $``$, which satisfies $`J_0(0)<\mathrm{}`$ for some $`\beta >0`$ and which satisfies the assumptions of Theorem 3.1. Then one has
$$\underset{t\mathrm{}}{lim}_{ct}^{\mathrm{}}𝑑x|u(x,t)|^2e^{2\beta (xct)}=\mathrm{\hspace{0.17em}0},$$
$`(4.2)`$
for all $`c>\left(\epsilon ^2+4\beta ^2+8\beta ^4\right)/\beta `$.
Remark. If one is willing to pay a price of slightly more complicated formulations and proofs, one can omit the condition on $`J_0(0)`$ in Theorem 4.1. One would then assume the pointwise bounds of Theorem 3.1 fore some $`\beta >0`$ and work throughout the proof with a $`J_\xi (t)`$ defined with some $`\beta ^{}<\beta `$, but arbitrarily close to it, since the condition on $`c`$ is open.
Proof. We define $`v_\xi (x,t)=u(x,t)e^{\beta (x\xi )}`$, so that $`J_\xi (t)=_\xi ^{\mathrm{}}𝑑x|v_\xi (x,t)|^2`$, and $`v_\xi `$ solves the equation
$$_tv_\xi (x,t)=\epsilon ^2v_\xi (x,t)\left(1+(_x\beta )^2\right)^2v_\xi (x,t)v_\xi ^3(x,t)e^{2\beta (x\xi )}.$$
$`(4.3)`$
Since $`u`$ is real, the absolute values in the definition of $`J_\xi (t)`$ can be omitted. Differentiating (4.1) with respect to time, we get
$$\begin{array}{cc}\hfill \frac{1}{2}_tJ_\xi (t)& =_\xi ^{\mathrm{}}𝑑xv_\xi (x,t)_tv_\xi (x,t).\hfill \end{array}$$
Since $`\xi `$ is fixed throughout the calculation, we omit the index of $`v_\xi `$. We also omit the arguments $`(x,t)`$. Note that by Theorem 3.1, $`lim_x\mathrm{}_x^jv_\xi (x,t)=0`$, for $`j=0,\mathrm{},3`$, so that we can freely integrate by parts in the following calculation. We find, using $`_xv=v^{}`$:
$$\begin{array}{cc}\hfill \frac{1}{2}_tJ_\xi (t)& =_\xi ^{\mathrm{}}𝑑xv\left(\epsilon ^2v\left(1+(_x\beta )^2\right)^2vv^3e^{2\beta (x\xi )}\right)\hfill \\ & =_\xi ^{\mathrm{}}𝑑xv\left(\epsilon ^2v\left(1+_x^22\beta _x+\beta ^2\right)^2vv^3e^{2\beta (x\xi )}\right)\hfill \\ & =_\xi ^{\mathrm{}}dxv(\epsilon ^2v_x^4v+4\beta _x^3v2(1+3\beta ^2)_x^2v+4\beta (1+\beta ^2)_xv(1+\beta ^2)^2v\hfill \\ & v^3e^{2\beta (x\xi )})\hfill \\ & =_\xi ^{\mathrm{}}dx((\epsilon ^2(1+\beta ^2)^2)v^2e^{2\beta (x\xi )}v^42(1+3\beta ^2)vv^{\prime \prime }\hfill \\ & +v^{}v^{\prime \prime \prime }4\beta v^{}v^{\prime \prime })\hfill \\ & +\left(vv^{\prime \prime \prime }4\beta vv^{\prime \prime }2\beta (1+\beta ^2)v^2\right)|_{x=\xi ,t}.\hfill \end{array}$$
We integrate by parts some more and get
$$\begin{array}{cc}\hfill \frac{1}{2}_tJ_\xi (t)& =_\xi ^{\mathrm{}}𝑑x\left(\left(\epsilon ^2(1+\beta ^2)^2\right)v^2e^{2\beta (x\xi )}v^42(1+3\beta ^2)v^{\prime \prime }v(v^{\prime \prime })^2\right)\hfill \\ & +\left(vv^{\prime \prime \prime }4\beta vv^{\prime \prime }2\beta (1+\beta ^2)v^2v^{}v^{\prime \prime }+2\beta (v^{})^2\right)|_{x=\xi ,t}.\hfill \end{array}$$
$`(4.4)`$
We write $`B_\xi (t)`$ for the boundary term obtained above:
$$B_\xi (t)=\left(vv^{\prime \prime \prime }4\beta vv^{\prime \prime }2\beta (1+\beta ^2)v^2v^{}v^{\prime \prime }+2\beta (v^{})^2\right)|_{x=\xi ,t}.$$
Finally, we rewrite (4.4) by completing a square:
$$\begin{array}{cc}\hfill \frac{1}{2}_tJ_\xi (t)& =_\xi ^{\mathrm{}}dx((\epsilon ^2(1+\beta ^2)^2+(1+3\beta ^2)^2)v^2e^{2\beta (x\xi )}v^4\hfill \\ & (v^{\prime \prime }+(1+3\beta ^2)v)^2)+B_\xi (t).\hfill \end{array}$$
$`(4.5)`$
Note that (4.5) leads immediately to a differential inequality:
$$\frac{1}{2}_tJ_\xi (t)G(\beta )J_\xi (t)+B_\xi (t),$$
$`(4.6)`$
with
$$G(\beta )=\epsilon ^2(1+\beta ^2)^2+(1+3\beta ^2)^2=\epsilon ^2+4\beta ^2+8\beta ^4.$$
$`(4.7)`$
This is the origin of the polynomial in (1.7). We bound first the boundary term.
Lemma 4.2. There is a $`C_9`$ such that for all $`u_0`$, all $`\xi `$, and all $`t>0`$ one has
$$B_\xi (t)C_9.$$
$`(4.8)`$
Proof. Recall that $`v_\xi (x,t)=e^{\beta (x\xi )}u(x,t)`$. Using elementary calculus, we find
$$_x^jv_\xi (x,t)=\underset{k=0}{\overset{j}{}}\left(\genfrac{}{}{0pt}{}{j}{k}\right)\beta ^je^{\beta (x\xi )}_x^{jk}u(x,t).$$
Therefore,
$$_x^jv_\xi (\xi ,t)=\underset{k=0}{\overset{j}{}}\left(\genfrac{}{}{0pt}{}{j}{k}\right)\beta ^j_\xi ^{jk}u(x,t)|_{x=\xi },$$
and the assertion follows because $`u`$.
Using Lemma 4.2, we conclude from (4.6) that
$$_tJ_\xi (t)\mathrm{\hspace{0.17em}2}G(\beta )J_\xi (t)+2C_9.$$
Solving the differential inequality from $`t`$ to $`t^{}`$, we obtain for $`t^{}>t`$,
$$J_\xi (t^{})e^{2G(\beta )(t^{}t)}J_\xi (t)+2\frac{e^{2G(\beta )(t^{}t)}1}{2G(\beta )}C_9.$$
$`(4.9)`$
We need this inequality in a slightly different form. Note that for $`\xi ^{}>\xi `$, one has
$$\begin{array}{cc}\hfill J_\xi ^{}(t)=_\xi ^{}^{\mathrm{}}𝑑xu^2(x,t)e^{2\beta (x\xi ^{})}& =e^{2\beta (\xi ^{}\xi )}_\xi ^{}^{\mathrm{}}𝑑xe^{2\beta (x\xi )}u^2(x,t)\hfill \\ & e^{2\beta (\xi ^{}\xi )}_\xi ^{\mathrm{}}𝑑xe^{2\beta (x\xi )}u^2(x,t)\hfill \\ & =e^{2\beta (\xi ^{}\xi )}J_\xi (t).\hfill \end{array}$$
$`(4.10)`$
Combining this with (4.9) we get for $`\xi ^{}>\xi `$ and $`t^{}>t`$:
$$J_\xi ^{}(t^{})e^{2\beta (\xi ^{}\xi )}\left(e^{2G(\beta )(t^{}t)}J_\xi (t)+\frac{e^{2G(\beta )(t^{}t)}1}{G(\beta )}C_9\right).$$
$`(4.11)`$
To complete the proof of Theorem 4.1, it suffices to set $`\xi ^{}=c\tau `$, $`t^{}=\tau `$, $`\xi =0`$ and $`t=0`$ in (4.11). Then we get
$$J_{c\tau }(\tau )e^{2(G(\beta )\beta c)\tau }\left(J_0(0)+\frac{C_9}{G(\beta )}\right).$$
$`(4.12)`$
Clearly, if $`c>G(\beta )/\beta `$, then $`J_{c\tau }(\tau )0`$ as $`\tau \mathrm{}`$. Thus, if $`J_0(0)<\mathrm{}`$, and the assertion of Theorem 4.1 follows.
Remark. One can do a little better than (4.12). Namely, consider the case where $`c=G(\beta )/\beta `$, that is, the case of a critical speed. Then one finds from (4.11) that
$$J_{c\tau +\lambda }(\tau )e^{2\beta \lambda }\left(J_0(0)+\frac{C_9}{G(\beta )}\right),$$
and in particular $`lim_\lambda \mathrm{}J_{c\tau +\lambda }(\tau )=0`$, if $`J_0(0)`$ is finite. This means that in the frame moving with exactly the critical speed, no amplitude ‘‘leaks’’ far ahead in that frame in $`L^2(e^{2\beta x}dx)`$. One can compare this with the results of Bramson \[B\] who showed (for positive solutions of the Ginzburg-Landau equation) that such a leakage is only possible if the initial data decay like $`e^xx^\alpha `$ with $`\alpha >1`$. In that case, he gets positive amplitudes at $`ct+(\alpha 1)\mathrm{log}t`$. Note that the condition $`J_0(0)<\mathrm{}`$ can only hold for $`\alpha <\frac{1}{2}`$, and then the correction term will push the amplitude behind the position of $`ct`$. Thus, in the case of the Ginzburg-Landau equations the two results are consistent.
5. An example of a non-linear velocity bound Consider the semi-linear parabolic equation
$$_tu=P(_x)u+f(u),$$
$`(5.1)`$
where $`P`$ is a real polynomial, $`ReP(ik)`$ diverges to $`\mathrm{}`$ as $`|k|\mathrm{}`$ and $`Im(ik)`$ is a polynomial of lower order.The complex Ginzburg-Landau equation is somewhat more complicated because in that case $`P`$ is a $`2\times 2`$ matrix polynomial. But it is covered by our methods. We also assume that $`f`$ is a $`𝒞^2`$ function for which $`f(0)=0`$, and $`f^{}(0)=0`$. This implies that $`u=0`$ is an unstable fixed point of (5.1) . We also assume that
$$\underset{|u|\mathrm{}}{lim\; sup}\frac{f(u)}{u}<\mathrm{\hspace{0.17em}0}.$$
This assumption ensures global existence and regularity of the semiflow (see \[CE\]). (If $`\stackrel{}{u}`$ is vector valued we impose $`lim\; sup_\stackrel{}{u}\mathrm{}\stackrel{}{u}\stackrel{}{f}(u)/\stackrel{}{u}^2<0`$.)
Define
$$\sigma =\underset{u}{sup}\frac{f(u)}{u}.$$
This is a finite positive quantity from the above assumptions (if $`\stackrel{}{u}`$ is vector valued we define it as the sup of $`\stackrel{}{u}\stackrel{}{f}(u)/\stackrel{}{u}^2`$ .) Note that one can have $`\sigma >f^{}(0)`$, and if this happens Aronson and Weinberger \[AW\] showed that the minimal speed is bounded above by $`\sqrt{4\sigma }`$, when $`P(ik)=k^2`$. In this section we show that the same result can be recovered for this, and many other equations using the methods of Section 4, again without any recourse to the maximum principle.
In this case, Eq.(4.7) becomes
$$G(\beta )=Q(\beta )+\sigma ,$$
where $`Q`$ is given by
$$Q(\beta )=\underset{k_\beta ^{}}{sup}ReP(\beta +ik_\beta ^{})$$
where the $`k_\beta ^{}`$ are the solutions of
$$\frac{dReP(\beta +ik)}{dk}|_{k=k_\beta ^{}}=\mathrm{\hspace{0.17em}0}.$$
The remainder of the proof is the same, except that in (4.5) the term $`\mathrm{exp}(2\beta (x\xi ))v^4`$ is replaced by
$$e^{\beta (x\xi )}vf\left(e^{\beta (x\xi )}v\right)\sigma v^2$$
After this modification the proof proceeds as before.
Appendix: The determination of the critical speed Let $`P`$ be a real polynomial for which $`ReP(ik)`$ diverges to $`\mathrm{}`$ as $`|k|\mathrm{}`$ and $`ImP(ik)`$ is of lower order. In the case of SH, we have $`P(z)=\epsilon ^2(1+z^2)^2`$. For $`\beta >0`$ we consider $`P(\beta +ik)`$, take the real part and look for an extremum in $`k`$. In other words, we solve
$$\frac{dReP(\beta +ik)}{dk}=\mathrm{\hspace{0.17em}0},$$
in the unknown $`k`$. Since $`P`$ is analytic, one can write this as
$$0=Im\left(\frac{dP(z)}{dz}|_{z=\beta +ik}\right).$$
$`(A.1)`$
For each $`\beta `$ we find solutions $`k_\beta ^{}`$. The velocity $`c_\beta ^{}`$ is related to the critical value of $`P`$ in (A.1) by
$$c_\beta ^{}=\underset{k_\beta ^{}}{sup}ReP(\beta +ik_\beta ^{})/\beta .$$
$`(A.2)`$
Then, the minimal speed is
$$c_{}=\underset{\beta (0,\mathrm{}]}{inf}c_\beta ^{},$$
which is determined by $`\mathrm{}`$ccc. To simplify the discussion, we will assume from now on that for all $`k_\beta ^{}`$ one obtains the same critical value. This is the case for the Ginzburg-Landau and Swift-Hohenberg equations.
Note that there is at least one $`\beta _{}`$ solving
$$_\beta \left(ReP(\beta +ik_\beta ^{})/\beta \right)|_{\beta =\beta _{}}=0,$$
$`(A.3)`$
for which $`c_{}=c_\beta _{}^{}`$.
In the approach of \[BBDKL\] the authors consider $`\omega _0(k)=P(ik)`$. They determine $`\overline{k}(c)𝐂`$ by
$$(d\omega _0/d\overline{k})|_{k=\overline{k}(c)}=ic,$$
$`(A.4)`$
and then $`c_{}𝐑`$ by the condition
$$Re\left(\omega (\overline{k}(c_{}))i\overline{k}(c_{})c_{}\right)=\mathrm{\hspace{0.17em}0}.$$
$`(A.5)`$
To compare the two approaches, note that $`P(\beta +ik)=\omega _0(k+i\beta )`$. Clearly the equations (A.2) and (A.5) are equivalent. To see that (A.1) and (A.4) say the same thing, note that since $`c`$ is real one has
$$Im(P^{}(z))=Re(\omega _0^{}(iz))=Re(\omega _0^{}(iz)+ic).$$
$`(A.6)`$
In particular, if $`\omega _0`$ is an even function, the relation $`Re(\omega _0^{}(\overline{k})ic)=0`$ is equivalent to requiring $`\omega _0^{}(\overline{k})=ic`$, which is (A.4). Using (A.6), we conclude that the solution $`\overline{k}`$ of $`\omega _0^{}(\overline{k})=ic`$ of \[BBDKL\] is the same as $`i`$ times the solution $`z`$ of $`Im(P^{}(z))=0`$, which is (A.1). Therefore $`\overline{k}=k_\beta _{}^{}+i\beta _{}`$. Finally, to find $`c_\beta _{}^{}`$ one can solve
$$\begin{array}{c}\text{ }0=Re(P(\beta _{}+ik_\beta _{}^{})\beta _{}c_\beta _{}^{})=Re(\omega _0(k_\beta _{}^{}+i\beta _{})ic_\beta _{}^{}(k_\beta _{}^{}+i\beta _{}))\text{ }\hfill \\ \text{ }=Re(\omega _0(\overline{k})ic_\beta _{}^{}\overline{k}).\text{ }\hfill \end{array}$$
Remark. The same kind of calculation can be done for multi-component problems (such as reaction diffusion), where $`P`$ would be a matrix. The example of the SH equation. In this case
$$P(z)=\epsilon ^2(1+z^2)^2,$$
and so
$$\omega _0(k)=P(ik)=\epsilon ^2+(1k^2)^2.$$
In \[BBDKL\], it is found that
$$\begin{array}{cc}\hfill \overline{k}& =\overline{k}_1+i\overline{k}_2,\hfill \\ \hfill \overline{k}_2& =\frac{\sqrt{\sqrt{1+6\epsilon ^2}1}}{12}=\frac{\epsilon ^2}{4}+\mathrm{},\hfill \\ \hfill \overline{k}_1& =\mathrm{\hspace{0.17em}1}+3\overline{k}_2^2,\hfill \\ \hfill c_{}& =\mathrm{\hspace{0.17em}8}\overline{k}_2(1+4\overline{k}_2^2)=\mathrm{\hspace{0.17em}4}\epsilon +\mathrm{}.\hfill \end{array}$$
$`(A.7)`$
In our formulation, we find
$$P(\beta +ik)=\epsilon ^2\left(1+(ik\beta )^2\right)^2.$$
The real part of the derivative w.r.t. $`k`$ yields
$$\frac{dReP(\beta +ik)}{dk}=\mathrm{\hspace{0.17em}4}k4k^3+12k\beta ^2.$$
The solutions of $`dReP(\beta +ik)/dk=0`$ are $`k_\beta ^{}=\pm \sqrt{1+3\beta ^2}`$ (and $`k_\beta ^{}=0`$ which leads to less stringent bounds). Substituting back into $`ReP`$, we get
$$ReP(\beta +ik_\beta ^{})=\epsilon ^2+4\beta ^2+8\beta ^4,$$
which is what we announced in (1.6) and got as a result of integration by parts in Eqs.(4.5)--(4.7). Solving now
$$ReP(\beta +ik_\beta ^{})c\beta =\mathrm{\hspace{0.17em}0},$$
for $`c=c_\beta ^{}`$ leads to $`c_\beta ^{}=(\epsilon ^2+4\beta ^2+8\beta ^4)/\beta `$. To find the absolutely minimal speed, we find that $`\beta `$ for which $`c_\beta ^{}`$ is extremal, that is $`_\beta c_\beta ^{}=0`$. The only positive solution is
$$\beta _{}=\frac{\sqrt{3\sqrt{1+6\epsilon ^2}3}}{6},$$
and hence,
$$c_\beta _{}^{}=\frac{4\left(\sqrt{1+6\epsilon ^2}1+6\epsilon ^2\right)}{3\sqrt{3\sqrt{1+6\epsilon ^2}3}}.$$
This quantity is the same as $`inf_{\beta 𝐑}c_\beta `$ where $`c_\beta `$ is given by (1.7). Acknowledgments. We thank W. van Saarloos for some help with the references, and M. Hairer and G. van Baalen for a critical reading of the manuscript. This work was partially supported by the Fonds National Suisse.
References
\[AW\]Aronson, D. and H.F. Weinberger: Multidimensional nonlinear diffusion arising in population genetics. Adv. Math. 30, 33--76 (1978).
\[BBDKL\]Ben-Jacob, E., H. Brand, G. Dee, L. Kramers, and J.S. Langer: Pattern propagation in non-linear dissipative systems. Physica 14D, 348--364 (1985).
\[B\]Bramson, M.: Convergence of solutions of the Kolmogorov equation to travelling waves.. Mem. Amer. Math. Soc. 44, No. 285, iv+190pp (1985).
\[CE\]Collet, P. and J.-P. Eckmann: Instabilities and Fronts in Extended Systems, Princeton University Press (1990).
\[DL\]Dee, G. and J. Langer: Propagating pattern selection. Phys. Rev. Lett. 50, 383--386 (1983).
\[LMK\]Langer, J. and H. Müller-Krumbhaar: Mode-selection in a dendrite-like nonlinear system. Phys. Rev. A 27, 499--514 (1983).
\[vS\]van Saarloos, W.: Front propagation into unstable states: Marginal stability as a dynamical mechanism for velocity selection. Phys. Rev. A 37, 211--229 (1988).
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# Path Dependent Option Pricing: the path integral partial averaging method
## 1 Introduction
Financial derivatives (eg. options and futures) derive their value from an underlying traded financial security, whose price is modeled by some stochastic process. In their most general form, the option payoff is path dependent since it depends on the entire future path traversed by the underlying security. Path dependent options are defined using either discrete or continuous price sampling. For continuous sampling closed from solutions are often available, but in practice most traded path dependent options are discretely sampled. It is known that the application of these closed form solutions leads to substantial pricing errors for discretely sampled options . This feature has necessitated the development of practical and efficient computational methods for the evaluation of path dependent options . Most research has focussed on either partial differential equation, Monte Carlo or tree based methods. In contrast to these approaches, in this paper I will develop an alternative based on the path integral formulation of the pricing problem.
For many years, theoretical physicists have been developing and applying the path integral method for calculating expectations similar to those now being encountered in the evaluation of financial derivatives (via the risk-neutral valuation formula). A path integral is an infinite dimensional Riemannian integral as the integral is performed over a set of functions or paths. The path integral method is unique in that it gives a global formulation of the problem in question. This global description provides a powerful tool for deriving analytical approximation and numerical solution schemes that are difficult or impossible to formulate in other ways. Formally, the path integral method is easily extended to multi-dimensional problems. Much of the driving force behind the development of path integral numerical methods in physics has been that, compared to other methods, they show only a slow increase in computational complexity as the dimensionality of the problem increases. Path integral methods were first introduced by Feynman in 1942 as an alternative formulation of quantum physics . They have found wide application in the evaluation of both the real time dynamics and equilibrium statistical mechanics of quantum many-body systems . They are also a popular and natural tool for the analysis of diffusion processes , including non-Markovian systems where there is a lack of practical alternatives .
The application of path integral methods to financial derivatives was pioneered by Dash who developed a path integral framework for pricing bonds and options within one factor term structure models . More recently, Linetsky was the first to show how path dependent options could be formally priced in a path integral framework. He considered several examples of one and two factor path dependent options. Baaquie has shown how path integral methods can be used to price vanilla options with stochastic volatility. The same author has also recast the popular Heath-Jarrow-Morton model of forward interest rates as a problem in path integration . Similar to Dash, Otto has more recently shown how to use path integration to price bonds and bond options under general short rate models . Bennati et-al have focussed on a multi-dimensional path integral formalism for solving general financial problems based on systems of stochastic equations. All these preceding authors focussed on general formalism and exactly solvable models. They have shown that path integrals constitute a natural framework for describing the evaluation of general multi-factor derivatives. Although path integral methods offer an attractive way of obtaining exact solutions, they cannot find exact solutions not available using more standard methods. The greatest promise of path integral methods will be in the development of new numerical and approximation methods for addressing pricing problems where exact solutions are impossible.
Path integral numerical methods involve the evaluation of a multi-dimensional integral, either by deterministic or Monte Carlo methods. A review of deterministic and Monte Carlo methods, developed for physics applications, can be found in Drozdov and Makri respectively. An independent and much more recent development has been the application of path integral numerical methods to financial derivatives. This was pioneered by Eydeland who, using fast Fourier transform methods, has derived a deterministic path integral algorithm for calculating the generating function of a random variable defined as the time integral of a general diffusion process. He pointed to a number of potential financial applications of this algorithm. Chiarella and El-Hassan have applied the method of Eydeland to the pricing of American bond options in the Heath-Jarrow-Morton framework . More recently, Chiarella et-al have devised a deterministic path integral algorithm based on Fourier-Hermite series expansions and applied it to the pricing of American options and point barrier options . They report computational times that are a significant improvement over the standard binomial and finite difference methods. In contrast to these deterministic methods, Makivic has initiated research into the path integral Monte Carlo evaluation of financial derivatives. He broadly outlines a computational approach based on the standard Metropolis algorithm. Rosa-clot and Taddei have discussed both a deterministic method and the Monte Carlo approach and applied them to several examples. Some of these previous papers point out that the path integral method is superior to the the traditionally used lattice methods, because the underlying asset price is left continuous rather than being discretized. This has several important advantages. The option price is obtained more accurately since all possible price paths are included in the simulation. Option Greeks are obtained more reliably since the method avoids the need for numerical differentiation. A further advantage of path integral methods is that they can be easily and efficiently extended to evaluate multi-factor financial derivatives.
The application of path integral methods here is fundamentally different from all the previous cited works. We use the path integral representation of the option price to show that rather generally, it is possible to perform analytically a partial averaging over the underlying risk-neutral diffusion process. This key result will greatly reduce the computational burden placed on the subsequent numerical evaluation. The application of this method is inspired by a technique first developed for computation problems in chemical physics . Conceptually, the partial averaging corresponds to averaging over the high frequency fluctuations of the risk-neutral diffusion process. For short-medium term options it leads to a general approximation formula that only requires the evaluation of a one dimensional numerical integral. Longer term options can be evaluated deterministically or most generally by standard Monte Carlo methods. In this case, the partial averaging method will greatly reduce the required dimension of the Monte Carlo simulation.
The outline of this paper is as follows. In section 2 we develop a formal path integral representation for the evaluation of rather general path dependent options. In section 3 we show how the previous path integral can be numerically evaluated after performing analytically a partial averaging over the underlying risk-neutral diffusion process. In section 4 some examples are presented. For clarity of presentation and completeness, the path integral methods used in this paper are developed from first principles in the three appendices.
## 2 Path dependent option theory
In this section we will present a formal path integral framework general enough to price a wide range of path dependent options. We will assume the standard geometric Brownian motion (gbm) model of the asset price process. More general diffusion processes present no special difficulties. In this case the risk-neutral price process is given by the Ito stochastic differential equation
$$dS_t=(rq)S_tdt+\sigma S_tdW_t,$$
(2.1)
where $`r`$ is the interest rate and $`q`$ is the continuous dividend yield. Using Ito’s lemma we can show
$$dx_t=\mu dt+\sigma dW_t,$$
(2.2)
where
$$x_t=\mathrm{ln}S_t,\mu =rq\frac{1}{2}\sigma ^2.$$
(2.3)
Consider a path dependent option with price $`C_F`$ at expiry $`u`$ given by
$$C_F(S_u,,u)=F(S_u,)=F(e^{x_u},),$$
(2.4)
where $`F`$ is the option payoff function which depends on some path dependent random variable $``$. In this paper we will assume that $``$ can be written as
$$=_t^u𝑑sw(s)f(x_s,s),$$
(2.5)
which is a time integral over an arbitrary function of the risk-neutral diffusion process (2.2). For continuous sampling $`w(s)=1`$, but for discrete sampling (which is more realistic in practice)
$$w(s)=\underset{i}{}w_i\delta (ss_i),$$
(2.6)
where $`w_i`$ are the sampling weights and $`s_i`$ are the sampling times. The above definition of $``$ was used before by Wilmott et-al whose focus was on partial differential equation methods. It was shown to be general enough to include Asian, barrier and lookback options which are 3 qualitatively different path dependent options. We will present examples in section 4.
In a risk-neutral framework, the option price at inception time $`t`$ is given by
$$C_F(S_t,t)=e^{rT}E_{x_t}[F(e^{x_u},)],$$
(2.7)
where $`T=ut`$ and the expectation is with respect to the transformed risk-neutral price process (2.2) conditioned on the initial value $`x_t`$. The price at inception can then be expressed as
$$C_F(S_t,t)=e^{rT}_{\mathrm{}}^{\mathrm{}}𝑑x_u_{\mathrm{}}^{\mathrm{}}𝑑P(x_u,|x_t)F(e^{x_u},),$$
(2.8)
where $`P(x_u,|x_t)`$ is the joint probability density function (pdf) of $`x_u`$ and the path dependent random variable $``$. In appendix A we show for a general diffusion process how the joint pdf can be formally computed as a path integral. For the special case of gbm, we show in appendix A that the joint pdf is given by
$$P(x_u,|x_t)=\frac{1}{2\pi }\mathrm{exp}\left[\frac{\mu x_u}{\sigma ^2}\frac{\mu x_t}{\sigma ^2}\frac{\mu ^2T}{2\sigma ^2}\right]_{\mathrm{}}^{\mathrm{}}𝑑ke^{ik}K(x_u,x_t;T),$$
(2.9)
where $`\mu `$ is the constant defined in (2.3) and $`K`$, which we refer to as the propagator is defined by
$$K(x_u,x_t;T)=_{x_t}^{x_u}𝒟x_s\mathrm{exp}\left[\frac{1}{2\sigma ^2}_t^u𝑑s\left(\dot{x}_s^2+V(x_s,s)\right)\right].$$
(2.10)
In (2.10), the integration measure $`𝒟x_s`$ denotes a path integral which is defined precisely in appendix A. It describes an infinite dimensional integral over all paths connecting $`x_u`$ at the expiry time and $`x_t`$ at the initial time. We refer to the function $`V`$ in (2.10) as the potential function. For the gbm model it has the form
$$V(x_s,s)=2ik\sigma ^2w(s)f(x_s,s).$$
(2.11)
We see that in this case the potential is imaginary with a functional form determined by the path dependent random variable (2.5). We show in appendix A that for more general risk-neutral diffusion processes the potential function is complex.
The propagator (2.10), up to a boundary term, is the characteristic function of the joint pdf. In physics, the path integral (2.10) is equivalent to the path integral representation for the equilibrium quantum statistical density matrix of a particle in a complex potential $`V(x)`$ . In this case temperature replaces the role of time in (2.10). Under an imaginary time transformation, (2.10) becomes equivalent to that which gives the quantum mechanical propagator of a one dimensional quantum particle in a complex potential $`V(x)`$. These identifications with standard problems in theoretical physics are of great value because we can then use the methods and results of theoretical physics for performing these path integrals. Tables of known exact path integrals of the form (2.10), for various potential functions, are listed by Grosche . Such tables along with the formulation provided here provide an easy way to obtain exact solutions to path dependent option prices.
### 2.1 Seasoned path dependent options
In this paper we consider the option at its inception time. More generally, the option price at time $`t^{}`$ ($`t<t^{}<u`$) for a seasoned path dependent option is given by
$$C_F(S_t^{},_t^t^{},t^{})=e^{r(ut^{})}E_{x_t^{}}[F(e^{x_u},_t^t^{}+_t^{}^u)],$$
(2.12)
where we use the more detailed notation
$$_t^u=_t^u𝑑sw(s)f(x_s,s).$$
(2.13)
We then find
$$C_F(S_t^{},_t^t^{},t^{})=e^{r(ut^{})}_{\mathrm{}}^{\mathrm{}}𝑑x_u_{\mathrm{}}^{\mathrm{}}𝑑_t^{}^uP(x_u,_t^{}^u|x_t^{})F(e^{x_u},_t^t^{}+_t^{}^u).$$
(2.14)
We see that this case only affects the payoff function in a simple way and the joint pdf we need to find is the same problem as before. Therefore all the final results can be trivially extended to seasoned options.
## 3 Partial Averaging
In the previous section we showed that for the gbm model, equations (2.5) and (2.9-11) define a formal path integral representation for the joint pdf. The option price is then obtained from this joint pdf via (2.8). In this section we use the previous path integral formulation to show that it is possible to perform analytically a partial averaging over the underlying risk-neutral diffusion process. The option price can then be more efficiently evaluated by numerical methods.
First we must discretize in time (2.2), by defining the discrete time $`s_n=n\epsilon +t`$, where $`n=0,1,\mathrm{},N`$ and $`\epsilon =T/N`$ with $`T=ut`$. The propagator (2.10) can then be decomposed as
$$K(x_u,x_t;T)=_{\mathrm{}}^{\mathrm{}}𝑑x_{N1}\mathrm{}𝑑x_1\underset{n=1}{\overset{N}{}}K(x_n,x_{n1};\epsilon )$$
(3.1)
where $`x_Nx_u,x_0x_t`$ and a general form for the short-time propagator is, as shown in appendix C,
$$K(x_n,x_{n1};\epsilon )=\left(\frac{1}{2\pi \sigma ^2\epsilon }\right)^{1/2}\mathrm{exp}\left[\frac{(x_nx_{n1})^2}{2\sigma ^2\epsilon }\gamma (x_n,x_{n1};\epsilon )\right],$$
(3.2)
where $`\gamma `$ will be defined below. Substituting (3.2) into (3.1) we find that the propagator becomes
$$K(x_u,x_t;T)=\left(\frac{1}{2\pi \sigma ^2\epsilon }\right)^{\frac{N}{2}}_{\mathrm{}}^{\mathrm{}}𝑑x_{N1}\mathrm{}𝑑x_1\mathrm{exp}\left[\underset{n=1}{\overset{N}{}}\frac{(x_nx_{n1})^2}{2\sigma ^2\epsilon }\underset{n=1}{\overset{N}{}}\gamma (x_n,x_{n1};\epsilon )\right].$$
(3.3)
As $`\epsilon 0`$, its possible to show that
$$\gamma (x_n,x_{n1};\epsilon )=\epsilon \frac{(V(x_n,s_n)+V(x_{n1},s_{n1}))}{4\sigma ^2}$$
(3.4)
yields the correct short-time propagator when substituted into (3.2).
We will refer to equation (3.2), with (3.4), as the primitive short-time propagator. It is the standard short-time propagator used in the numerical evaluation of path integrals. A key feature of this propagator is that it is not correct to first order in $`\epsilon `$. It is in fact only valid as $`\epsilon 0`$. Clearly, the dimension of the integral in (3.1) could be made much smaller by searching for short-time propagators accurate over larger time-steps. This observation has motivated the search for improved short-time propagators for use in path integral calculations for physics applications . In appendix C, we show that its possible in general to write
$$\gamma (x_n,x_{n1};\epsilon )=\underset{m=1}{\overset{\mathrm{}}{}}\frac{1}{m!}\left(\frac{\epsilon }{2\sigma ^2}\right)^mC_m(x_n,x_{n1};\epsilon ),$$
(3.5)
where $`C_m`$ describes a cumulant structure with each cumulant of order $`(\sigma ^2\epsilon )^{m1}`$. The key result is, we can obtain a short-time propagator formally correct to second order in $`\epsilon `$, by truncating all cumulants beyond the first. This truncation corresponds to retaining an averaging over only the high frequency fluctuations of the risk-neutral diffusion process; i.e. a partial averaging. The improved short-time propagator will lead to a much more efficient numerical evaluation compared to that obtained by using the primitive short-time propagator. In appendix C, we calculate the first 2 cumulants exactly for a general potential and derive an expansion of the propagator to third order in $`T`$ (a result only valid for smooth potentials).
Using the results from appendix C, with the gbm potential (2.11), we find that
$$\gamma (x_n,x_{n1};\epsilon )i\epsilon k\alpha (x_n,x_{n1};\epsilon )+\frac{\epsilon ^2k^2}{2}\beta (x_n,x_{n1};\epsilon )+o(\sigma ^4\epsilon ^5),$$
(3.6)
where
$$\alpha (x_n,x_{n1};\epsilon )=_0^1𝑑\tau w(\tau )_{\mathrm{}}^{\mathrm{}}𝑑p_\tau P(p_\tau )f(\overline{x}_\tau +p_\tau ,\tau ),$$
(3.7)
$$P(p_\tau )=\frac{1}{\sqrt{2\pi \nu _\tau ^2}}\mathrm{exp}(p_\tau ^2/2\nu _\tau ^2)$$
(3.8)
and
$$\nu _\tau ^2=\sigma ^2\epsilon (1\tau )\tau ,\overline{x}_\tau =\tau (x_nx_{n1})+x_{n1},\tau =(ss_{n1})/\epsilon .$$
(3.9)
Equation (3.7) is a key equation as it contains the partial averaging. The origin of $`\alpha `$ is the first cumulant in the expansion (3.5) and its evaluation will lead to a short-time propagator correct to second order in $`\epsilon `$. Expanding $`\alpha `$ to order $`\epsilon `$ is consistent with the truncation of the higher cumulants in (3.5). The origin of $`\beta `$ is the second cumulant which is formally calculated in appendix C. All we need to know here is that $`\beta `$ is of order $`\sigma ^2\epsilon `$, since it will be set to zero at the end. We keep it to determine the order of the first correction term due to the truncation of all higher cumulants. After combining (3.6),(3.3) and (2.9) and performing the integration over $`k`$, we find that the joint pdf becomes
$`P(x_u,|x_t)`$ $``$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}dx_{N1}\mathrm{}dx_1P[x_N,..,x_1|x_0]`$ (3.10)
$`\times `$ $`\mathrm{exp}\left[{\displaystyle \frac{\left(\epsilon _{n=1}^N\alpha (x_n,x_{n1};\epsilon )\right)^2}{2\epsilon ^2_{n=1}^N\beta (x_n,x_{n1};\epsilon )}}\right]\left(2\pi \epsilon ^2{\displaystyle \underset{n=1}{\overset{N}{}}}\beta (x_n,x_{n1};\epsilon )\right)^{1/2}.`$
where $`x_Nx_u`$, $`x_tx_0`$ and the discrete path pdf is given by
$$P[x_N,\mathrm{},x_1|x_0]=\left(\frac{1}{2\pi \sigma ^2\epsilon }\right)^{\frac{N}{2}}\mathrm{exp}\left[\underset{n=1}{\overset{N}{}}\frac{(x_nx_{n1}\mu \epsilon )^2}{2\sigma ^2\epsilon }\right].$$
(3.11)
Equation (3.11) describes the probability density of realizing a particular discrete path of the risk-neutral stochastic process (2.2). In the limit that $`\beta `$ tends to zero, the joint pdf (3.10) becomes a delta function in $``$ and we can show, using (2.8), that the option price becomes
$$C_F(S_t,t)e^{rT}_{\mathrm{}}^{\mathrm{}}𝑑x_N\mathrm{}𝑑x_1P[x_N,\mathrm{},x_1|x_0]F(e^{x_N},\epsilon \underset{n=1}{\overset{N}{}}\alpha (x_n,x_{n1};\epsilon ))+o(\epsilon ^2\sigma ^2T).$$
(3.12)
The order of the correction term follows from (3.10) and a saddle point expansion of the Gaussian integral over $``$ in (2.8).
Equation (3.12) is the major result of this paper. The multi-dimensional integral can be evaluated by deterministic methods or more generally by standard Monte Carlo methods. In this case we use the observation that (3.12) is equivalent to an expectation of the option payoff function $`F`$, with respect to the discretely sampled risk-neutral diffusion process defined by (3.11), or equivalently by (2.2). In (3.12), the discretization time-scale $`\epsilon `$ is independent of any discrete option sampling time-scale. If we choose $`\epsilon `$ to match the interval between discrete option sampling, we find that (3.7) will reduce to a primitive short-time propagator and (3.12) will become
$$C_F(S_t,t)=e^{rT}_{\mathrm{}}^{\mathrm{}}𝑑x_N\mathrm{}𝑑x_1P[x_N,\mathrm{},x_1|x_0]F(e^{x_N},\underset{n=0}{\overset{N}{}}w_nf(x_n)).$$
(3.13)
This is equivalent to a direct discretization of (2.7), consistent with a discretely sampled path dependent random variable described by (2.5) and (2.6). In this case no analytical partial averaging has been performed and the Monte Carlo evaluation of (3.13) is completely standard. The great advantage of the partial averaging method is that in (3.12), the discrete time interval $`\epsilon `$ can be chosen to be much larger than the option sampling time-scale. The partial averaging is performed in (3.7) where we must average over Gaussian fluctuations about the straight line path $`\overline{x}_\tau `$ connecting $`x_n`$ and $`x_{n1}`$. This corresponds to averaging over the high frequency fluctuations of the risk-neutral diffusion process. In practice, as will be seen in the next section, the partial averaging integral is easily performed analytically. It is simply the Gaussian transform of the function $`f`$, which defines the path dependent random variable in question via (2.5). Of most practical interest will be discretely sampled path dependent options, where $`w(s)`$ is given by (2.6). In this case the subsequent integral over $`\tau `$ in (3.7) reduces to a discrete sum which presents no problems.
For the special case $`N=1`$, for which $`\epsilon =T`$, $`x_1x_u`$ and $`x_0x_t`$, we find that (3.12) becomes
$$C_F(S_t,t)\frac{e^{rT}}{\sqrt{2\pi \sigma ^2T}}_{\mathrm{}}^{\mathrm{}}𝑑x_u\mathrm{exp}\left[\frac{(x_ux_t\mu T)^2}{2\sigma ^2T}\right]F(e^{x_u},T\alpha (x_u,x_t;T))+o(\sigma ^2T^3).$$
(3.14)
This describes rather generally an approximate path dependent option price. It can be simply evaluated as a one dimensional numerical integral. As a measure of the accuracy of (3.14), we ask at what time to maturity $`T`$, does the error term in (3.14) become $`1\%`$ of the true option price. We assume a typical market volatility of $`\sigma =0.25`$ and that the unknown coefficient of the correction term is equal to the true option price. We find that $`T0.5`$ gives a $`1\%`$ error, while $`T0.25`$ gives an error of approximately $`0.1\%`$. These estimates begin to illustrate the power of the method presented here.
## 4 Examples
The previous formulation is rather general and can be applied to a range of path dependent options. In this section we will show how two important and qualitatively different classes of path dependent options fit into this framework.
### 4.1 Average rate options
The payoff of the geometric Asian option is some function of the path dependent random variable given by
$$=_t^u𝑑sw(s)x_s,$$
(4.1)
where $`x_s`$ is related to the risk-neutral asset price by (2.3). From (2.5) we can identify $`f`$ with $`x_s`$. After performing the partial averaging (3.7), we find
$$\alpha (x_n,x_{n1};\epsilon )=_0^1𝑑\tau w(\tau )\overline{x}_\tau .$$
(4.2)
For the continuous sampling ($`w(\tau )=1`$) we find
$$\alpha (x_n,x_{n1};\epsilon )=(x_n+x_{n1})/2.$$
(4.3)
For this simple example $`\alpha `$ is just the primitive short-time propagator.
The payoff of the arithmetic Asian option will be some function of the path dependent random variable
$$=_t^u𝑑sw(s)e^{x_s}.$$
(4.4)
From (2.5) we identify $`f`$ with $`e^{x_s}`$. After performing the partial averaging (3.7) we find
$$\alpha (x_n,x_{n1};\epsilon )=_0^1𝑑\tau w(\tau )e^{\overline{x}_\tau +\nu _\tau ^2/2}.$$
(4.5)
We can expand (4.5) to order $`\sigma ^2\epsilon `$ without a significant loss of accuracy. This is consistent with the order of the truncation of the cumulant expansion (3.5). We then find
$$\alpha (x_n,x_{n1};\epsilon )_0^1𝑑\tau w(\tau )e^{\overline{x}_\tau }\left(1+\sigma ^2\epsilon (1\tau )\tau /2+o(\sigma ^4\epsilon ^2)\right).$$
(4.6)
The final result will depend on whether we use discrete or continuous sampling. For continuous sampling we have $`w(\tau )=1`$ and (4.6) becomes
$$\alpha (x_n,x_{n1};\epsilon )e^{x_{n1}}\left[\frac{1}{a}\left(e^a1\right)+\frac{\sigma ^2\epsilon }{2a^3}\left(e^a(a2)+a+2\right)+o(\sigma ^4\epsilon ^2)\right],$$
(4.7)
where $`a=x_nx_{n1}`$. For discrete sampling we can perform the necessary summations analytically so the method is equally effective. It is instructive to compare (4.7) with the $`\alpha `$ that generates the primitive short-time propagator. This is obtained by approximating (3.7) by
$$\alpha (x_n,x_{n1};\epsilon )\frac{1}{2}\left(f(x_n)+f(x_{n1})\right).$$
(4.8)
One can see that (4.7) will only reduce to (4.8) in the limit $`a0`$ (when $`f=e^{x_s}`$).
We can now use (3.12) to perform a Monte Carlo evaluation of the continuously sampled Asian option. Using (4.7) allows us to obtain accurate results by simulating only relatively low dimensional random paths of (2.2). This will avoid the problems associated with the Monte Carlo evaluation of these options .
### 4.2 Occupation time derivatives
A large and important class of path dependent options are those where the payoff depends on the time the asset price spends within a given region. This class of options have been referred to as occupation time derivatives and they have been discussed in detail by Linetsky and Hugonnier . The valuation of debt and contingent claims with default risk can also be placed in this class .
As a simple example, consider an option with a payoff that depends on the time spent below a possibly time dependent barrier $`B_s`$. This time is a path dependent random variable which can be expressed as (and similarly for the above case)
$$=_t^u𝑑sw(s)H(B_se^{x_s}),$$
(4.9)
where $`H`$ is a simple step function. We can therefore identify $`H`$ with the function $`f`$ in (2.5). After performing the partial averaging (3.7) we find
$$\alpha (x_n,x_{n1};\epsilon )=_0^1𝑑\tau w(\tau )N\left((\mathrm{ln}B_s\overline{x}_\tau )/\nu _\tau \right),$$
(4.10)
where $`N`$ is the cumulative normal distribution function. This function is difficult to integrate analytically. However the real case of practical interest is discrete sampling. In this case the integration reduces to a discrete summation and presents no problems. Barrier options are the most basic and well known example of options that fall under this framework. An interesting generalization is the step option which has a finite knock-out rate. This is motivated by risk-management arguments and a variety of possible payoff functions have been discussed . All variations are included in this framework since the payoff function is an arbitrary function of the path dependent random variable.
Consider the example of derivatives depending on the occupation time between two barriers $`B_1`$ and $`B_2`$.<sup>1</sup><sup>1</sup>1the occupation time outside these barriers can be trivially constructed from this time Included in this class are double barrier options and the range products such as range notes and corridor options . This time is a path dependent random variable which can be expressed as
$$=_t^u𝑑sw(s)\left[H(B_2e^{x_s})H(B_1e^{x_s})\right].$$
(4.11)
From (2.5) we can identify the function $`f`$ with the difference of two step functions. After performing the partial averaging (3.7) we find
$$\alpha (x_n,x_{n1};\epsilon )=_0^1𝑑\tau w(\tau )\left[N\left((\mathrm{ln}B_2\overline{x}_\tau )/\nu _\tau \right)N\left((\mathrm{ln}B_1\overline{x}_\tau )/\nu _\tau \right)\right].$$
(4.12)
Our framework is easily generalized to include derivatives dependent on several occupation times.
## 5 Discussion and Conclusion
The aim of this paper was to present a new approach to evaluating the price of path dependent options. We considered options with the general payoff function (2.4), contingent on a path dependent random variable expressible in the form (2.5). The key results of this paper were the general evaluation formula (3.12) and the short time to expiry approximation (3.14). They give the option price after performing analytically a partial averaging, defined by equations (3.7-9), over the underlying risk-neutral diffusion process (2.2). Since the method is analytically based, it also gives the order of the error made by choosing a finite time discretization. Specific examples were presented in section 4.
In (3.12), the partial averaging allows one to choose the discrete time interval $`\epsilon `$ to be much larger than the option sampling time-scale. We could, for example, imagine choosing $`\epsilon `$ to be one month for an option with daily sampling. Clearly, the partial averaging method can greatly reduce the dimension of the integral in (3.12). This integral can be evaluated most generally by standard Monte Carlo methods. In this case the partial averaging will greatly reduce the dimension of the random paths to be simulated. In effect, it allows random simulations to be replaced with deterministic calculations. Standard methods to increase the efficiency of the Monte Carlo evaluation can still be used. These include variance reduction techniques , the simulation of sample paths using the Brownian bridge process or the use of quasi Monte Carlo sampling . Interestingly, quasi Monte Carlo sampling is known to be more advantageous for low dimensional numerical integrals. The partial averaging method can therefore increase the relative gains made by the implementation of quasi Monte Carlo methods.
The framework presented here can be extended to multi-factor path dependent options and to path dependent options dependent on several path dependent random variables. Recent work has shown that a path-independent option price, in a stochastic volatility environment, can be approximated by pricing a more complex path dependent option in the usual Black-Scholes framework. The volatility risk is included by an extra continuous path-dependent payout function which has the same form as (2.5) . This work suggests that the computational framework presented here might also be used to price path dependent options consistent with the market implied volatility. Further work is required in this direction.
Path integral methods have long been developed and used as a computational tool in theoretical and chemical physics. Hopefully the work presented here will stimulate more interest in the application of these methods to problems in computational finance.
## Appendix A Path Integral Representation
Consider a stochastic process $`x_s`$, which obeys the stochastic differential equation <sup>2</sup><sup>2</sup>2 Note that any one-dimensional risk-neutral diffusion process can be cast into this form by a change of variable
$$dx_t=g^{}(x_t)dt+\sigma dW_t.$$
(A.1)
In our notation $`x_sx(s)`$. We will rewrite this equation as
$$\dot{x}_t=g^{}(x_t)+\sigma \zeta _t$$
(A.2)
where $`\zeta _t=\frac{dW_t}{dt}`$ and a dot denotes the derivative with respect to time. This notation is more suited for the path integral formulation. In (A.1) we assume that $`\sigma `$ in independent of $`x_t`$ (additive noise). The path integral formulation of the more general case (multiplicative noise) can be found in .
Consider a general continuous Gaussian noise process $`\chi (s)`$ (note that small $`s`$ denotes the time history variable between time $`t`$ and $`u`$ and should not be confused with the price $`S_t`$). We discretize this process by defining a discrete time $`s_n=n\epsilon +t`$, where $`n=0,1,\mathrm{},N`$ and $`\epsilon =T/N`$ with $`T=ut`$. Note that $`\epsilon `$ is equivalent to $`ds`$ when $`N\mathrm{}`$. The now discrete Gaussian process is fully defined by the normalized probability density functional
$$𝒫[\chi _1,\mathrm{},\chi _N]=\left(2\pi \right)^{n/2}\left(detR\right)^{1/2}\mathrm{exp}\left[\frac{1}{2}\underset{n,m=1}{\overset{N}{}}\chi _nR_{nm}^1\chi _m\right],$$
(A.3)
where $`\chi _n=\chi (s_n)`$, $`E[\chi _n\chi _m]=R_{nm}`$ and $`R_{nm}^1`$ denotes the inverse matrix. Equation (A.3) fully defines the probability density of the whole history of the discrete Gaussian process. This normalization condition implies
$$_{\mathrm{}}^{\mathrm{}}𝒫[\chi _1,\mathrm{},\chi _N]𝑑\chi _1\mathrm{}𝑑\chi _N=1.$$
(A.4)
Since $`dW_t`$ has a variance $`dt`$, we know $`\zeta _t`$ will have a variance $`dt^1`$. We then find that for the discrete time white noise process $`\zeta _n`$, we have $`E[\zeta _n\zeta _m]=\epsilon ^1\delta _{nm}`$ where $`\delta _{nm}`$ is the unit diagonal matrix. Using (A.3) we then find
$$𝒫[\zeta _1,\mathrm{},\zeta _N]=\left(\frac{\epsilon }{2\pi }\right)^{N/2}\mathrm{exp}\left[\frac{1}{2}\underset{n=1}{\overset{N}{}}\epsilon \zeta _n^2\right].$$
(A.5)
In the continuous limit this becomes
$$𝒫[\zeta _s]=\left(\frac{\epsilon }{2\pi }\right)^{N/2}\mathrm{exp}\left[\frac{1}{2}_t^u𝑑s\zeta _s^2\right],$$
(A.6)
where the continuous time expressions are written with the understanding that $`\epsilon 0`$ and $`N\mathrm{}`$ such that $`N\epsilon =T`$.
We wish to use the probability density functional (A.6) to find a probability density functional for the stochastic process $`x_s`$. To do this we need to first discretize in time (A.2). In discrete time (A.2) becomes
$$\frac{x_nx_{n1}}{\epsilon }=g^{}(\stackrel{~}{x}_n)+\sigma \zeta _n,$$
(A.7)
where
$$\stackrel{~}{x}_n=\varphi x_n+(1\varphi )x_{n1}$$
(A.8)
and $`\varphi `$ is a discretization parameter between 0 and 1. From (A.7) we see that a path $`\{\zeta _1,\mathrm{},\zeta _N\}`$ maps to a unique path $`\{x_1,\mathrm{},x_N\}`$ *as long as* $`x_0`$ *is given*. We will therefore write, by virtue of (A.4)
$$_{\mathrm{}}^{\mathrm{}}𝒫[x_1,\mathrm{},x_N|x_0]𝒥𝑑x_1\mathrm{}𝑑x_N=1,$$
(A.9)
where $`𝒥`$, defined by
$$𝒥=\mathrm{det}\left|\frac{\zeta _n}{x_m}\right|,m,n=1,\mathrm{},N$$
(A.10)
is the Jacobian of the change in coordinates. In the continuous limit we can substitute (A.2) into (A.6) to obtain
$$𝒫\left[x_s\right]=\left(\frac{\epsilon }{2\pi }\right)^{N/2}\mathrm{exp}\left[\frac{1}{2\sigma ^2}_t^u𝑑s\left(\dot{x}_s+g^{}(x_s)\right)^2\right].$$
(A.11)
We show in appendix B that the Jacobian (A.10) becomes in the continuous limit
$$𝒥=\left(\sigma \epsilon \right)^N\mathrm{exp}\left[\varphi _t^u𝑑sg^{\prime \prime }\left(x_s\right)\right].$$
(A.12)
The only sensible choice is $`\varphi =1/2`$ . Combining the Jacobian and (A.11) we obtain a new probability density functional
$$𝒫_x\left[x_s\right]=\mathrm{exp}\left[\frac{1}{2\sigma ^2}_t^u𝑑s\left(\dot{x}_s+g^{}(x_s)\right)^2+\frac{1}{2}_t^u𝑑sg^{\prime \prime }(x_s)\right],$$
(A.13)
which is normalized with respect to the functional measure $`𝒟_s`$ which is the continuous limit of
$$𝒟x_s=\left(2\pi \epsilon \sigma ^2\right)^{N/2}dx_1\mathrm{}dx_{N1}.$$
(A.14)
This means the conditional probability density function (pdf) of (A.2) is given by the path integral
$$P(x_u,x_t)=_{x_t}^{x_u}𝒟x_s𝒫_x[x_s].$$
(A.15)
The expectation of a functional $`[x_s]`$, conditional on the initial value $`x_t`$, can now be expressed in the path integral form
$$E_{x_t}\left[[x_s]\right]=_{\mathrm{}}^{\mathrm{}}𝑑x_u_{x_t}^{x_u}𝒟x_s𝒫_x[x_s][x_s].$$
(A.16)
Using (2.7), we can now write the option price in the path integral form
$$C_F(S_t,t)=e^{rT}_{\mathrm{}}^{\mathrm{}}𝑑x_u_{x_t}^{x_u}𝒟x_s𝒫_x[x_s]F(e^{x_u},).$$
(A.17)
We will not deal directly with the above path integral representation of the option price. From (2.8) we see that we can extract the payoff function out of the path integral if we instead focus on the path integral representation of the joint pdf $`P(x_u,|x_t)`$. In the next subsection we will see how to calculate this function.
### A.1 Calculating the joint pdf
The joint pdf introduced in (2.8) can be simply expressed as the path integral
$$P(x_u,|x_t)=_{x_t}^{x_u}𝒟\widehat{x}_s𝒫_x[\widehat{x}_s]\delta (\widehat{}).$$
(A.18)
Using the Fourier representation of the delta function and (2.5), we find that (A.18) becomes
$$P(x_u,|x_t)=\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}𝑑ke^{ik}_{x_t}^{x_u}𝒟\widehat{x}_s𝒫_x[\widehat{x}_s]\mathrm{exp}\left[ik_t^u𝑑sw(s)f(\widehat{x}_s,s)\right].$$
(A.19)
Substituting (A.13) into (A.19) and using
$$_t^u𝑑s\dot{x}_sg^{}(x_s)=g(x_u)g(x_t),$$
(A.20)
we find that (A.19) becomes
$$P(x_u,|x_t)=\frac{1}{2\pi }\mathrm{exp}\left[\frac{g(x_t)g(x_u)}{\sigma ^2}\right]_{\mathrm{}}^{\mathrm{}}𝑑ke^{ik}K(x_u,x_t;T),$$
(A.21)
where $`K`$, which we will refer to as the propagator, is defined by
$$K(x_u,x_t;T)=_{x_t}^{x_u}𝒟x_s\mathrm{exp}\left[\frac{1}{2\sigma ^2}_t^u𝑑s\left[\dot{x}_s^2+V(x_s,s)\right]\right]$$
(A.22)
and what we will call the potential function $`V`$ is
$$V(x_s,s)=g^2(x_s)\sigma ^2g^{\prime \prime }(x_s)2ik\sigma ^2w(s)f(x_s,s).$$
(A.23)
Of special interest is the geometric Brownian motion model defined by (2.2). Comparing (2.2) with (A.1), we can identify $`g^{}(x_t)`$ with $`\mu `$ and we find the joint pdf (A.21) becomes (2.9) with the potential (2.11). Because the drift term is constant in this case, it can be extracted out of the path integral and we are left with an imaginary potential.
## Appendix B Calculating the Jacobian
In this appendix we will calculate the Jacobian (A.10). In discrete time the stochastic differential equation (sde) (A.2) becomes
$$\frac{x_nx_{n1}}{\epsilon }=g^{}(\stackrel{~}{x}_n)+\sigma (\stackrel{~}{x}_n)\zeta _n,$$
(B.1)
where it has been generalized to include a non-constant diffusion coefficient. From this discrete sde we see that
$$\frac{\zeta _n}{x_m}=0,\mathrm{for}m>n.$$
(B.2)
Therefore the Jacobian matrix is triangular and from (A.10) we obtain
$$𝒥=\underset{n=1}{\overset{N}{}}\frac{\zeta _n}{x_n}.$$
(B.3)
From (B.1) we find
$$\frac{1}{\epsilon }\frac{x_n}{\zeta _n}=\varphi g^{\prime \prime }(\stackrel{~}{x}_n)\frac{x_n}{\zeta _n}+\varphi \zeta _n\sigma ^{}(\stackrel{~}{x}_n)\frac{x_n}{\zeta _n}+\sigma (\stackrel{~}{x}_n),$$
(B.4)
where we have used
$$\frac{\stackrel{~}{x}_n}{\zeta _n}=\frac{\stackrel{~}{x}_n}{x_n}\frac{x_n}{\zeta _n}=\varphi \frac{x_n}{\zeta _n}$$
(B.5)
obtained from (A.8). Multiplying (B.4) by $`\epsilon \frac{\zeta _n}{x_n}`$ we obtain
$$\frac{\zeta _n}{x_n}=\frac{1+\epsilon \varphi g^{\prime \prime }(\stackrel{~}{x}_n)\epsilon \varphi \zeta _n\sigma ^{}(\stackrel{~}{x}_n)}{\epsilon \sigma (\stackrel{~}{x}_n)}.$$
(B.6)
For small $`\epsilon `$ (B.6) becomes
$`{\displaystyle \frac{x_n}{\zeta _n}}`$ $``$ $`\epsilon \sigma (\stackrel{~}{x}_n)\left(1\epsilon \varphi g^{\prime \prime }(\stackrel{~}{x}_n)+\epsilon \varphi \zeta _n\sigma ^{}(\stackrel{~}{x}_n)\right)`$ (B.7)
$``$ $`\epsilon \sigma (\stackrel{~}{x}_n)\mathrm{exp}\left[\epsilon \varphi \left(\sigma ^{}(\stackrel{~}{x}_n)\zeta _ng^{\prime \prime }(\stackrel{~}{x}_n)\right)\right].`$
Using this result we find
$$\underset{n=1}{\overset{N}{}}\frac{x_n}{\zeta _n}\epsilon ^N\left(\underset{n=1}{\overset{N}{}}\sigma (\stackrel{~}{x}_n)\right)\mathrm{exp}\left[\epsilon \varphi \underset{n=1}{\overset{N}{}}\left(\sigma ^{}(\stackrel{~}{x}_n)\zeta _ng^{\prime \prime }(\stackrel{~}{x}_n)\right)\right].$$
(B.8)
In the continuous limit $`\epsilon 0`$, we therefore find
$$𝒥^1=\epsilon ^N\left(\underset{n=1}{\overset{N}{}}\sigma (\stackrel{~}{x}_n)\right)\mathrm{exp}\left[\varphi _t^u𝑑s\left(\sigma ^{}(x_s)\zeta _sg^{\prime \prime }(x_s)\right)\right].$$
(B.9)
From the continuous sde (A.2) we have
$$\zeta _s=\left(\dot{x}_s+g^{}(x_s)\right)/\sigma (x_s).$$
(B.10)
Substituting this into (B.9) we find
$$𝒥=\epsilon ^N\left(\underset{n=1}{\overset{N}{}}\sigma (\stackrel{~}{x}_n)\right)^1\mathrm{exp}\left[\varphi _t^u𝑑s\left((\dot{x}_s+g^{}(x_s))\sigma ^{}(x_s)/\sigma (x_s)g^{\prime \prime }(x_s)\right)\right].$$
(B.11)
This agrees with equation 2.4.37 in Stratonovich who derives the Jacobian for a general multi-factor system of sde’s. For a constant diffusion coefficient (B.11) reduces to (A.12).
## Appendix C Improved Short-time Propagator
The goal of this appendix is to obtain an expansion of $`K(x_u,x_t;T)`$ in powers of $`T`$. The derivation here is inspired by the cumulant method used in connection with the quantum statistical density matrix .
We must first write the propagator in the Fourier path integral representation (fpir). We decompose the paths as
$$x_\tau =\overline{x}_\tau +\left(\frac{4\sigma ^2T}{\pi }\right)^{1/2}\underset{n=1}{\overset{\mathrm{}}{}}\frac{z_n\mathrm{sin}(n\pi \tau )}{n}$$
(C.1)
where
$$\overline{x}_\tau =\tau (x_ux_t)+x_t,\tau =(st)/T.$$
(C.2)
The term $`\overline{x}_\tau `$ in (C.1) is the straight line path connecting $`x_t`$ and $`x_u`$. The remaining terms are the harmonic perturbations about the straight line path. With this we can now write
$$\frac{1}{2\sigma ^2}_t^u𝑑s\left[\dot{x}_s^2+V(x_s,s)\right]=\frac{(x_ux_t)^2}{2\sigma ^2T}+\pi \underset{n=1}{\overset{\mathrm{}}{}}z_n^2+\frac{T}{2\sigma ^2}_0^1𝑑\tau V(x_\tau ,\tau ).$$
(C.3)
Summing over all paths between $`x_t`$ and $`x_u`$ is equivalent to integrating over all possible values of the Fourier coefficients $`\{z_n\}`$. This means we can write
$$_{x_t}^{x_u}𝒟x_s=(\mathrm{constant})\times _{\mathrm{}}^{\mathrm{}}𝑑z_1\mathrm{}𝑑z_{\mathrm{}},$$
(C.4)
where the constant is some Jacobian factor resulting from the change of integration variables. Substituting (C.4) and (C.3) into (2.10), we obtain the fpir of the propagator
$$K(x_u,x_t;T)=K_f(x_u,x_t;T)_{\mathrm{}}^{\mathrm{}}𝑑z_1\mathrm{}𝑑z_{\mathrm{}}\mathrm{exp}\left[\pi \underset{n=1}{\overset{\mathrm{}}{}}z_n^2\frac{T}{2\sigma ^2}_0^1𝑑\tau V(x_\tau ,\tau )\right]$$
(C.5)
where $`K_f(x_u,x_t;T)`$ is given by
$$K_f(x_u,x_t;T)=\left(\frac{1}{2\pi \sigma ^2T}\right)^{1/2}\mathrm{exp}\left(\frac{(x_ux_t)^2}{2\sigma ^2T}\right).$$
(C.6)
The constant Jacobian factor in (C.4) must equal the prefactor of (C.6) to obtain the correct propagator when $`V=0`$ (clearly the Jacobian is independent of the potential). We can rewrite (C.5) as
$$K(x_u,x_t;T)=K_f(x_u,x_t;T)E\left[\mathrm{exp}\left(\frac{T}{2\sigma ^2}_0^1𝑑\tau V(x_\tau ,\tau )\right)\right]_{\{z_n\}},$$
(C.7)
where the expectation is with respect to the set of independent Gaussian random variables $`\{z_n\}`$ with zero mean and variance $`1/2\pi `$. The fpir can also be set up using a reference potential that is quadratic .
We can use the fpir (C.7) to derive an expansion in time for the short time propagator. We have using the cumulant expansion
$$K(x_u,x_t;T)=K_f(x_u,x_t;T)\mathrm{exp}\left(\underset{m=1}{\overset{\mathrm{}}{}}\frac{1}{m!}\left(\frac{T}{2\sigma ^2}\right)^mC_m(x_u,x_t;T)\right),$$
(C.8)
where
$$C_1=M_1,C_2=M_2M_1^2,C_3=M_33M_1M_2+2M_1^3,\mathrm{}$$
(C.9)
and
$$M_m=E\left[\left(_0^1𝑑\tau V(x_\tau ,\tau )\right)^m\right]_{\{z_n\}}.$$
(C.10)
The propagator (C.8) has the general structure defined by (3.2) and (3.5). Consider the first cumulant which is
$$C_1(x_u,x_t;T)=_0^1𝑑\tau E\left[V(x_\tau ,\tau )\right]_{\{z_n\}}.$$
(C.11)
We know from (C.1) that $`x_\tau `$ is a sum of Gaussian random variables which itself must be Gaussian distributed. So when calculating the expectation (C.11), we can replace $`x_\tau `$ by
$$x_\tau =\overline{x}_\tau +p_\tau ,$$
(C.12)
where $`p_\tau `$ is a single Gaussian random variable with variance
$$\nu _\tau ^2=\frac{2\sigma ^2T}{\pi ^2}\underset{n=1}{\overset{\mathrm{}}{}}\frac{\mathrm{sin}^2(n\pi \tau )}{n^2}=\sigma ^2T(1\tau )\tau .$$
(C.13)
Using this result we find that
$$E\left[V(x_\tau ,\tau )\right]_{\{z_n\}}=_{\mathrm{}}^{\mathrm{}}𝑑p_\tau P(p_\tau )V(\overline{x}_\tau +p_\tau ,\tau )$$
(C.14)
where
$$P(p_\tau )=\frac{1}{\sqrt{2\pi \nu _\tau ^2}}\mathrm{exp}(p_\tau ^2/2\nu _\tau ^2).$$
(C.15)
The first cumulant then becomes
$$C_1(x_u,x_t;T)=_0^1𝑑\tau _{\mathrm{}}^{\mathrm{}}𝑑p_\tau P(p_\tau )V(\overline{x}_\tau +p_\tau ,\tau ).$$
(C.16)
Expanding the potential around $`p_\tau =0`$ and performing the Gaussian integrals we obtain
$$E\left[V(x_\tau ,\tau )\right]_{\{z_n\}}V(\overline{x}_\tau ,\tau )+\frac{\nu _\tau ^2}{2}V^{\prime \prime }(\overline{x}_\tau ,\tau )+\frac{\nu _\tau ^4}{8}V^{\prime \prime \prime \prime }(\overline{x}_\tau ,\tau )+o(T^3),$$
(C.17)
which from (C.13) is an expansion in time $`T`$.
Consider now the second cumulant
$$C_2(x_u,x_t;T)=_0^1𝑑\tau 𝑑\tau ^{}\left(E\left[V(x_\tau ,\tau )V(x_\tau ^{},\tau ^{})\right]_{\{z_n\}}E\left[V(x_\tau ,\tau )\right]_{\{z_n\}}E\left[V(x_\tau ^{},\tau ^{})\right]_{\{z_n\}}\right).$$
(C.18)
In (C.18) we can replace $`x_\tau `$ and $`x_\tau ^{}`$ with
$$x_\tau =\overline{x}_\tau +p_\tau ,x_\tau ^{}=\overline{x}_\tau ^{}+p_\tau ^{}$$
(C.19)
where $`p_\tau `$ and $`p_\tau ^{}`$ are two correlated Gaussian random variables with variances
$$\nu _\tau ^2=\sigma ^2T(1\tau )\tau ,\nu _\tau ^{}^2=\sigma ^2T(1\tau ^{})\tau ^{}.$$
(C.20)
Using (C.1) we find the covariance is
$$c(\tau ,\tau ^{})=E[p_\tau p_\tau ^{}]=\frac{4\sigma ^2T}{\pi }\underset{n,n^{}=1}{\overset{\mathrm{}}{}}E[z_nz_n^{}]_{\{z_n\}}\frac{\mathrm{sin}(n\pi \tau )\mathrm{sin}(n^{}\pi \tau ^{})}{nn^{}}.$$
(C.21)
Using
$$E[z_nz_n^{}]_{\{z_n\}}=\frac{1}{2\pi }\delta _{nn^{}}$$
(C.22)
and the identity
$$\frac{2}{\pi ^2}\underset{n=1}{\overset{\mathrm{}}{}}\frac{\mathrm{sin}(n\pi \tau )\mathrm{sin}(n\pi \tau ^{})}{n^2}=\tau _l(1\tau _g),$$
(C.23)
we find
$$c(\tau ,\tau ^{})=\sigma ^2T\tau _l(1\tau _g)$$
(C.24)
where $`\tau _l`$ is the lesser of $`\tau `$ and $`\tau ^{}`$ and $`\tau _g`$ is the greater of $`\tau `$ and $`\tau ^{}`$. Using (A.3), we can write the probability distribution of $`p_\tau `$ and $`p_\tau ^{}`$ as
$$P(p_\tau ,p_\tau ^{})=\frac{1}{2\pi }\frac{1}{\sqrt{\nu _\tau ^2\nu _\tau ^{}^2c^2}}\mathrm{exp}\left[\frac{(\nu _\tau ^{}^2p_\tau ^22cp_\tau p_\tau ^{}+\nu _\tau ^2p_\tau ^{}^2)}{2(\nu _\tau ^2\nu _\tau ^{}^2c^2)}\right].$$
(C.25)
We therefore find that
$$E\left[V(x_\tau ,\tau )V(x_\tau ^{},\tau ^{})\right]_{\{z_n\}}=_{\mathrm{}}^{\mathrm{}}𝑑p_\tau 𝑑p_\tau ^{}P(p_\tau ,p_\tau ^{})V(\overline{x}_\tau +p_\tau ,\tau )V(\overline{x}_\tau ^{}+p_\tau ^{},\tau ^{}).$$
(C.26)
Substituting this result and (C.14) into (C.18), we find that the second cumulant is
$`C_2(x_u,x_t;T)`$ $`=`$ $`{\displaystyle _0^1}d\tau d\tau ^{}({\displaystyle _{\mathrm{}}^{\mathrm{}}}dp_\tau dp_\tau ^{}P(p_\tau ,p_\tau ^{})V(\overline{x}_\tau +p_\tau ,\tau )V(\overline{x}_\tau ^{}+p_\tau ^{},\tau ^{})`$ (C.27)
$``$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}dp_\tau dp_\tau ^{}P(p_\tau )P(p_\tau ^{})V(\overline{x}_\tau +p_\tau ,\tau )V(\overline{x}_\tau ^{}+p_\tau ^{},\tau ^{})).`$
Expanding the second cumulant around $`p=0`$ we find that the first order remaining term is
$$C_2(x_u,x_t;T)\sigma ^2T_0^1𝑑\tau 𝑑\tau ^{}V^{}(\overline{x}_\tau ,\tau )V^{}(\overline{x}_\tau ^{},\tau ^{})\tau _l(1\tau _g)+o(\sigma ^4T^2).$$
(C.28)
This follows since the integral we have to do is just the covariance $`c`$.
It can be shown that the $`m`$th cumulant is of order $`(\sigma ^2T)^{m1}`$. Therefore from (C.8) we see that $`C_2`$ starts contributing at order $`T^3`$, while $`C_3`$ starts contributing at order $`T^5`$. This means that to get the expansion (C.8) correct to order $`T^3`$, we need only expand $`C_1`$ to order $`T^2`$ and $`C_2`$ to order $`T`$ as has been done in (C.17) and (C.28). This means that $`C_1`$ alone will give the correct propagator to order $`T^2`$. Substituting (C.28) and (C.17) into (C.8) we find
$`K(x_u,x_t;T)`$ $``$ $`K_f(x_u,x_t;T)\mathrm{exp}[{\displaystyle \frac{T}{2\sigma ^2}}{\displaystyle _0^1}d\tau V(\overline{x}_\tau ,\tau ){\displaystyle \frac{T^2}{4}}{\displaystyle _0^1}d\tau \tau (1\tau )V^{\prime \prime }(\overline{x}_\tau ,\tau )`$ (C.29)
$``$ $`{\displaystyle \frac{\sigma ^2T^3}{16}}{\displaystyle _0^1}𝑑\tau \tau ^2(1\tau )^2V^{\prime \prime \prime \prime }(\overline{x}_\tau ,\tau )`$
$`+`$ $`{\displaystyle \frac{T^3}{8\sigma ^2}}{\displaystyle _0^1}d\tau d\tau ^{}V^{}(\overline{x}_\tau ,\tau )V^{}(\overline{x}_\tau ^{},\tau ^{})\tau _l(1\tau _g)+o(T^4)],`$
which is our desired expansion of the propagator to third order in time.
Acknowledgements: I would like to thank the Australia Research Council for their generous support of this research through an Australian Postdoctoral Research Fellowship. I would also like to thank Science & Finance and the School of Economics and Finance, University of Technology, Sydney, where part of this research was carried out.
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# Dust Streamers in the Virgo Galaxy M86 from Ram Pressure Stripping of its Companion VCC 882
## 1 Introduction
M86 (NGC 4406) is a bright elliptical (E3/S0) galaxy located in the Virgo Cluster at a distance of $`18.3`$ Mpc (Capaccioli et al. 1990). It has a redshift of $`227`$ km s<sup>-1</sup> (Binggeli et al. 1985), while the mean heliocentric velocity of the cluster as a whole is 1050 km s<sup>-1</sup> (Binggeli et al. 1993). M86 is thought to be on an orbit passing through the core of the cluster approximately every 5 billion years (Forman et al. 1979). It is an X-ray object and appears to be a weak radio source (Laing et al. 1983; Fabbiano et al. 1992; Rangarajan et al. 1995). A plume of X-ray, HI, and infrared emission from M86 suggests that its interstellar medium was swept back by the ram pressure from its motion through the intracluster medium (Forman et al. 1979; Fabian, Schwarz, & Forman 1980; Takeda, Nulsen, & Fabian 1984; Bregman & Roberts 1990; Knapp et al. 1989; White et al. 1991). There is also an optical asymmetry to M86 that gives it a slightly enhanced emission along the plume (Nulsen & Carter 1987).
Here we discuss a nucleated dwarf elliptical galaxy, VCC 882 (NGC 4406B; Binggeli et al. 1985), that lies just to the northeast of M86, inside its projected stellar halo. Deep CCD images show a 28 kpc long dust trail inside M86 that appears to follow VCC 882 in its orbit. This trail is possibly the result of ram pressure stripping of gas originally inside VCC 882 that was removed by the high pressure of its motion through the hot gaseous halo of M86. The gas mass obtained from the extinction in the trail is consistent with this former connection to VCC 882. Other evidence for an interaction between M86 and VCC 882 could be an isophotal twist in the central 8 to 80 arcsec of M86 (Bender & Mollenhoff 1987), and the asymmetric outer isophotes of M86 (Nulsen & Carter 1987).
Gas stripping is pervasive within this core region of the Virgo cluster. Three spirals are close to the M86/VCC 882 pair in projection: NGC 4438, NGC 4388, and NGC 4402. There is also a nearby dwarf galaxy, IC 3355. The spirals have been severely stripped of their outer HI disks (Warmels 1986; Hoffman, Helou, & Salpeter, 1988; Cayatte et al. 1990), but the dwarf looks normal. The peculiar negative velocity of M86 is also shared by NGC 4402, NGC 4438, and IC 3355, as if at least some of these galaxies are comoving like a group through the Virgo cluster (Kotanyi & Ekers 1983). Detailed studies of one of these galaxies, NGC 4438, suggest that ram pressure from its motion through the Virgo intracluster medium has visibly distorted its outer disk (Arp 1966), pushing the gas and star formation off the normal plane, and producing a short diffuse trail of radio continuum and X-ray emission (Kotanyi & Ekers 1983; Kotanyi, van Gorkom & Ekers 1983). A second galaxy in this group, NGC 4388, has been studied by Pogge (1988), Petitjean & Durret (1993) and Veilleux et al. (1999), with mixed results on a stripped origin for extraplanar material. The spiral galaxy NGC 4569, at an equal distance from M87 but on the eastern side, has a negative velocity too, along with an anemic classification (van den Bergh 1976) and depletion in HI (Cayatte et al. 1990). Another dwarf galaxy, IC 3475, is closer in projection to M87 than IC 3355 and is highly stripped of HI (Vigroux et al. 1986), while another anemic spiral, NGC 4548, is about twice as far from M87 as these others and has a distorted outer HI disk (Vollmer et al. 1999).
Ram pressure seems to have stripped the Virgo spiral NGC 4694 also. This galaxy is far to the east of the core region in Virgo. It has a 36 kpc long trail of HI streaming off to the west, with a linear velocity gradient along the trail that smoothly connects it with the galaxy (van Driel & van Woerden 1989). A very faint dwarf galaxy is also in the trail, where the HI column density peaks.
All of these cases suggest stripping from the motion of a galaxy through the hot, low-density, intracluster medium inside Virgo. The case discussed here differs because VCC 882 was apparently stripped by its motion through the much denser hot gas associated with the elliptical galaxy M86. Its relative speed is just as large as in the other cases, exceeding 1000 km s<sup>-1</sup>, but the ram pressure it felt must have been much larger because of the higher ambient density.
A similar case is the Virgo dwarf galaxy UGC 7636, which was apparently stripped as it moved through the giant elliptical, NGC 4472. The evidence for this is: a gas cloud to the side of UGC 7636 (Sancisi, Thonnard & Ekers 1987) that has the right mass to have been formerly part of UGC 7636 (Patterson & Thuan 1992; Irwin & Sarazin 1996); a 30 kpc long optical trail of luminous debris adjacent to the cloud and the dwarf (McNamara et al. 1994); absorption of the elliptical galaxy X-ray radiation by the HI cloud (Irwin & Sarazin 1996), and an oxygen abundance in an HII region of the cloud that is consistent with the abundance expected for UGC 7636 (Lee et al. 2000).
Previous photometric studies of Virgo cluster galaxies were made by Binggeli et al. (1985, 1993), Bender & Mollenhoff (1987), and Caon et al. (1994). Katsiyannis et al. (1998) co-added thirteen Schmidt exposures to produce a deep R-band image of the Virgo southeast region, from which they studied extended and overlapping halos. The dust trail discussed here did not appear in these other images.
## 2 Observations and Data Reduction
We obtained 92 images in B, V, and I over an approximately 1 square degree field in the southeast center of the Virgo Cluster using the Burrell Schmidt 0.6-m telescope at Kitt Peak National Observatory on 17-22 March 1999. The plate scale is 1.6 arcsec per pixel, or about 142 pc per pixel at the assumed distance of 18.3 Mpc. The images were reduced using standard IRAF procedures. Flat fields were made from global sky flats taken over the entire observing run. The images were mosaicked to produce combined images with total exposure times of approximately 6.5 hours in each band.
Figures 1a and 1b show a mosaic of the central square degree of the images in B band. The left-hand image is low-contrast, showing the giant ellipticals M84 and M86 as well as the small companion VCC 882, which is normally lost in the bright light of M86. The right-hand image is shown with high contrast to emphasize the outer extents of the large galaxies.
Enhanced images of M86 are in Figure 2. The top left frame shows the logarithm of the intensity in B band with the angular scale indicated. Dark dust features are labeled. Feature A is 2.05 arcmin east of center, and Feature B is 3.7 arcmin southeast of center, midway between two bright foreground stars. These features are more prominent in the other enhanced images. In the top right, which is a (B-I) color map, Feature A has a fork at its mid-point and a faint extension, marked C, towards VCC 882. The lower left frame is an unsharp-masked image, made by smoothing the original with a Gaussian filter 5 pixels wide and subtracting this from the original. Most of the stars have been removed from this image by replacing the corresponding pixels with the average surrounding background. A star at the top of the eastern dust feature remains, as well as another due west of the center of M86 and one near VCC 882. The lower right-hand frame is an ”embossed” B-band image made in Adobe Photoshop; the three-dimensional perspective is the result of an apparent ”illumination” from the west. Here, the dust features are black against the gray background.
Previous published images of M86 did not show the dust features present in Figure 2. They are too faint to appear on high-contrast prints. Bregman & Roberts (1990) mapped this region in HI and reported a private communication with J. L. Tonry, who noticed dust patches $`100^{\prime \prime }`$ east of center. These are probably what we see here. The associated gas did not show up in H I emission because it was below Bregman & Roberts’ detection limit.
At the assumed distance of 18.3 Mpc, Features A and B lie at projected distances of 10.9 kpc and 19.7 kpc from the center of M86. The latter is comparable to the 19.5 kpc radius of M86 at a surface brightness of 25 mag arcsec<sup>-1</sup> (de Vaucouleurs et al. 1991). The outermost halo of M86 shown in the right-hand image of Figure 1 has a much larger radius of 62 kpc. Feature A is approximately 130<sup>′′</sup> ($`=11.5`$ kpc) long from north to south, and its average width is 6.4<sup>′′</sup> ($`=570`$ pc). Feature B is $`63^{\prime \prime }`$ ($`=5.5`$ kpc) long and $`19^{\prime \prime }`$ ($`=1.7`$ kpc) wide, while Feature C is $`27^{\prime \prime }`$ ($`=2.4`$ kpc) long and $`4.8^{\prime \prime }`$ ($`=430`$ pc) wide. The total length of the trail is 28 kpc.
In order to estimate the extinction of the dust features, we determined the magnitude difference between them and their surrounding regions in each passband. East-west intensity cuts were made at 28 points along the length of the trail; 17 of these cuts that are on Feature A are shown in Figure 3. The midpoint of the cut was at the approximate position of the trail. The dashed lines are averages of all 17 cuts. The magnitude differences between the dust and the surrounding regions varied from 0 to 0.2 in B band and from 0 to 0.15 in V band, with uncertainties of about 0.02 mag. The dust features do not show up well in I band.
Figure 4 shows the magnitude differences in the dust features for the B and V bands as a function of position from south to north. The ratio of these magnitude differences, $`\mathrm{\Delta }m_B/\mathrm{\Delta }m_V`$, is shown at the top. Least squares fits to feature A are shown as dotted lines. The endpoints of the curves are the southeastern dust Feature B and the northernmost extension, Feature C, near VCC 882. The magnitude differences decrease from south to north in both B and I bands. The ratio of the magnitude differences in B and V bands is slightly larger than in a Whitford reddening law, where it would be 1.3. The data are not accurate enough to tell if this ratio changes along the length of the trail. It should decrease with greater depth inside the halo of M86 because foreground stars wash out the color difference in the trail, but such a decrease can be offset by changes in intrinsic extinction. The opacity and depth of the dust trail are estimated in Section 3 from simple radiative transfer.
Figure 5 shows the B and I-band radial profiles along the major and minor axes of VCC 882. The profiles were made after first removing the underlying light from M86 by rotating the image $`180^{}`$ and subtracting it from the original. The profiles are well-fit by an exponential, which is the case for many Virgo dwarf ellipticals (Vader & Chaboyer 1994; Ryden et al. 1999). The scale length is 3.8$`{}_{}{}^{\prime \prime }=340`$ pc. The (B-I) color profiles along the major and minor axes are shown at the bottom of Figure 5 as solid and dashed lines. The southern side of VCC 882, nearest M86, is bluer than the northwest side, while the color profile along the major axis is symmetric and flat.
The isophotal contours of VCC 882 are in Figure 6. They are symmetric like the color profile, also showing no evidence of a tidal tail. In contrast, the interaction of the dwarf galaxy UGC 7636 and the giant elliptical NGC 4472 shows a blue tidal tail in the dwarf (Patterson & Thuan 1992). An ellipse fit for VCC 882 gives a position angle of $`116^{}\pm 4^{}`$ and an ellipticity of $`0.26\pm 0.02`$. The inner $`2^{\prime \prime }`$ region in Figure 6 shows an isophotal twist.
## 3 Grain Size, Opacity and Geometric Depth
The magnitude difference between the dust feature and the surrounding field can be converted into an opacity and a relative depth through the M86 halo stars if the intrinsic ratio of $`\mathrm{\Delta }m_B/\mathrm{\Delta }m_V`$ is known. Denoting this ratio by $`C`$ ($`=4/3`$ for a Whitford reddening law), we can write the observed intensities in V and B bands in terms of the intensities of the halo light behind and in front of the cloud, $`I_0`$ and $`I_1`$, respectively, and in terms of the opacity in $`V`$, $`\tau _V`$:
$$I_V^{cloud}=I_0e^{\tau _V}+I_1;I_B^{cloud}=XI_0e^{C\tau _V}+XI_1$$
(1)
where $`X`$ is the ratio of B to V intensity to the side of the cloud:
$$I_V^{side}=I_0+I_1;I_B^{side}=XI_0+XI_1.$$
(2)
The magnitude differences may be written in the form $`I_V^{cloud}/I_V^{side}=10^{0.4\mathrm{\Delta }m_V}`$. Thus we can express the opacity $`\tau _V`$ in terms of the magnitude differences:
$$\frac{1e^{C\tau _V}}{1e^{\tau _V}}=\frac{110^{0.4\mathrm{\Delta }m_B}}{110^{0.4\mathrm{\Delta }m_V}},$$
(3)
and we can express the relative depth of the cloud in the M86 halo as
$$\frac{I_1}{I_0}=\frac{10^{0.4\mathrm{\Delta }m_V}e^{\tau _V}}{110^{0.4\mathrm{\Delta }m_V}}.$$
(4)
Figure 7 shows the V-band opacity and the relative depth of feature A for four values of $`C`$, using the linear fits to $`\mathrm{\Delta }m_V`$ and $`\mathrm{\Delta }m_B`$ that are in figure 4. The opacity generally decreases to the north, and the depth increases. The $`C`$ values are all required to be larger than the maximum of $`\mathrm{\Delta }m_B/\mathrm{\Delta }m_V`$, because this maximum is the value $`C`$ would have for a foreground cloud. In the figure, the value of $`C`$ is given in terms of this maximum, which is 2.36. The solution is not reliable when $`C`$ is close to the maximum, as shown by the sudden decrease of the depth in the northern part of the trail for the $`C=max(\mathrm{\Delta }m_B/\mathrm{\Delta }m_V)`$ case. Large values of $`C`$ suggest that the dust grains are smaller in the trail than they are in the Milky Way. The measurements are too inaccurate to be conclusive, however.
## 4 Discussion
The dwarf galaxy VCC 882 appears to have lost much of its interstellar medium during a recent passage through the X-ray emitting gas of M86. The variation of extinction and depth along the trail of this debris, and the positive relative velocity of VCC 882, suggest that the orbit of the dwarf took it from the lower-left foreground of M86 to the upper middle or background. Ram pressure stripping is the mostly likely cause for the gas removal, rather than tidal stripping, because the stellar distribution in the dwarf is hardly affected by the encounter, and ram-pressure stripping affects only the gas.
The total mass of the dust features was estimated by assuming a standard HI column density of $`1.9\times 10^{21}`$ cm<sup>-2</sup> for A<sub>V</sub>=1 mag (Bohlin, Savage & Drake 1978), and multiplying the resulting column density by the total area of the features. For a magnitude difference of A<sub>V</sub>=0.1 (Fig. 4), the above measurements give a total mass of $`3.4\times 10^7`$ M. If the dwarf is deficient in dust, as is common for small galaxies with low metallicities, then the gas mass in the trail will be larger. For a metallicity of \[Fe/H\]=$`0.59\pm 0.42`$ (Brodie & Huchra 1991), a proportionately lower extinction-to-gas ratio would increase the trail mass by a factor of $`4`$ to $`1.4\times 10^8`$ M. We shall use this mass below.
If the radiative transfer model is correct, then the opacity in the trail can be larger than 0.1 mag; Figure 7 suggests it might be as large as 1 mag. The gas mass would then increase in proportion. We do not consider this model to be accurate, however, because of measurement errors in the magnitude differences and because the results are very sensitive to the unknown ratio $`C`$. More accurate measurements of the trail should improve this situation.
The timescale for the stripping can be estimated from the trail length, which is $`28`$ kpc, and the relative speed of VCC 882 inside M86, $`v_{rel}=1328`$ km s<sup>-1</sup> (the heliocentric velocity of VCC 882 and M86 are 1101 km s<sup>-1</sup> and -227 km s<sup>-1</sup>, respectively). The ratio of these numbers, 15 My, is a measure of the timescale for the trail to be deposited, aside from projection effects.
There is no measure of the mass of VCC 882, but we estimate the absolute magnitude in B band to be $`16`$ mag based on a comparison of our counts with photometric sources in the field. For comparison, observations by Binggeli, Sandage, & Tammann (1985), Harris (1991), and Cohen (1988) convert to absolute B magnitudes of $`14.6`$, $`16.2`$, and $`15.2`$, respectively, using a distance of 18.3 Mpc. Our estimate of $`16`$ mag makes the galaxy luminosity $`L_B3.8\times 10^8`$ L. If the gas in the trail was formerly part of the VCC 882 galaxy, then the ratio of the HI mass to the stellar luminosity was about unity, which is comparable to that for dwarf irregular galaxies (Roberts 1969; Swatters 1999) but high for nucleated ellipticals. This makes the proposed VCC 882 predecessor resemble ESO 359-G29, which is a gas-rich nucleated, dwarf elliptical (Sandage & Fomalont 1993). For a normal stellar population with $`M/L_B2`$ M/L (Sandage & Fomalont 1993; Hirashita, Takeuchi, & Tamura 1998), the galaxy mass would be $`M_{gal}8\times 10^8`$ M. This, combined with the exponential scale length of $`R_{gal}=340`$ pc, gives a characteristic virial speed of $`GM_{gal}/(5R_{gal})^{1/2}=45`$ km s<sup>-1</sup>.
The presence of globular clusters around VCC 882 (Cohen 1988), as well as the symmetric isophotal contours in the outer parts of this galaxy (Fig. 6), suggest that ram pressure stripping of the gas may not have been accompanied by significant tidal disruption of the stars. However, tidal disruption detaches stars and globular clusters slowly from the gravitational pull of VCC 882; it need not have moved the stars very far yet from the galaxy center. If we multiply the 15 My interaction time by the 45 km s<sup>-1</sup> virial speed inside VCC 882, we get a plausible drift distance of $`675`$ kpc, to within a factor of$`3`$ uncertainty from projection effects. This is not a large enough distance for the globulars to have migrated significantly from the center of VCC 882, considering that globulars are usually seen in the outer parts of galaxies anyway. Thus the dwarf could have been exposed to a significant tidal force from M86, but we would not necessarily have noticed the effects of this force yet.
The ram pressure force on the gas can have a very different effect than the tidal force on the stars. The average density in the X-ray halo of M86 is $`n0.01`$ cm<sup>-3</sup> (Thomas et al. 1987), and the VCC 882 orbital speed discussed above is $`v_{rel}=1328`$ km s<sup>-1</sup>. Thus the average ram pressure on the VCC 882 gas is $`P=n\mu v_{rel}^23\times 10^6`$ K cm<sup>-3</sup> for mean atomic weight $`\mu =2.2\times 10^{24}`$ g. Such a pressure can strip the ISM from VCC 882 all at once if it exceeds the gravitational binding energy density that the gas would have had there. This energy density comes from the former gas density multiplied by the square of the escape velocity in the dwarf. This condition for stripping may be written
$$\rho _{\mathrm{external}}v_{rel}^2>\rho _{\mathrm{internal}}v_{\mathrm{escape}}^2,$$
(5)
which is comparable to the conditions written by Gunn & Gott (1972) for disk stripping and Takeda et al. (1984) for spheroid stripping. We take the escape speed to be $`2^{1/2}`$ times the virial speed derived above.
The effective density that the gas had when it was inside VCC 882, if this is where it came from, can be estimated from the current gas mass $`M1.4\times 10^8`$ M in the trail and from the exponential scale length $`R_{gal}=340`$ pc of the galaxy. This naively gives a density of $`3M/(4\pi R_{gal}^3\mu )26`$ cm<sup>-3</sup> for a spherical distribution. This is high for an average ISM density, but it is not supposed to be the real density, only the effective density for the stripping calculation. With this density, the gravitational energy density gives a self-binding pressure of $`1.7\times 10^7`$ K cm<sup>-3</sup>, which is $`5\times `$ larger than the ram pressure.
We get about the same result if we consider a more realistic disk-like structure for the gas before it was stripped. Then the expression for self-binding pressure is $`2\pi G\sigma _{tot}\sigma _{gas}`$ (Gunn & Gott 1972) for total and gaseous column densities $`\sigma _{tot}`$ and $`\sigma _{gas}`$. For an exponential disk with central surface density $`\sigma _0`$ and an exponential scale length $`R_{gal}`$, the total mass out to infinity is $`2\pi \sigma _0R_{gal}^2`$. Thus $`\sigma _{tot}=M_{tot}/(2\pi R_{gal}^2)`$, $`\sigma _{gas}=M_{gas}/(2\pi R_{gal}^2)`$, and the self-binding pressure is $`\left(GM_{gal}/R_{gal}\right)M_{gas}/\left(2\pi R_{gal}^3\right)`$. This is $`(5/3)\rho _{gas}v_{esc}^2`$ for $`\rho _{gas}=3M_{gas}/\left(4\pi R_{gal}^3\right)`$ as above. Thus the stripping condition given by equation (5) with an effective galactic gas density crudely derived for a sphere the size of the exponential scale length is the same to within 60% as Gunn & Gott’s condition for a disk with the gas distributed exponentially.
In either case, the results suggest that the ram pressure from the motion of VCC 882 through M86 was not overwhelming compared to the internal gravitational energy density of the former interstellar medium in that galaxy. If it were, it would probably have led to a more rapid loss of gas from the dwarf, in one or two big clumps (as for UGC 7636). Instead, the stripping was apparently mild, and the morphology of the stripped gas more filamentary, like a long trail of debris removed somewhat steadily.
In fact, this morphological dependence on stripping strength is in agreement with general theory. Stevens, Acreman, & Ponman (1999) found that long gas trails like the one we observe here result if the stripping pressure is modest rather than overwhelming, and if the galactic stars continuously replenish the interstellar medium at normal rates through winds and supernovae. Earlier calculations by Lea & De Young (1976) got a trail for this case too; they explained it as the result of expansion of galactic gas into the low pressure zone downstream in the trail. Takeda, Nulsen & Fabian (1984) found a rapid initial loss of gas for their models, but this was followed at long times by a trail of more steady ablation consisting of gas from the stellar replenishment. Balsara, Livio & O’Dea (1994) got a trail in the mass-replenishment case too, but found also that stripping occurred in bursts. For steady state models, trails occur when there is a low rate of gas replenishment and a high ram pressure (Portnoy, Pistinner, & Shaviv 1993).
The VCC 882 case may have had stripping bursts too, in addition to a more steady gas loss that formed the overall trail. Figure 2 suggests that the trail is not perfectly uniform, but has three prominent condensations. These are approximately equally spaced, and consistent with the suggestion by Balsara, Livio & O’Dea (1994) and Stevens, Acreman & Ponman (1999) that the bursts of mass loss would be periodic. The timescale for these bursts should depend on the turbulent or sound crossing time inside the dwarf galaxy, out to the shocks at the interacting surfaces (cf. Lea & De Young 1976), because internal adjustments and pulsations in the interstellar medium of the dwarf have this characteristic scale. The positions of these shocks depend on the flow speed, so the oscillation timescales should also scale with the flow crossing time outside the galaxy, which is $`2R_{replenish}/v_{rel}`$. Here $`R_{replenish}`$ is the mass-replenishment half radius defined by Stevens et al. (1999). Specific timescales given in these previous models range from 3 to 5 $`\times 10^7`$ years, which was several flow-crossing times (see also Abadi, Moore, & Bower 1999). Here the flow-crossing time is much shorter, $`2R_{gal}/v_{rel}0.5`$ My, so several flow times is perhaps slightly over 1 My. Considering that the whole trail may be only 15 My old (uncertain because of projection effects), this oscillation timescale is short enough to be consistent with the occurrence of three condensations along the trail behind VCC 882.
Other aspects of the trail morphology might result from variable pressures in the M86 halo, i.e., increasing with depth, or from Kelvin-Helmholtz instabilities along the trail of gas. The wiggles in Feature A cannot result from the rotation of VCC 882 because the rotation period at one exponential scale length is $`50`$ My, which is longer than the expected trail formation time.
The long term evolution of the trail is presumably dominated by evaporation and mixing with the hot halo of M86. The evaporation time given by Veilleux et al. (1999) provides a useful guide:
$$t_{evap}=1000n_{trail}R_{trail,pc}^2\left(T_{M86}/10^7K\right)^{5/2}\left(\mathrm{ln}\mathrm{\Lambda }/30\right)yr;$$
(6)
$`\mathrm{\Lambda }`$ is the Coulomb logarithm. The trail density, $`n_{trail}`$ is unknown, but we can use a column density of $`1.9\times 10^{20}`$ cm<sup>-2</sup> from the estimated extinction of A$`{}_{V}{}^{}=0.1`$ mag to give $`n_{trail}R_{trail,pc}=30`$ (half the total because $`R`$ is a radius rather than the full depth of the trail). For $`R_{trail,pc}`$ comparable to the observed half-width of $`250`$ pc, the evaporation time is $`7.5`$ My, which is half the estimated trail age. Thus the trail length could be determined in part by evaporation.
It is conceivable that VCC 882 is bound to M86 and produces a continuous trail from replenished gas that always evaporates by the time it reaches $`1020`$ My downstream. This can be checked from the orbit speed in the potential of M86. If we assume that the two galaxies are currently separated by their projected distance of $`D=86.9^{\prime \prime }=7.7`$ kpc, and that VCC 882 is in a parabolic orbit with a speed of $`v_{rel}=1328`$ km s<sup>-1</sup> (ignoring projection effects), then the mass of M86 becomes $`0.5Dv_{rel}^2/G1.6\times 10^{12}`$ M. The log of the blue luminosity of M86 is 10.65 (Tully 1988), which then gives an M/L ratio of 35 M/L. Although the average M/L for elliptical galaxies is 5 to 10, the M/L ratio for M86 was calculated to be 63 based on mass estimates from the X-ray temperatures and densities (Forman et al. 1985). Our value is consistent with their result, considering the uncertainties here. This suggests that much of the apparent speed of VCC 822 could be from its motion inside M86. The visible trail could be only the non-evaporated remnant of a longer trail that has periodically removed fresh gas from VCC 882 as it plunged into the denser regions of M86.
This situation is similar to the model proposed by Takeda et al. (1984) for elliptical galaxies in clusters, applied here on a smaller scale. However, for VCC 882, the gas mass in the trail is probably too large to be entirely the result of stellar wind accumulation in an orbit time around M86. For the replenishment rate of $`dM/dt=1.6\times 10^{11}M_{gal}`$ yr<sup>-1</sup> used by Balsara et al. (1994), the time for $`1.4\times 10^8`$ M of gas to accumulate inside an $`8\times 10^8`$ M galaxy is $`10^{10}`$ years. This would have to be the orbit time if the stripping process were repetitive, but it is such a long time compared to the Hubble time that VCC 882 would have to be on its first orbit. If the gas accumulation time is not shorter than this, then the trail behind VCC 882 is probably a first stripping event. Even longer replenishment times were assumed for elliptical galaxies by Gaetz, Salpeter & Shaviv (1987).
The distinction between these two models of first-time and repetitive stripping is relevant to the problem of the origin of nucleated dwarf ellipticals. If VCC 882 is now experiencing its first stripping event, then it is just now becoming a gas-poor elliptical galaxy and its formation mechanism would have been one involving stripping from a gas-rich dwarf irregular (Lin & Faber 1983), or preferably a BCD, considering its compact bright nucleus (Davidge 1989; Sandage & Hoffman 1991; Vader & Chaboyer 1994). If the gas accumulation rate is faster than the Balsara et al. estimate by a factor of 10, then VCC 882 could have formed as an elliptical long ago by other mechanisms and now be shedding only the gas mass that comes from evolved stars. The concentration of dE,N galaxies in the center of the Virgo cluster suggest an early formation anyway (Gallagher & Hunter 1989), consistent with dwarf formation models in a cold dark matter universe (Silk, Wyse, & Shields 1987).
The origin of dE’s in Virgo is unclear although their properties are well studied. Nucleated dE’s like VCC 882 concentrate toward the center of Virgo, unlike the non-nucleated dE’s and the dI’s (van den Bergh 1986; Ferguson & Sandage 1989). This suggests either an early origin for the dE’s or stripping from a previous gas-rich state. The centralized dE,N’s are also brighter and larger than the peripheral dE’s (Ichikawa et al. 1988), so stripping from a common ancestor will not produce both types of dE’s. The nucleated dE’s like VCC 882 have more globular clusters per unit luminosity than non-nucleated dE’s (Harris 1991), and this also suggests a primordial origin for the dE,N’s, unless the globulars were made during the stripping event. The exponential light profile in VCC 882 is typical for dE’s, which are like dI’s in this respect (Faber & Lin 1983; Caldwell 1983; Binggeli, Sandage & Tarenghi 1984), but the dE’s have more metals than the dI’s (Thuan 1985) so either the dE’s had a distinct origin or early stripping events triggered substantial amounts of star formation. No star formation is seen in VCC 882, however; there is only a slight B-I color gradient along the minor axis (Fig. 5). dE’s are also distinct from E’s and globular clusters on correlation diagrams considering galactic structure parameters (Binggeli, Sandage & Tarenghi 1984; Kormendy 1985), but not considering color (Caldwell 1983; Bothun & Caldwell 1984; Wirth & Gallagher 1984; Zinnecker et al. 1985),
These properties of dE,N galaxies suggest that most of them acquired their morphologies early on, possibly passing through an early BCD phase to get the bright nucleus and/or globular clusters through intense star formation. Indeed, dE,N nuclei resemble bright globulars (Phillips et al. 1996). But in this case the presence of what appears to be a trail following VCC 882 is an anomaly. This galaxy could not have been a BCD only $`15`$ My ago, which is the stripping time considering the trail length and relative velocity of M86, because that would make VCC 882 too blue now. It might have had a mixed morphology before, like ESO 359-G29 (Sandage & Fomalont 1993) or several other dwarf galaxies (Vigroux et al. 1986; Davidge 1989; Sandage & Hoffman 1991; Vader & Chaboyer 1994). Or it could have an anomalously high gas replenishment rate, which would make the trail come from stripped stellar ejecta in a more continuous flow. The lack of obvious blue stars or young stellar clumps in what was supposedly a gas-rich system only several tens of millions of years ago is difficult to understand. Perhaps the trail did not come from VCC 882 at all.
Observations of the metallicity, dust-to-gas ratio, grain properties, and velocities in the trail would be useful, as would more detailed studies of the current stellar populations in VCC 882. High resolution X-ray observations of the region around and behind VCC 882 would be interesting too, to distinguish the trail from the plume and other X-ray emitting gas connected with M86.
## 5 Conclusions
Deep B, V, and I images of the southeast core region of the Virgo cluster reveal a dust trail in the central region of M86, spanning 28 kpc with a thickness of $`500`$ pc. The dust appears to connect to the nucleated dwarf galaxy VCC 882 located just north of M86. We estimate that the dust has a V-band extinction of 0.1 mag, with a total mass of $`10^8`$ M. This dust and the associated gas may have been ram-pressure stripped from VCC 882 over the last 10–20 My as the dwarf galaxy passed through the hot X-ray emitting gas of M86. The mass of gas is somewhat consistent with its former presence inside VCC 882, and the morphology of a trail is consistent with theoretical predictions for relatively weak ram pressures and periodic outflows. The similarity between the evaporation time of the trail and the dwarf orbit time suggest that evaporation is important. The length of the trail could even be determined by steady evaporation at the back end.
A problem with this interpretation is that VCC 882 is a classical nucleated dwarf. It seems to have no recent star formation as if it just had a gas disk, and the gas mass estimate is much too high for such a morphology. If the trail really did come from VCC 882, then it would have had a mixed morphology before the stripping, like several other rare cases that have been studied elsewhere. More observations will be required to determine the origin and history of this gas.
We gratefully acknowledge summer student support from the Wm. Keck Foundation to the Keck Northeast Astronomy Consortium, and thank Vassar College for a research publication grant.
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# 1 Introduction and motivation
## 1 Introduction and motivation
Lately, there has been a renewed interest in field theories on manifolds with boundaries. In general one would expect a nontrivial relation between bulk and boundary dynamics. One such model is the open string theory with a constant $`B`$-field. It turns out that the bulk and boundary properties of this model are quite different . Recently, in the attempt to provide a solid ground for the quantization of the model, it has been proposed to treat boundary conditions as Hamiltonian constraints within the Dirac approach . This idea seems quite powerful and could be applied to a wide range of models.
Let us recall that for systems with boundaries, traditionally the boundary conditions have to be imposed to properly define the functional derivatives in the theory. Specifically, for the Hamiltonian treatment one needs a boundary conditions for the proper definition of the Poisson brackets (symplectic structure on the phase space). However, there is an alternative algebraic approach to the definition of symplectic structures. Using three basic properties of the Poisson bracket (antisymmetry, Leibniz rule and Jacobi identity) and the canonical brackets for momenta and coordinates one can calculate any bracket. Of course a function on the phase space should now be understood as a formal power expansion in momenta and coordinates. This approach to the Poisson bracket is in the spirit of quantum mechanics where the algebraic definitions are the basic ones.
Thus applying the algebraic definition of Poisson bracket one does not have to impose the boundary conditions in order to do concrete calculations. As a matter of fact one sees that to define momentum, hamiltonian and primary constraints formally in many models there is no need to use the boundary conditions. Thus the natural question arises: what is the status of boundary conditions in this framework. In this paper we try to answer this question. The main point which we are going to make is that the boundary conditions can arise in a purely algebraic fashion as Hamiltonian constraints localized on the boundary. As a result all the Dirac machinery may be applied<sup>2</sup><sup>2</sup>2That is true up to certain technical problems which, we believe, can be resolved. to the boundary conditions vs the Hamiltonian boundary constraints. Indeed we need those boundary constraints to make the whole Hamiltonian treatment consistent.
Let us make a few technical remarks. The basic idea is rather naive. Since one can formally define the momentum, hamiltonian, the primary constraints and do calculations with the Poisson brackets without using the boundary conditions we may proceed in a formally along Dirac’s lines using the canonical Poisson brackets . However now we are not allowed to throw away the total derivative terms (the boundary terms). These boundary terms produce the corresponding constraints on the boundary. The usual consistency conditions have to be required for these constraints. To handle the technical side of the idea it is useful to work with constraints $`\mathrm{\Phi }`$ smeared with test functions $`N(x)`$
$$\mathrm{\Phi }[N]=d^dx\mathrm{\Phi }(x)N(x).$$
(1.1)
This is often a convenient notation, especially when one wants to keep track of partial integrations in a calculation. For the case of boundary constraint one assumes that the smearing function has support on the boundary only. (Those assumptions do not effect the formal calculations in any way.) In all calculations we will avoid the functional questions and will concentrate attention on the algebraic aspect of the computations. We will see that there is no ambiguity as soon as a calculation is done using test functions. However to define the Dirac brackets one has to do calculations without the test functions (as far as the author knows) and this may lead to trouble in some cases. To make sense of those boundary Dirac’s bracket in certain cases one should give a mathematically rigorous definition of the relevant objects. We do not do this and just present the formal answer with short comments. We clarify this point by considering concrete examples.
It is worthwhile to make some remarks concerning the status of boundary conditions also within the Lagrangian formalism. The presence of a boundary can spoil some properties which hold in the situation without a boundary. For instance, two classically equivalent actions (in the sense that they reproduce the same equations of motion) can give rise to different boundary conditions. A nontrivial example of this situation is the relation between the Howe-Tucker and Nambu-Goto actions for the open branes. These actions give slightly different boundary conditions. To relate them is a subtle task and depends on the dimension of the background space-time<sup>3</sup><sup>3</sup>3For instance, the strings in two dimensional space-time: The Nambu-Goto action is linear and it gives rise only the Dirichlet boundary conditions. However the Howe-Tucker (Polyakov) action is quadratic and it might produce as well the Neumann boundary conditions.. This is another reason for looking at the Hamiltonian treatment of the boundary conditions.
The main motivation behind the present work is to understand the status of the boundary conditions in a quantum theory especially an interacting one. The Hamiltonian approach has certain advantages when it comes to quantizing a theory. (At least in principle it is clear what one should do.) At the end we will comment on the possible quantum applications of our results.
Let us briefly comment on the literature. In mathematical physics the Hamiltonian systems with boundaries is an old subject (see for a list of references). The main attention has been on possible modifications of the Poisson bracket by surface terms to fulfil general axiomatic properties in the presence of boundaries. A modifed Poisson bracket was defined in and later generalized in . Unlike the present discussion is not formal. Our attitude is conservative. We calculate the Poisson brackets in the standard way and keep track of the boundary terms which we interpret as Hamiltonian constraints. Eventually these Hamiltonian boundary constraints should lead to a Dirac bracket which gives the right symplectic structure for the model.
Since the general idea by itself is simple one and it is difficult to give any general theorems, examples are quite helpful. Thus throughout the paper we consider four examples: the topological sigma model with a boundary, the open string theory with and without a constant $`B`$-field and four dimensional $`U(1)`$ gauge theory on a manifold with a boundary. There is a section for every example and at the end we summarize the results and discuss the problems. In the first section we consider the topological sigma model with boundaries. This example is rather simple and particular. It demonstrates that there is a difference between the functional and algebraic approaches to the Poisson bracket and that the former approach misses some interesting information about the boundary. In next two section we consider the open string theory. We show that the boundary constraints can arise in a purely algebraic fashion from the algebra of constraints. We hope that the discussion in the fourth section will clarify some points in earlier analysises of the problem . We also point out that some problems may arise in the definition of modified symplectic structure on the boundary. The last section is devoted to Euclidean electrodynamics with a topological term on a manifold with boundaries. In the spirit of open string theory with B-field we derive the modification of the symplectic structure on the boundary.
## 2 Topological sigma model with boundaries
Let us consider a topological sigma model with boundaries defined on 2d-dimensional smooth manifold which admits a symplectic structure<sup>4</sup><sup>4</sup>4$`\omega =\omega _{\mu \nu }dX^\mu dX^\nu `$ is a symplectic structure if $`d\omega =0`$ and $`\omega _{\mu \nu }`$ is not degenerate. $`\omega `$
$$S=\frac{1}{2}\underset{\mathrm{\Sigma }}{}d^2\xi \omega _{\mu \nu }(X)_\alpha X^\mu _\beta X^\nu ϵ^{\alpha \beta },$$
(2.2)
where $`\mathrm{\Sigma }`$ is a two-dimensional world-sheet with boundary. In the bulk this model is purely topological and has no local degrees of freedom . For the present purposes we ignore the topological aspects of the model and thus assume that the background space-time manifold can be covered by one patch. It means that we can think of the sympelctic form as an exact two form $`\omega =dA`$ and the action (2.2) becomes
$$S=\underset{\mathrm{\Sigma }}{}𝑑\tau A_\mu (X)\dot{X}^\mu ,$$
(2.3)
which describes the boundary dynamics. These boundary dynamics are trivial ($`\dot{X}^\mu =0`$) and there is just a modification of the Poisson brackets for $`X^\mu `$ since there are 2d second class constraints.
Now we want to try to extract the information about boundary dynamics from the Hamiltonian treatment, starting from the action (2.2). A variation of the action (2.2) gives
$$\delta S=\underset{\mathrm{\Sigma }}{}𝑑\xi ^\alpha \omega _{\mu \nu }\delta X^\mu _\alpha X^\nu +\frac{1}{2}\underset{\mathrm{\Sigma }}{}d^2\xi (d\omega )_{\mu \nu \rho }_\alpha X^\mu _\beta X^\nu ϵ^{\alpha \beta }\delta X^\rho ,$$
(2.4)
where the last term vanishes by itself (see footnote) and the boundary term should vanish as well. Since $`\omega _{\mu \nu }`$ is nondegenerate, one should impose the Dirichlet boundary condition
$$\delta X^\mu |_\mathrm{\Sigma }=0,$$
(2.5)
which simply means that the naive functional derivative with respect to $`X`$ is not defined on the boundary. Thus the functional approach cannot be used to find the boundary dynamics. Instead we may use an algebraic approach to the problem. The action (2.2) produces $`2d`$ constraints
$$\mathrm{\Phi }_\mu [N^\mu ]=𝑑\sigma N^\mu (P_\mu \omega _{\mu \nu }(X)X^\nu ),$$
(2.6)
which give the following Poisson bracket algebra
$$\{\mathrm{\Phi }_\mu [N^\mu ],\mathrm{\Phi }_\nu [M^\nu ]\}=𝑑\sigma \left[(d\omega )_{\mu \nu \rho }N^\mu M^\nu X^\rho \right]N^\mu M^\nu \omega _{\mu \nu }|_0^\pi ,$$
(2.7)
where $`N^\mu `$ and $`M^\nu `$ are test functions and $`\sigma [0,\pi ]`$. In the present calculation (and as well as in the next two sections) we use the following formula
$$\underset{0}{\overset{\pi }{}}𝑑\sigma \underset{0}{\overset{\pi }{}}𝑑\sigma ^{}f(\sigma )g(\sigma ^{})_\sigma (\delta (\sigma \sigma ^{}))=f(\sigma )g(\sigma )|_0^\pi \underset{0}{\overset{\pi }{}}𝑑\sigma f^{}(\sigma )g(\sigma ).$$
(2.8)
which can be easily motivated. The constraints (2.6) are first class in the bulk and second class constraints on the boundary. To simplify the calculations we can do the following. Since we are working on one patch one can assume that the symplectic structure $`\omega _{\mu \nu }`$ is a constant matrix of a special form (due to the Darboux theorem in one patch there are always special coordinates where the symplectic form can be brought to canonical form). Thus because of (2.7) there is a suitable modification of the symplectic structure on the boundary
$$\{X^\mu ,X^\nu \}|_\mathrm{\Sigma }=\omega ^{\mu \nu },$$
(2.9)
where $`\omega ^{\mu \nu }\omega _{\nu \rho }=\delta _\rho ^\mu `$. Furthermore,
$$\{P_\mu ,P_\nu \}|_\mathrm{\Sigma }=\frac{1}{4}\omega _{\mu \nu },\{X^\mu ,P_\nu \}|_\mathrm{\Sigma }=\frac{1}{2}\delta _\nu ^\mu ,\{X^\mu ,X^\nu \}|_\mathrm{\Sigma }=\frac{1}{2}\omega ^{\mu \nu }.$$
(2.10)
It is easy to check that all these brackets have the desired properties (everything should have a trivial bracket with constraints on the boundary). Proceeding along standard lines one finds that $`\delta X^\mu =\dot{X}^\mu `$ equals $`N^\mu `$ in the bulk and zero on the boundary.
The present model is trivial, nevertheless it contains the essence of the general situation of Hamiltonian models with boundaries. It shows that there is a difference between the functional and the algebraic approaches to the symplectic structure. In the algebraic approach the right boundary conditions arise by themselves in a consistent way. In the next sections we consider less trivial examples of this situation.
## 3 Open string theory without a $`B`$-field
Let us consider an open string theory in a flat space-time ($`\eta _{\mu \nu }=(1,1,\mathrm{},1)`$) without antisymmetric background field. The model has the following action
$$S=\frac{1}{2}d^2\sigma \sqrt{h}h^{\alpha \beta }_\alpha X^\mu _\beta X^\nu \eta _{\mu \nu },$$
(3.11)
where $`h^{\alpha \beta }`$ is an auxilary metric. The treatment of the theory is presented in string theory textbooks (for instance ). To the author’s knowledge the canonical treatment of the open string has only been given in the lectures by Henneaux . In this section we would like to have a new look at some well-known facts about open strings. The action (3.11) produces the following boundary condition
$$(\sqrt{h}h^{10}\dot{X}^\mu +\sqrt{h}h^{11}X^\mu )|_{0,\pi }=0.$$
(3.12)
If one starts from the Nambu-Goto action instead then the general boundary condition is $`\dot{X}^\mu X^\mu |_{0,\pi }`$ which is equivalent to (3.12). The condition (3.12) can be rewritten as follows in phase space
$$(\eta _{\mu \nu }X^\nu +\sqrt{h}h^{01}P_\mu )|_{0,\pi }=0,$$
(3.13)
which states that $`\eta _{\mu \nu }X^\nu `$ and $`P_\mu `$ are proportional to each other on the boundary.
Now let us turn to the Hamiltonian analysis of the system. For the model (3.11) the constraints are well known
$$_1[N]=\underset{0}{\overset{\pi }{}}𝑑\sigma P_\mu X^\mu N,[M]=\underset{0}{\overset{\pi }{}}𝑑\sigma (P_\mu \eta ^{\mu \nu }P_\nu +X^\mu \eta _{\mu \nu }X^\nu )M,$$
(3.14)
and they hold at all points including the boundary. Since the system is generally covariant the naive Hamiltonian vanishes identically. Both constraints (3.14) are first class and they correspond to reparametrizations of the two dimensional world sheet. The constraints obey the following Poisson bracket algebra
$`\{_1[N],_1[M]\}=_1[NM^{}N^{}M],`$ (3.15)
$`\{_1[N],[M]\}=[NM^{}N^{}M]+NM(P_\mu \eta ^{\mu \nu }P_\nu X^\mu \eta _{\mu \nu }X^\nu )|_0^\pi ,`$ (3.16)
$`\{[N],[M]\}=_1[4(NM^{}N^{}M)].`$ (3.17)
The bracket between $`_1`$ and $``$ gives rise the boundary term which should be set to zero to make the Hamiltonian treatment consistent. Since $`_1`$ and $``$ hold everywhere we must require the following constraints on the boundary
$$P_\mu \eta ^{\mu \nu }P_\nu |_{0,\pi }=0,X^\mu \eta _{\mu \nu }X^\nu |_{0,\pi }=0,P_\mu X^\mu |_{0,\pi }=0.$$
(3.18)
One might call them the boundary constraints. The next step should be to check whether the algebra of new constraints is closed or not. As we said before all calculations can be done in a formal way avoiding questions of regularization. For example let us introduce the following notation for the boundary constraints
$$\varphi _1[N]=\underset{0}{\overset{\pi }{}}𝑑\sigma NP_\mu \eta ^{\mu \nu }P_\nu ,\varphi _2[M]=\underset{0}{\overset{\pi }{}}𝑑\sigma MX^\mu \eta _{\mu \nu }X^\nu $$
(3.19)
where $`N`$ and $`M`$ might be thought as test functions localized on the boundary (or around boundary if there is some regularization assumed). This kind of assumptions does not effect the formal calculations. For instance we calculate the following brackets
$$\{\varphi _1[N],\varphi _2[M]\}=4\underset{0}{\overset{\pi }{}}𝑑\sigma NM^{}P_\mu X^\mu +4\underset{0}{\overset{\pi }{}}𝑑\sigma NMP_\mu X^{\prime \prime \mu }4NMP_\mu X^\mu |_0^\pi ,$$
(3.20)
and see that secondary constraints arise. However the constraints (3.18) can be resolved since $`P`$ and $`X^{}`$ are null vectors on the boundary and they are orthogonal to each other there is a proportionality relation on the boundary
$$(\alpha P_\mu +\eta _{\mu \nu }X^\nu )_{0,\pi }=0,$$
(3.21)
where $`\alpha `$ is some proportionality constant which is subject to gauge condition (since it relates world-sheet density to the world-sheet vector). The conditions (3.21) give us the same information as one would get from the Lagrangian formalism (3.13). Hence the whole system can be described as two first class constraints $`_1`$, $``$ plus a set of second class boundary constraints (3.21). The constraints (3.21) are second class because of the non vanishing brackets
$$\{\mathrm{\Phi }_\mu [N^\mu ],\mathrm{\Phi }_\nu [M^\nu ]\}=\alpha \underset{0}{\overset{\pi }{}}𝑑\sigma [N^\nu M^\mu N^\mu M^\nu ]\eta _{\mu \nu },$$
(3.22)
where $`\mathrm{\Phi }_\mu [N^\mu ]`$ is (3.21) smeared with the test function $`N^\mu `$. Proceeding formally for the second class constraints (3.21) we define the corresponding Dirac brackets
$$\{X^\mu (\sigma ),X^\nu (\sigma ^{})\}=\frac{\alpha }{2}\eta ^{\mu \nu }\frac{1}{_\sigma }\delta (\sigma \sigma ^{}),$$
(3.23)
as well as the brackets
$$\{X^\mu (\sigma ),P_\nu (\sigma ^{})\}=\frac{1}{2}\delta _\nu ^\mu \delta (\sigma \sigma ^{}),\{X^\mu (\sigma ),X^\nu (\sigma ^{})\}=\frac{\alpha }{2}\eta ^{\mu \nu }\delta (\sigma \sigma ^{}).$$
(3.24)
We are interested in the restriction of these brackets to the boundary and it is not clear how to find this, especially for the non-local bracket (3.23). The point is that this question can not be answered unless our description of the model is supplemented with a certain amount of additional information. The extra information concerns the restrictions on the behaviour of the fields in order to make operator $`_\sigma `$ invertable (in general there is a constant zero mode for this operator). Therefore to make further progress one needs more insight into the model. It would be interesting to quantize the free open string theory in a nonconformal gauge (where $`\alpha 0`$) and calculate the commutators (3.23), (3.24) explicitly on the boundary. Resolving this kind of questions can lead to the proper understanding of the foundations of Witten’s open string field theory where the noncommutativity of the ends of strings plays a crucial role.
## 4 Open string theory with a constant $`B`$-field
Now let us turn to the open string theory with a constant $`B`$-field. The model has the following action
$$S=\frac{1}{2}d^2\sigma (\sqrt{h}h^{\alpha \beta }_\alpha X^\mu _\beta X^\nu \eta _{\mu \nu }ϵ^{\alpha \beta }_\alpha X^\mu _\beta X^\nu B_{\mu \nu }).$$
(4.25)
This system has attracted much attention recently because of the noncommutative properties of the end points of the string. A treatment of the model has been given in (also see for the quite full list of references). Let us just recall that in the Lagrangian formalism one should impose the boundary conditions
$$(\sqrt{h}h^{10}\eta _{\mu \nu }\dot{X}^\nu +\sqrt{h}h^{11}\eta _{\mu \nu }X^\nu +B_{\mu \nu }\dot{X}^\nu )|_{0,\pi }=0,$$
(4.26)
which have the following form in phase space
$$(B_\mu {}_{}{}^{\nu }P_{\nu }^{}+G_{\mu \nu }X^\nu +\sqrt{h}h^{01}P_\mu )|_{0,\pi }=0,$$
(4.27)
where $`G_{\mu \nu }=\eta _{\mu \nu }B_{\mu \sigma }B^\sigma _\nu `$. For the sake of simplicity we assume that $`B`$ is a non-degenerate matrix (for the degenerate case one can easily generalize all the following arguments).
Now we turn to the Hamiltonian formalism. In the usual fashion the constraints are
$$_1[N]=\underset{0}{\overset{\pi }{}}𝑑\sigma P_\mu X^\mu N,$$
(4.28)
$$[M]=\underset{0}{\overset{\pi }{}}𝑑\sigma (P_\mu \eta ^{\mu \nu }P_\nu 2P_\mu B^\mu {}_{\nu }{}^{}X_{}^{\nu }+X^\mu G_{\mu \nu }X^\nu )M.$$
(4.29)
These are first class constraints and they hold everywhere including at the boundary points. Next we calculate the algebra keeping track of the boundary terms. The constraints obey the following Poisson bracket algebra
$`\{_1[N],_1[M]\}=_1[NM^{}N^{}M],`$ (4.30)
$`\{_1[N],[M]\}=[NM^{}N^{}M]+NM(P_\mu \eta ^{\mu \nu }P_\nu X^\mu G_{\mu \nu }X^\nu )|_0^\pi ,`$ (4.31)
$`\{[N],[M]\}=_1[4(NM^{}N^{}M)].`$ (4.32)
To make the theory consistient one should set the boundary term to zero. Since $`_1`$ and $``$ hold everywhere there is a boundary constraint
$$X^\mu (B_\mu {}_{}{}^{\nu }P_{\nu }^{}+G_{\mu \nu }X^\nu )|_{0,\pi }=0,$$
(4.33)
which is the difference between $`_1`$ and the boundary term in (4.31). One cannot solve the system as simply as before. Therefore we proceed along Dirac’s lines . We look at possible secondary and tertiary constraints and then try to separate them into first and second class constraints. Sometimes, before separating them into different classes it is helpful to solve some of them.
We thus have to calculate brackets of all constraints including the boundary one and see if new constraints arise. We will perform the calculations in a formal way and introduce the following notation for the boundary constraint
$$\mathrm{\Phi }[N]=\underset{0}{\overset{\pi }{}}𝑑\sigma NX^\mu (B_\mu {}_{}{}^{\nu }P_{\nu }^{}+G_{\mu \nu }X^\nu ),$$
(4.34)
where $`N`$ is a test function. As a result of the computations some new constraints will arise. Let us look at some of them to see the pattern. We have
$$\{\mathrm{\Phi }[N],\mathrm{\Phi }[M]\}=\underset{0}{\overset{\pi }{}}d\sigma [NM^{}N^{}M]X^\mu B_\mu {}_{}{}^{\rho }(B_\rho {}_{}{}^{\nu }P_{\nu }^{}+G_{\rho \nu }X^\nu ).$$
(4.35)
Introducing the following notation for the new constraint
$$\mathrm{\Phi }_1[N]=\underset{0}{\overset{\pi }{}}d\sigma NX^\mu B_\mu {}_{}{}^{\rho }(B_\rho {}_{}{}^{\nu }P_{\nu }^{}+G_{\rho \nu }X^\nu ),$$
(4.36)
we get
$$\{\mathrm{\Phi }_1[N],\mathrm{\Phi }_1[M]\}=\underset{0}{\overset{\pi }{}}d\sigma [NM^{}N^{}M]X^\mu B_\mu {}_{}{}^{\delta }B_{\delta }^{}{}_{}{}^{\sigma }B_{\sigma }^{}{}_{}{}^{\rho }(B_\rho {}_{}{}^{\nu }P_{\nu }^{}+G_{\rho \nu }X^\nu ),$$
(4.37)
and
$$\{\mathrm{\Phi }_1[N],\mathrm{\Phi }[M]\}=\underset{0}{\overset{\pi }{}}d\sigma [NM^{}N^{}M]X^\mu B_\mu {}_{}{}^{\delta }B_{\delta }^{}{}_{}{}^{\rho }(B_\rho {}_{}{}^{\nu }P_{\nu }^{}+G_{\rho \nu }X^\nu ),$$
(4.38)
and so on. This suggests the following boundary conditions
$$X^\mu M_\mu {}_{}{}^{\sigma }(B_\sigma {}_{}{}^{\nu }P_{\nu }^{}+G_{\sigma \nu }X^\nu )|_{0,\pi }=0,$$
(4.39)
where $`M`$ is some power of $`B`$. Since $`B`$ is nondegenerate and antisymmetric all these conditions can be replaced by the following one
$$(B_\mu {}_{}{}^{\nu }P_{\nu }^{}+G_{\mu \nu }X^\nu +\beta P_\mu )|_{0,\pi }=0,$$
(4.40)
where $`\beta `$ is the coefficient of proportionality which is subject to a gauge condition (like $`\alpha `$ in the previous section). We will see that (4.40) are second class constraints. Introducing the notation
$$𝒦_\mu [N^\mu ]=\underset{0}{\overset{\pi }{}}𝑑\sigma N^\mu (B_\mu {}_{}{}^{\nu }P_{\nu }^{}+G_{\mu \nu }X^\nu +\beta P_\mu )$$
(4.41)
it is easy to check the brackets
$$\{𝒦_\mu [N^\mu ],𝒦_\nu [M^\nu ]\}=\underset{0}{\overset{\pi }{}}𝑑\sigma [N^\mu M^\nu N^\nu M^\mu ](B_\mu {}_{}{}^{\rho }G_{\rho \nu }^{}+\beta G_{\mu \nu })$$
(4.42)
where we have nondegenerate matrix on the right-hand side. Therefore we conclude that to make the whole Hamiltonian treatment consistent one must impose the boundary conditions (4.27) which play the role of second class constraint on the boundary. Otherwise the algebra (4.32) would not be closed. Thus the bracket algebra has to be modified on the boundary. The Poisson bracket must be replaced by the Dirac bracket. For example on the boundary the coordinates have the following bracket
$$\{X^\mu ,X^\nu \}_\mathrm{\Sigma }=B^\mu {}_{\sigma }{}^{}(G^1)_{}^{\sigma \nu }+\beta (nonlocalpart)$$
(4.43)
where the non local part has the same structure as in the previous section. For the case $`\beta =0`$ (for instance, conformal gauge or static gauge) the brackets (4.43) are well defined. A discussion of the modified brackets is given in .
## 5 Electrodynamics with topological term
As a last example we consider theory with a nonvanishing Hamiltonian. We will take a look at four dimensional Euclidean electrodynamics with a topological term. The action is defined by
$$S=\frac{1}{2g^2}\underset{}{}FF+\frac{i\theta }{4\pi ^2}\underset{}{}FF,$$
(5.44)
where we use differential forms. Equivalently, in components,
$$S=\frac{1}{4g^2}\underset{}{}d^4xF_{\mu \nu }F^{\mu \nu }+\frac{i\theta }{16\pi ^2}\underset{}{}d^4xϵ_{\mu \nu \rho \sigma }F^{\mu \nu }F^{\rho \sigma }.$$
(5.45)
The theory is defined on a manifold $``$ with non empty boundary $``$. Since we are interested in the Hamiltonian treatment we assume that $`=R\times \mathrm{\Sigma }`$ where $`\mathrm{\Sigma }`$ is a spatial manifold. For the sake of simplicity we further assume that $`\mathrm{\Sigma }`$ is closed set in $`R^3`$ and thus it carries a flat metric. This assumption is not essential and the whole logic can be generalized to the general curved case.
Before looking at the Hamiltonian formalism we briefly consider the Lagrangian formalism. To the author’s knowledge this system has not been separately studied, except in . The action (5.44) gives the following equations of motion
$$dF=0,$$
(5.46)
which should be supplemented by the boundary condition
$$\underset{}{}\delta A(\frac{1}{g^2}F+\frac{i\theta }{2\pi ^2}F)=0.$$
(5.47)
To proceed further let us write (5.47) in components
$$𝑑t\underset{\mathrm{\Sigma }}{}d^2s\left[n^a(\frac{1}{g^2}E_a+\frac{i\theta }{2\pi ^2}B_a)\delta A^0n_bϵ^{abc}\delta A_c(\frac{1}{g^2}B_a+\frac{i\theta }{2\pi ^2}E_a)\right],$$
(5.48)
where we introduce the standard notation $`E_aF_{0a}`$ and $`B_a=\frac{1}{2}ϵ_{abc}F^{bc}`$ and $`n^a`$ is a vector normal to $`\mathrm{\Sigma }`$. Now it is straightforward to read off the boundary conditions
$$n^a(\frac{1}{g^2}E_a+\frac{i\theta }{2\pi ^2}B_a)|_{}=0or\delta A^0|_{}=0,$$
(5.49)
$$n_bϵ^{abc}(\frac{1}{g^2}B_a+\frac{i\theta }{2\pi ^2}E_a)|_{}=0or\delta A_c|_{}=0.$$
(5.50)
Thus one of these two sets of conditions should be imposed to make the Lagrangian treatment consistent<sup>5</sup><sup>5</sup>5Apart from this one can think about other physical requirements such as absence of energy-momentum flow through the boundary. In fact the conservation of energy-momentum requires the conditions on the left hand side of (5.49) and (5.50) ..
Let us rewrite the boundary conditions (5.49), (5.50) in phase space. The momentum is defined as follows
$$\pi _a=\frac{1}{g^2}E_a+\frac{i\theta }{2\pi ^2}B_a,$$
(5.51)
and there is the usual constraint $`\pi _0=0`$ which we will discuss later on. Using (5.51) the boundary conditions (5.49), (5.50) become
$$n^a\pi _a|_{}=0or\delta A^0|_{}=0,$$
(5.52)
$$n_bϵ^{abc}\left(\frac{i\theta }{2\pi ^2}g^2\pi _a+\left(\frac{1}{g^2}(\frac{i\theta }{2\pi ^2})^2g^2\right)B_a\right)|_{}=0or\delta A_c|_{}=0.$$
(5.53)
There are one normal condition (on the left handside (5.52)) and two tangential conditions (on the left handside (5.53)). We will keep these in mind. We hope to find them as boundary constraints required to make the whole Hamiltonian treatment consistent. Let us assume that one can choose a coordinate system such that the normal vector has the form $`\stackrel{}{n}=(1,0,0)`$.
We now turn to the Hamiltonian treatment. Using (5.44) and (5.51) one defines the Hamiltonian
$$H=\underset{\mathrm{\Sigma }}{}d^3x\left[\frac{g^2}{2}\pi _a\pi ^a\frac{i\theta }{2\pi ^2}g^2\pi _aB^a\frac{1}{2}\left(\frac{1}{g^2}(\frac{i\theta }{2\pi ^2})^2g^2\right)B_aB^a+(_aA_0)\pi ^a\right].$$
(5.54)
Thus one defines the Hamiltonian $`H`$ and primary constraint $`\pi _0`$ without using the boundary conditions. Introducing the notation
$$\mathrm{\Pi }[\mathrm{\Lambda }]=\underset{\mathrm{\Sigma }}{}d^3x\mathrm{\Lambda }(x)\pi ^0(x),$$
(5.55)
one has the following bracket
$$\{\mathrm{\Pi }[\mathrm{\Lambda }],H\}=𝒢[\mathrm{\Lambda }]\underset{\mathrm{\Sigma }}{}d^2s\mathrm{\Lambda }(n^a\pi _a),$$
(5.56)
where $`𝒢[\mathrm{\Lambda }]`$ is the Gauss law constraint
$$𝒢[\mathrm{\Lambda }]=\underset{\mathrm{\Sigma }}{}d^3x\mathrm{\Lambda }(x)_a\pi ^a(x).$$
(5.57)
Consistency then implies that the right hand of (5.56) must be equal to zero. Thus we are getting the standard Gauss law and as well the boundary constraint $`n^a\pi _a=0`$ (if one assumes that $`A^0`$ is zero on the boundary there is no boundary term in (5.56)).
Following the standard prescription one should look at the time evolution of Gauss law (5.57) and the new boundary constraint
$$\mathrm{\Phi }_1[\mathrm{\Lambda }]=\underset{\mathrm{\Sigma }}{}d^3x\mathrm{\Lambda }n^a\pi _a,$$
(5.58)
where $`\mathrm{\Lambda }`$ can be thought as test function with support on $`\mathrm{\Sigma }`$. The formal computation gives us the following result
$$\{𝒢[\mathrm{\Lambda }],H\}=\underset{\mathrm{\Sigma }}{}d^2s[n_b_c\mathrm{\Lambda }ϵ^{abc}]\left(\frac{i\theta }{2\pi ^2}g^2\pi _a+\left(\frac{1}{g^2}(\frac{i\theta }{2\pi ^2})^2g^2\right)B_a\right),$$
(5.59)
$$\{\mathrm{\Phi }_1[\mathrm{\Lambda }],H\}=\underset{\mathrm{\Sigma }}{}d^3x[_c(\mathrm{\Lambda }n_b)ϵ^{abc}]\left(\frac{i\theta }{2\pi ^2}g^2\pi _a+\left(\frac{1}{g^2}(\frac{i\theta }{2\pi ^2})^2g^2\right)B_a\right),$$
(5.60)
where we have used the higher dimensional analog of equation (2.8). The bracket (5.59) is localized on the boundary and it gives us the tangential boundary constraints which exactly coincide with the boundary conditions (5.53). The same boundary constraint is given by the bracket (5.60) and it is localized on the boundary since $`\mathrm{\Lambda }`$ has support on the boundary only. Therefore we introduce the new boundary constraint
$$\mathrm{\Phi }_a[N^a]=d^3xN_a[\frac{i\theta }{2\pi ^2}g^2\pi ^a+\left(\frac{1}{g^2}(\frac{i\theta }{2\pi ^2})^2g^2\right)B^a],$$
(5.61)
where it is assumed that $`N_a=(0,N_2,N_3)`$. To decide on the status of boundary constraints $`\mathrm{\Phi }_1`$, $`\mathrm{\Phi }_2`$ and $`\mathrm{\Phi }_3`$ one should calculate the following brackets
$$\{\mathrm{\Phi }_1[\mathrm{\Lambda }],\mathrm{\Phi }_b[N^b]\}=\left(\frac{1}{g^2}(\frac{i\theta }{2\pi ^2})^2g^2\right)\underset{\mathrm{\Sigma }}{}d^3x\mathrm{\Lambda }n^a^bN^cϵ_{abc},$$
(5.62)
$$\{\mathrm{\Phi }_a[N^a],\mathrm{\Phi }_b[M^b]\}=\frac{i\theta }{2\pi ^2}g^2\left(\frac{1}{g^2}(\frac{i\theta }{2\pi ^2})^2g^2\right)\underset{\mathrm{\Sigma }}{}d^3x^a(N^bM^c)ϵ_{abc}.$$
(5.63)
One notices that the brackets (5.63) of $`\mathrm{\Phi }_a`$ ($`a=2,3`$) are non-zero because of the boundary term on the right hand side of (5.63). In the bulk such brackets are zero since the constraints $`\mathrm{\Phi }_a`$ are generalization of the chiral condition for two forms . Therefore the brackets of all boundary constraints give a field independent antisymmetric matrix with rank 2. It turns out that there is one first class boundary constraint $`n^a\pi _a`$ which makes us able to gauge away the normal component of the connection on the boundary. The boundary constraints $`\mathrm{\Phi }_2`$ and $`\mathrm{\Phi }_3`$ are second class constraints which lead to the the following Dirac bracket on the boundary
$$\{A_2(x),A_3(y)\}|_\mathrm{\Sigma }=\frac{i\theta }{2\pi ^2}g^2\left(\frac{1}{g^2}+\frac{\theta ^2}{4\pi ^4}g^2\right)^1\delta ^{(2)}(xy).$$
(5.64)
In analogy with the models considered in the previous section we see that at the boundary the Poisson brackets should be replaced by the corresponding Dirac bracket. In general the boundary Dirac bracket will depend on the geometry of the boundary. We will discuss this elsewhere. The physical interpretation of (5.64) is unclear. It seems that one will have problems with localizing photons on the boundary. Certainly this subject deserves an independent study and we do no analyse the boundary theory further here.
## 6 Discussion and problems
In this paper we made an attempt to understand the status of the boundary conditions within the Hamiltonian formalism motivated by the quantum theory. We have shown that boundary conditions can arise in a purely algebraic fashion as Hamiltonian boundary constraints. Their existence is necessary to make the whole Hamiltonian treatment consistent. Our arguments were based on four examples: the topological sigma model, the open string theory with and without a B-field and electrodynamics with a topological term. For some systems it is important to motivate that the boundary conditions can be treated as Hamiltonian constraints. This type of systems has non trivial boundary conditions which mix momenta and coordinates. Such boundary conditions change the canonical brackets on the boundary drastically and therefore they are very important for the quantization of the system as whole. However as we saw in some instances (e.g.,(3.23)) problems can arise with the definition of the Dirac bracket on the boundary. To resolve those problems one needs more insight into the models. In other cases there is no ambiguity in defining the modified symplectic structure (e.g., (2.9), local part of (4.43) and (5.64)).
As can be seen from the last example the boundary conditions give rise not only to second class constraints but also to first class constraint. It is unclear how this kind of boundary constraint should be applied in the quantum theory. We hope to return to this question and do explicit calculations for this model in the presence of a simple boundary.
From a technical point of view the present approach is based on rules for dealing with the following expression
$$\underset{\mathrm{\Sigma }}{}d^dx\underset{\mathrm{\Sigma }}{}d^dx^{}f(x)g(x^{})D_xD_x^{}\delta (xx^{})$$
(6.65)
where $`D_x`$ ($`D_x^{}`$) is some differential operator. As was pointed out in the expression (6.65) is not in general well defined on a closed domain. In the present paper we considered simple models with a few derivatives and therefore it was no problem to define (6.65) in a reasonable way. However in general one should address this question more carefully.
It would be interesting to study an interacting theory like Yang-Mills theory and gravity systems in this context. We hope to treat these questions elsewhere.
Acknowledgments
It is pleasure to thank Ingemar Bengtsson, who has promoted my interest in the subject and who has helped a lot during the preparation of this work. I am grateful to Ingermar Bengtsson and Ulf Lindström for reading and commenting on the manuscript. I thank M.M.Sheikh-Jabbari and V.O.Soloviev for bringing the relevant references to my attention.
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# Theoretical Results for Sandpile Models of SOC with Multiple Topplings
## I Introduction
Sandpile models of stick-slip dynamics have received considerable attention as canonical models of self-organized criticality (SOC) . SOC refers to the widespread tendency of many extended, dissipative dynamical systems to evolve inevitably towards a complex state with power-law correlations in space and time: a “critical” state. Of course, a critical state is only one possible example of complex phenomena that can emerge in large, self-organizing systems composed of many strongly interacting parts. No doubt there are other types of complex states that have not yet been so well characterized mathematically, e.g. for example in networks . From this viewpoint, the phenomena of SOC itself is a prototype for how complexity emerges in nature without fine tuning parameters. In spite of the gross simplicity of various cellular models that have been introduced, and hundreds if not thousands of numerical studies of SOC, only minimal analytic understanding has been achieved.
In fact, a survey of analytic works on sandpile models of SOC is exceedingly short. The model of SOC introduced by Bak, Tang, and Weisenfeld (BTW) has yielded to some analytic treatment associated with its abelian properties, primarily due to the work of Dhar and collaborators . The scaling properties of waves, where each site only topples, or releases grains, once has been understood by Priezzhev and collaborators . Nevertheless, the large scale properties of avalanches, where each site can topple many times in response to a single grain being added to the system, remain unsolved and the numerical situation controversial . The same is true for the Zhang model where some limited progress has been made using methods from dynamical systems theory . Recurrent, multiple topplings within an avalanche also appear in most other unsolved sandpile models, such as the stochastic Manna model , the universality class represented by the Oslo rice pile model , cellular models of vortex dynamics , as well as one-dimensional trough models exhibiting multiscaling . The difficulties preventing progress in solving any of these simplified models in particular, or finding general analytic tools for granular systems exhibiting SOC, appear to be related, in part, to the existence of recurrent topplings.
This statement is further supported by the following facts: Dhar and Ramaswamy (DR) introduced a directed version of the BTW model, and solved for the avalanche distribution and many other properties exactly. In the DR model, it can be rigorously proven that no multiple topplings occur. (Consequently, the elegant DR solution, as it has been conceived thus far, does not address the full complexity of discrete or granular models of SOC.) The fixed scale transformation method of Pietronero and collaborators also explicitly ignores the presence of multiple topplings. One consequence of this fact is that this method puts the stochastic Manna model and the BTW model into the same universality class, which is not consistent with most numerical works (except ), including those measuring unequivocal differences in aging behaviors . Multiple topplings, where the activity can return an arbitrarily large number of times, do not appear in any mean field description , since in high enough dimensions, the avalanche activity is not recurrent, or able to return more than a finite number of times , at any site. Multiple topplings are a fluctuation effect associated with self-intersections of the avalanche cluster in space and time , and as such are relevant below some upper critical dimension.
Certainly, the intricacies associated with multiple topplings are not the only ones that present themselves in attempting an analytic treatment of granular models of SOC. For example, the fact that the dissipation process is confined to the boundary, which forces the system to self-organize, is an important and subtle point because the boundary cannot be scaled out in the limit of large system sizes as is usually done in statistical physics. In principle, the boundary is always important, because the incoming sand grains must be transported to it, no matter how large the system size. The broken translational invariance associated with the boundary often leads to long range boundary effects in the metastable states (see, for example, ). It might be useful to pry these complications apart, treating one issue at a time. Here we focus on the problem of recurrent or multiple topplings, and seek a model which does not present other difficulties.
Recently, Pastor-Satorras and Vespignani have studied numerically a stochastic directed sandpile model (SDM), which is a stochastic version of the exactly solvable model introduced by DR. This stochastic model is simpler and presumably unrelated to the directed models introduced and studied by Tadić and collaborators . Pastor-Satorras and Vespignani demonstrated by numerical simulations that the model exhibits multiple topplings which changes the universality class, making it distinct from the DR model. This was accomplished by numerically measuring and comparing various critical exponents characterizing the avalanches. Its close relation to the DR model, which has an exact solution, suggests to us an analytic study.
### A Summary
We proceed with an analysis of the SDM as follows: First we define the DR model and the SDM. For pedagogical reasons, in Section III, we review the proof that the critical state of the DR model is the set of all stable states with equal probability. We also review some necessary parts of Dhar’s construction of an operator algebra for stochastic models, which can be represented as deterministic models with a quenched array of random numbers. Combining these two works, we then show that for the SDM the critical state is also the set of all stable states with equal probability, described by a product measure. Using this fact, we show in Section IV that the SDM can be recast as a generalized branching process propagating in an uncorrelated environment, enabling a study of the infinite system. By carefully analyzing the microscopic dynamics of this process on the lattice, we explicitly derive a discrete dynamical equation for the propagation of flowing grains in avalanches. In Section V, coarse graining this discrete equation gives a continuum equation for avalanches that should describe the large scale properties of any microscopic model with the same symmetry, conservation of grains, and stochastic effects.
Notably, our equation is similar to the Edwards-Wilkinson (EW) equation except that the amplitude of the nonconservative noise is a Heaviside (theta) function of the local activity. Crucially, the noise amplitude is a threshold function, rather than being a constant, such as the temperature. The height of the interface represents the number of topplings in an avalanche. The steady state that is eventually reached in the limit of large times is always the state of no activity where the height of the interface is zero everywhere and the avalanche has died. Thus the equation describing avalanche dynamics corresponds to an absorbing state phase transition where the the transient state is governed by the EW equation in the region where it survives. Section V also describes an analysis of this nonlinear equation. We extract all the (nontrivial) critical exponents for avalanches, i.e. in $`d=1`$, $`D=7/4`$, $`\tau =10/7`$, $`z=2`$, $`\tau _t=D=7/4`$ distinct from the DR model. For $`d2`$, where multiple topplings are not relevant, the critical exponents are the same as in the DR model. All of these results agree perfectly with previous numerical works . We also write down the Fokker-Planck equation for the probability distribution of the number of topplings at each site in an avalanche, although we do not solve it. Finally, we conclude with a brief comment on some possibilities for future analytical work on absorbing state phase transitions and granular models of SOC.
## II Definition of Directed Models
Consider a two dimensional square lattice as shown in Fig. 1. The direction of propagation is labeled by $`t`$, with $`0t<T`$. The transverse direction is labeled by $`x`$, with periodic boundary conditions. Only sites with $`(x+t)`$ even are on the lattice, so that $`x`$ is a positive integer modulo $`2X`$, and the lattice has a total of $`TX`$ sites. On each site, an integer variable $`z(x,t)`$ is assigned. The $`i`$’th grain is added to a randomly chosen site $`x_i`$ on the top row $`t=0`$. There $`z(x_i,0)z(x_i,0)+1`$. When any site acquires a height greater than $`z_c=1`$ it topples, i.e. $`z(x,t)z(x,t)2`$ for $`z(x,t)>z_c`$.
The two models differ with respect to the transmission of grains out of a toppling site. In the DR model, one grain is transferred to the left downstream neighbor and one grain to the right so the toppling rule is for $`z(x,t)>z_c`$
$`z(x,t)`$ $``$ $`z(x,t)2`$ (1)
$`z(x1,t+1)`$ $``$ $`z(x1,t+1)+1`$ (2)
$`z(x+1,t+1)`$ $``$ $`z(x+1,t+1)+1.`$ (3)
For the SDM, on the other hand, each grain from a toppling site is given equal probability to go to any downstream nearest neighbor. In this case, when the site $`(x,t)`$ topples,
$$z(x,t)z(x,t)2$$
and
$`z(x1,t+1)`$ $``$ $`z(x1,t+1)+1`$ (4)
$`z(x+1,t+1)`$ $``$ $`z(x+1,t+1)+1`$ (5)
with probability 1/2, or
$`z(x1,t+1)`$ $``$ $`z(x1,t+1)+2`$ (6)
$`z(x+1,t+1)`$ $``$ $`z(x+1,t+1)`$ (7)
with probability 1/4, or
$`z(x1,t+1)`$ $``$ $`z(x1,t+1)`$ (8)
$`z(x+1,t+1)`$ $``$ $`z(x+1,t+1)+2.`$ (9)
with probability 1/4. Thus, the SDM is a directed version of the model introduced by Manna.
In both directed models, grains are conserved during each toppling event. This is true except at the open boundary $`t=T`$ where toppling sites simply discharge their grains out of the system. Sites are relaxed according to a parallel update until there are no more unstable sites, and the properties of the resulting avalanche are recorded. Then a new avalanche is initiated by adding a single grain to a randomly chosen site on the top row, $`t=0`$. An avalanche can be characterized by its longitudinal extent, $`t_c`$, the largest $`t`$ row affected, its width, $`x_c`$, the largest transverse distance from the avalanche origin to any site affected by the avalanche, its area, $`a`$, the total number of sites affected, its size, $`s`$, the total number of toppling events, and the maximum number of topplings at a site, $`n_c`$.
It is straightforward to generalize this definition to higher dimensions, with the number of directions transverse to the direction of propagation being $`d`$. In this case $`z_c=2d1`$. At a toppling site $`zzz_c1`$. In the DR case each downstream neighbor receives exactly one grain. In the stochastic case, each downstream neighbor has equal probability $`1/2d`$ to receive each grain. For simplicity of notation and concepts we will focus our discussion on the case $`d=1`$ unless otherwise noted.
## III States on the Attractor
For both directed models, any configuration satisfying $`0z(x,t)z_c`$ for all $`(x,t)`$ is stable. The total number of such configurations is $`(z_c+1)^{TX}`$. For clarity, we now review the argument showing that in the steady state, all such stable states are equally likely in the DR model.
### A Review of Some Exact Results by Dhar and Ramaswamy
Let $`C_0`$ be a starting configuration with the $`i`$’th particle added at site $`x_i`$, resulting in the new stable configuration $`C_i`$. Then $`C_i`$ is uniquely determined by the dynamics given $`x_i`$ and $`C_{i1}`$. The crucial point is that this dynamics is invertible. On the top row $`C_i`$ differs from $`C_{i1}`$ only at the site $`x_i`$, with $`z(x_i,0)`$ in $`C_i`$ being more than its value in $`C_{i1}`$ by one (mod2). Other rows in $`C_{i1}`$ are the same as in $`C_i`$ if there was no toppling at $`(x_i,0)`$; otherwise the $`z`$’s in the first row, $`t=1`$, in $`C_{i1}`$ are the same as in $`C_i`$, except at the two downstream neighbors, $`(x_i1,1)`$ and $`(x_i+1,1)`$ of $`(x_i,0)`$ whose heights are less by one (mod2) than their values in $`C_i`$. This obviously continues for subsequent rows. Thus given $`C_i`$ and $`x_i`$ we can uniquely determine $`C_{i1}`$.
For a given $`C_i`$, there are precisely $`X`$ distinct choices of $`C_{i1}`$ and $`C_{i+1}`$ corresponding to $`X`$ distinct, possible choices of $`x_i`$ and $`x_{i+1}`$. The master equation for the evolution of probabilities of configurations, is
$$dP(C)/dt=\underset{C^{}}{}T_{C^{}C}P(C)+\underset{C^{}}{}T_{CC^{}}P(C^{}).$$
(10)
Since there are $`X`$ distinct choices for the $`C^{}`$ into $`C`$ and also for the $`C^{}`$ out of $`C`$, each having probability $`1/X`$, the probability distribution $`P(C_0=a)=cons`$, independent of $`a`$, is invariant in time. Thus the probability distribution of states on the attractor is a product measure, with each site independently occupied with one particle with probability 1/2, otherwise being empty.
In a recent work, Dhar has shown that the stochastic Manna model also exhibits the Abelian property and is a special case of the Abelian Distributed Processors Model. Correspondingly some of the the analytic techniques of the BTW model also apply to the stochastic Manna model. It is only necessary to realize that for the stochastic models, instead of associating probabilities with each toppling, we can assign to each site an infinite stack of random numbers, uniformly distributed between zero and one, say. The quenched random numbers in each site’s stack then determine the allocation of grains during each toppling event. Thus, the $`q`$’th random number at (x,t) determines at the $`q`$’th toppling of that site where the grains will go. There is a one-to-one correspondence between any realization of the dynamics of the stochastic model, and the dynamics of a deterministic system with a random array (chosen appropriately to model the probability distribution of grain allocation), under the same condition of particle additions.
If we specify the height configuration of the sandpile as well as the infinite stack of random numbers at each site, Dhar shows that the model is also Abelian. It is easy to check that given any unstable configuration with two or more unstable sites, we get the same configuration by toppling at an unstable site $`i`$, and then at unstable site $`i^{}`$, as we would get if we first toppled at $`i^{}`$ and then at $`i`$, if the same list of random numbers in the array is provided. Iterating this until a metastable state is reached proves the Abelian property of the model.
### B New Results
The directed stochastic model is also equivalent to a deterministic directed model with an infinite stack of quenched random numbers at each site. Since the latter model is Abelian we can choose to relax each row, one site at a time, until it is completely stable, before going on to the next higher row. In this case, it is easy to see that the deterministic model with quenched random numbers shares the same property of invertibility as the DR model .
Let $`C_0`$ be a starting configuration and $`R(x,t,q)`$ be the infinite array of random numbers, with the initial pointers $`q_0(x,t)=0`$ for all entries $`(x,t)`$. The pointers $`q`$ in the array $`R`$ will move as the sequence of topplings proceeds. The $`i`$’th particle being added at site $`x_i`$ and the current pointers $`q_{i1}(x,t)`$ in the array $`R`$ known, this results in a new stable configuration $`C_i`$, and a new set of pointers $`q_i(x,t)`$ in the fixed array $`R`$. Invertibility follows. In this case we are given the current configuration $`C_i`$, the current set of pointers $`q_i(x,t)`$ in the fixed array $`R`$ and $`x_i`$. In order to prove invertibility we must determine both $`C_{i1}`$ and $`q_{i1}(x,t)`$.
On the top row $`C_{i1}`$ differs from $`C_i`$ only at the site $`x_i`$, with $`z(x_i,0)`$ in $`C_{i1}`$ being less than its value in $`C_i`$ by 1(mod2). If $`z(x_i,0)=1`$ in $`C_i`$, then no toppling occurred and $`C_{i1}`$ is the same as $`C_i`$ at all other sites; also the set of pointers $`\{q_{i1}=q_i\}`$. If $`z(x_i,0)=0`$ then one toppling occurred at that site. We locate the pointer $`q_i(x_i,0)`$ and move it back one step in the stack $`R(x_i,0,q_i(x_i,0))`$ giving $`q_{i1}(x_i,0)=q_i(x_i,0)1`$. This pointer now points to a number that tells us where the two grains were placed. The heights at the sites in the second row $`t=1`$ in configuration $`C_{i1}`$ are the same as those in $`C_i`$ except at the forward neighbors from $`x_i`$ that received a grain according to $`R(x_i,0,q_{i1}(x_i,0))`$. If both sites received a grain then we apply the same procedure to those sites as we applied to $`(x_i,0)`$. If one site receives two grains then that site must have toppled once. Its height in the previous configuration is the same as its height in the current one, and its pointer is moved back by one unit, determining which downstream neighbors receive grains. One continues in this fashion increasing the row $`t`$.
Unlike the DR model, eventually one can encounter a site receiving three or more grains from sites in the previous row. If the total number of grains received at a site, $`n`$, is even then the site must topple exactly $`n/2`$ times. The pointer at that site is moved back $`n/2`$ steps, so $`q_{i1}(x,t)=q_i(x,t)n/2`$, reading the intervening numbers in the stack at that site to determine where the grains from that site are sent. If $`n`$ is odd and in $`C_i`$ the height is one, then the site must have toppled $`(n1)/2`$ times, with its height in $`C_{i1}`$ being 0. Thus $`q_{i1}(x,t)=q_i(x,t)(n1)/2`$. Similarly if $`n`$ is odd and in $`C_i`$ the height is zero, then the site must have toppled $`(n+1)/2`$ times, with its height in $`C_{i1}`$ being 1. Thus $`q_{i1}(x,t)=q_i(x,t)(n+1)/2`$. One reads the intervening sequence in the array of random numbers for that site to determine how many grains each downstream neighbor receives, and so forth. Thus, given $`C_i`$, $`x_i`$, and $`q_i(x,t)`$, with a fixed array $`R(x,t,q)`$, we can uniquely determine $`C_{i1}`$ and $`q_{i1}(x,t)`$. This proves the invertibility of the dynamics of the SDM.
For a given array $`R`$ and set of pointers $`\{q_i\}`$, for any state $`C_i`$ there are precisely $`X`$ distinct choices of $`C_{i1}`$ and $`C_{i+1}`$ corresponding to the $`X`$ possible choices of $`x_i`$ and $`x_{i+1}`$. It then follows, as before, from the master equation for the evolution of probabilities of configurations that the state prepared with a uniform distribution over all stable states is invariant in time.
Thus for the directed Manna model, the self-organized critical state is the set of all stable states with equal likelihood; it is a product measure state, where the probability for a site to be empty is equal to the probability for it to have one grain, which are both equal to 1/2. This is exactly the same as in the DR model, so for the SDM the presence of multiple topplings does not lead to any correlations in the states on the attractor.
## IV Discrete Equation for Avalanches in the Critical State
The fact that the critical state is a product measure state leads to a significant simplification; namely the critical dynamics can be described as a type of generalized branching process. Thus one can simulate or describe avalanches in an infinite system as follows. Consider a site which we we will call the origin. The origin in the equivalent branching process represents the site that receives a grain in the critical state of the SDM. The height at that site is either one or zero with equal probability. Add one grain to it. If the height now is greater than one it topples. Then define the heights at sites (1,1) and (-1,1); they are one or zero with equal probability. They receive grains from the origin according to the stochastic rules of toppling in the directed model, and topple if they are unstable. In this way, one can always construct the lattice as the avalanche propagates and one can simulate the infinite system, albeit always for a finite time.
We define the quantity $`n(x,t)`$ to be the number of grains added to $`(x,t)`$ given that one grain was added to the origin. The total number of grains that leave a site $`n_{out}(x,t)`$ can at most differ by one from the number of grains going in. If $`n(x,t)`$ is even then $`n_{out}(x,t)=n(x,t)`$. If $`n(x,t)`$ is odd, then, since the number of grains which can leave any site is always even $`n_{out}(x,t)=n(x,t)\pm 1`$. The process is critical and the increase or decrease occur with equal probability. Thus we observe there is a source of bounded, nonconservative noise, with a threshhold, in the dynamics of $`n`$ during an avalanche that comes from the presence or absence of grains in the metastable states. Since the number of grains going into a site can only arise as a consequence of grains going into its immediate upstream neighbors we arrive at the following discrete equation
$`n(x,t+1)={\displaystyle \frac{1}{2}}\left(n(x1,t)+n(x+1,t)\right)`$ (11)
$`+\theta _o(n(x+1,t))\eta (x+1,t)+\theta _o(n(x1,t))\eta (x1,t)`$ (12)
$`j(x+1,t)+j(x1,t).`$ (13)
On average each site will get 1/2 of the grains going into its upstream neighbors. There are two sources of stochastic variations from the average. One is conservative: Each upstream neighbor may divide its out flowing grains unevenly between its two downstream sites, but what is taken away from one downstream neighbor is added to the other according to the binomial distribution. This gives a stochastic current $`j`$ which is either directed to the right (here defined as positive) or to the left (here defined as negative) for each site. The first two moments of the stochastic current of flowing grains are, from the binomial distribution,
$`j(x,t)`$ $`=`$ $`0`$ (15)
$`j(x,t)j(x^{},t^{})`$ $`=`$ $`n(x,t)\delta _o(x,x^{})\delta _o(t,t^{}).`$ (16)
Since this is a discrete equation, the Kronecker delta functions, $`\delta _o`$, are defined on the set of integers.
The nonconservative noise is the most interesting and, as we shall see, relevant noise. It is associated with the fact that the metastable states either add or absorb flowing grains from the avalanche. However, as mentioned before, the number of flowing grains can only change by one unit irrespective of the local number of flowing grains as long as it is nonzero. This gives rise to the discrete Heaviside step functions in Eq. 2 defined as $`\theta _o(u)=1`$ for $`u=1,2,3,\mathrm{}`$ and $`\theta (u)=0`$ otherwise. With this convention the nonconservative noise is at each point in space-time either $`\pm 1`$ with equal probability $`1/4`$ or 0 with probability $`1/2`$. Thus
$`\eta (x,t)`$ $`=`$ $`0`$ (17)
$`\eta (x,t)\eta (x^{},t^{})`$ $`=`$ $`{\displaystyle \frac{1}{2}}\delta _o(x,x^{})\delta _o(t,t^{}).`$ (18)
The appropriate initial condition to describe the avalanche is $`n(x,0)=\delta _o(x,0)`$. The avalanche propagates and spreads out; eventually it dies out. Then a new avalanche, represented by a new realization of the branching process is started.
## V Continuum Equation for the Avalanches
One could consider a rigorous derivation of the continuum limit of Eqs. 2-4, taking the lattice size in space, $`\mathrm{\Delta }_x`$, and time $`\mathrm{\Delta }_t`$, as well as the grain size, $`\mathrm{\Delta }_n`$, to zero. Instead, here we invoke the usual “hand-waving”, coarse graining procedure to obtain a smooth function $`n`$ of continuous variables $`x`$ and $`t`$. Expanding to leading order in gradients, and time derivatives, we arrive at
$$\frac{n(x,t)}{t}=\frac{1}{2}^2n(x,t)2\frac{j(x,t)}{x}+2\theta (n(x,t))\eta (x,t),$$
(19)
where the threshold function $`\theta (u)=0`$ for $`u0`$ and $`\theta (u)=1`$ for $`u>0`$. By the central limit theorem, the noise terms are both Gaussian with first and second moments
$`\eta (x,t)`$ $`=`$ $`0`$ (20)
$`\eta (x,t)\eta (x^{},t^{})`$ $`=`$ $`{\displaystyle \frac{1}{2}}\delta (xx^{})\delta (tt^{})`$ (21)
$`j(x,t)`$ $`=`$ $`0`$ (22)
$`j(x,t)j(x^{},t^{})`$ $`=`$ $`n(x,t)\delta (xx^{})\delta (tt^{}).`$ (23)
The appropriate initial condition for the avalanche is $`n(x,0)=\delta (x)`$. The avalanche grows by increasing or decreasing $`n`$ locally where $`n`$ is non-zero. Eventually the avalanche dies and $`n(x,t)=0`$ everywhere. This equation describes the transient out of an absorbing state associated with the avalanche. In particular, the state with no flowing grains $`n(x,t)=0`$ for all $`(x,t)`$ is stable, which is a requirement of any equation describing avalanche dynamics.
### A Analysis
Dimensional Analysis is the simplest tool we can apply, and the first step in any theoretical analysis. The dimension of the conservative noise is $`[j]^2/[x]^2=([n]/[t][x]^3)`$, and dimension of the nonconservative noise is $`[\eta ]^2=(1/[x][t])`$. Thus, as long as $`[n]<[x]^2`$ then the conservative noise is irrelevant with respect to the nonconservative noise. Ignoring this term we arrive at our main result
$$\frac{n(x,t)}{t}=\frac{1}{2}^2n(x,t)+2\theta (n(x,t))\eta (x,t).$$
(24)
In the region covered by the avalanche $`n(x,t)>0`$, the threshold function, $`\theta =1`$ and may be ignored, resulting in an avalanche dynamics described by the linear Edwards-Wilkinson equation . Dimensional analysis then gives the correct scaling of various quantities. Thus for the SDM, $`x_ct_c^{1/z}`$ with $`z=2`$ precisely as in the DR model. However, the Edwards-Wilkinson equation gives a rough surface in one dimension and the maximum number of topplings scales as the transverse extent of the avalanche as $`n_cx_c^{1/2}`$. This differs markedly from the DR model where $`n_c=1`$, independent of the transverse extent, $`x_c`$.
Continuing with our scaling analysis, the area covered by the avalanche is $`ax_ct_ct_c^{3/2}`$ (as in the DR model), but the size of the avalanche includes the extra effect of multiple topplings. The size scales as $`snx_ct_ct_c^{7/4}`$. Since on average for every grain added one grain must be transported the entire length of the system to the open boundary, we have that $`<s>=T`$. Since all the geometric quantities associated with avalanches exhibit scaling, it is reasonable to assume, and can perhaps be proven, that the distribution of avalanche sizes, times, and spatial extent are power laws, namely $`P(s,T)s^\tau f(s/T^D)`$, $`P_t(t,T)t^{\tau _t}g(t/T)`$, and $`P_x(x,T)x^{\tau _x}w(x/T^{1/z})`$. The constraint on the average size then gives $`1=D(2\tau )`$ or $`\tau =10/7`$. Similarly, from conservation of probability, $`\tau _t1=D(\tau 1)=(\tau _x1)/z`$, gives $`\tau _t=D=7/4`$ and $`\tau _x=5/2`$.
It is straightforward to check that Eqs. 2-4 also apply to the case where there are $`d`$ transverse dimensions. The only factors that are changed are various constants. Applying dimensional analysis, we find that in $`d`$ transverse dimensions, $`x_ct_c^{1/2}`$, so that $`z=2`$ and $`n_cx^{\frac{2d}{2}}`$. The upper critical dimension is $`d_c=2`$ above which the maximum number of topplings does not diverge with the size of the avalanche, and mean field results obtain. This corresponds to the fact that the surface described by the EW equation is flat above two dimensions rather than being rough. For $`d2`$, $`sn_cx_c^dtt^{3/2+d/4}`$, i.e. $`D=3/2+d/4`$ and the other exponents are obtained via the above scaling relations giving $`\tau =2(4/(6+d))`$, and $`\tau _t=3/2+d/4`$.
#### 1 The Threshold Term
Outside the region covered by the avalanche, the threshhold function $`\theta `$ has a major effect on the dynamics. In particular, in regions where $`n(x,t)=0`$, the interface is pinned and cannot move. The noise does not act where there are no flowing grains! This is completely different than the usual models of stochastic interfacial growth. The threshhold function importantly breaks the translational symmetry of the EW equation ($`nn+\text{const}.`$) and leads to an absorbing state. Typically absorbing state phase transitions have been considered where the amplitude of the noise depends on the activity $`n`$ to some positive power . Here we find a very weak effect simply distinguishing between having activity and not having it in terms of a threshold function. This effect is so weak that the scaling dimensions of the propagating avalanche are the same as the linear EW equation. Obviously if the threshhold function $`\theta (n)`$ were replaced by $`n^\alpha `$ in Eq. 5 that would no longer be the case. Thus Eqs. 5-6 is a hybrid combining interface dynamics (the number of topplings of the avalanche being the interface), and an absorbing state model. It describes a previously undiscovered absorbing state phase transition.
#### 2 Averaging over Noise
Averaging over avalanches corresponds to averaging Eq. 7 over noise and we arrive at a linear diffusion equation for the average propagation of flowing grains in response to a single grain being added at $`(0,0)`$ to the critical system:
$$\frac{n(x,t)}{t}=\frac{1}{2}^2n(x,t),$$
(25)
whose solution is
$$n(x,t)=\frac{1}{(2\pi t)^{1/2}}e^{x^2/2t}.$$
(26)
Obviously, this solution has the important property of conservation, namely $`𝑑xn(x,t)=1`$ which is required for stationarity. Note that the DR model also obeys exactly the same equation for the average propagation of activity. This equation is enforced by the local conservation and symmetry properties of the system and is in no way related to the presence or absence of multiple topplings. Thus we get that for all models with the same symmetry and conservation of grains, the dynamical exponent $`z=2`$ and the exponent $`\eta =0`$, since the average amount of activity remains constant in time.
#### 3 The Fokker-Planck Equation
Ideally one would like to determine the full probability distribution for the number of topplings $`n(x,t)`$ in avalanches. The dynamics of this probability distribution $`P[n;t]`$ is expressed by the Fokker-Planck equation. The Fokker-Planck equation can be obtained by straightforward means from the Langevin equation (Eqs. 6,7). It is
$`{\displaystyle \frac{P[n;t]}{t}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle }dx{\displaystyle \frac{\delta }{\delta n}}\left\{(^2n)P\right\}`$ (28)
$`+{\displaystyle 𝑑x\frac{\delta }{\delta n}\left\{\theta (n)\left\{\frac{\delta }{\delta n}(\theta (n)P)\right\}\right\}}.`$
Unfortunately we are not currently able to analyze this equation in any significant way.
## VI Outlook for Future Work
A major limitation of the present work is that applies only to a set of directed models where all stable states are equally likely. Unlike most models of SOC, there are no spatial or temporal correlations in the metastable states on the attractor. Even in this drastically simplified setting, the occurrence of multiple topplings has a profound affect on the critical properties of the system, changing the universality class. This fact suggests that any reasonable theory of avalanche dynamics in sandpile models of SOC must treat the effect of multiple topplings.
Avalanches are described by the dynamics of particles which exhibit an absorbing state phase transition. As shown in our main result, Eq. (7), this transition has the new feature that the amplitude of the noise is a threshhold function of the local activity rather than a being a power of the activity. This threshhold form of the noise amplitude has never before been discussed in the literature. It may be interesting to examine a broad range of absorbing state phase transitions with such a threshhold noise amplitude, as well as the implications of the threshhold noise amplitude on avalanche dynamics in SOC.
The picture of avalanches as reaction-diffusion systems with an absorbing state was first suggested in as applicable to SOC systems and later in . In the case discussed here, the particles, representing topplings, are known to propagate in an uncorrelated environment because the probability distribution of metastable states on the attractor is described by a product measure. In the general case, there will be important correlations from the background that must be included along with boundary effects, leading among other things, to correlations between avalanches. It seems to us that Dhar’s development of an operator algebra for stochastic models might provide a fruitful avenue to pursue further research.
Note Added: After submitting this manuscript for publication, another preprint appeared . They also obtain the critical exponents for the stochastic directed model we discuss.
We thank D. Dhar for helpful comments on the manuscript. This work was supported in part by NSF grant DMR-0074613.
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# Resolving the Controversy Over the Core Radius of 47 Tucanae (NGC 104)1footnote 11footnote 1Based on observations with the NASA/ESA Hubble Space Telescope, obtained at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5-26555.,2footnote 22footnote 2Lick Observatory Bulletin No. 139X.
## 1 Introduction
Globular clusters are excellent laboratories for studying the dynamics of a stellar system. The high density of stars near their centers results in frequent interactions—e.g., single star-single star, single star-binary, and binary-binary. Such interactions redistribute energy throughout the cluster and drive its global evolution on the so-called “two-body relaxation” timescale. This can sometimes be comparable to the typical orbital period of stars near the cluster center, and significantly shorter than the cluster age, implying that dynamical evolution is important. The orbital period or crossing/dynamical time is related to the core radius (the characteristic length scale associated with the inner density profile) and the velocity dispersion of the cluster. It is customary to characterize the radial distribution of various stellar populations in terms of the core radius (cf. Rasio (2000)).
The inner stellar density profile of a globular cluster is suggestive of its evolutionary state (Hut (1996)). Most clusters are characterized by constant density cores, well fit by models based on a relaxed, Maxwellian distribution function of stars out to the limiting (tidal) radius (King 1966a ). On the other hand, about twenty percent of all globular clusters appear to have undergone catastrophic gravothermal collapse as a result of runaway energy loss from the core due to two-body relaxation (Djorgovski & King (1986); see the recent review in Meylan & Heggie (1997)). Even these post–core-collapse (PCC) clusters, however, are not expected to develop central singularities as heating by binaries will drive a “quasi-steady post-collapse” stage followed by gravothermal oscillations. Theoretical modeling suggests that clusters in the post-collapse phase should have core radii in the range $`r_{\mathrm{core}}0.01`$$`\mathrm{\hspace{0.25em}0.04}r_h`$ (Goodman (1987); Gao et al. (1991)) or even as large as $`r_{\mathrm{core}}0.09r_h`$ (Vesperini & Chernoff (1994)), where $`r_h`$ is the half-mass radius. These calculations are based on a single-mass stellar population; models with a realistic stellar mass function could yield a larger core radius.
The globular cluster 47 Tucanae has been studied extensively, and has long been considered a prototypical relaxed globular cluster (King (1985)) with a large core radius: $`r_{\mathrm{core}}25^{\prime \prime }`$ (Djorgovski & King (1984)). Recent observations of a variety of exotic objects, such as millisecond pulsars (Robinson et al. (1995)), X-ray sources (Hertz & Grindlay (1983); Hasinger, Johnston, & Verbunt (1994)), and nine stars whose velocities differ from the cluster mean by $`30`$ km $`\mathrm{s}^1`$ (Gebhardt et al. (1995)), have raised the possibility that 47 Tuc may be approaching core collapse. The half-light (or half-mass) radius of 47 Tuc is $`r_h=174^{\prime \prime }`$ (Trager, Djorgovski, & King (1993)), so if it has undergone core collapse it should have a core radius in the range $`1\stackrel{}{\mathrm{.}}7`$ to $`16^{\prime \prime }`$ (following Goodman (1987); Gao et al. (1991); Vesperini & Chernoff (1994)). Two recent determinations of the core radius of 47 Tuc fall within the upper end of this range, as discussed in § 2 below. Moreover, Gebhardt & Fischer (1995) have constructed nonparametric dynamical models based on the surface brightness and velocity dispersion profiles, and conclude that the mass profile of 47 Tuc is as centrally concentrated as that of M15. The density profile slope of 47 Tuc beyond $`2^{}`$ is found to be similar to that seen in Fokker-Planck simulations of PCC clusters (Cohn (1980)).
Ground-based studies typically use integrated surface brightness measurements to determine a cluster’s density profile (cf. Djorgovski & King (1984)). A substantial fraction of the optical light of a cluster comes from a handful of the brightest red giant branch (RGB) stars so that the effective Poisson error associated with integrated light measurements is large. This kind of “sampling” error makes the measured density profile noisy, and any error in the center determination tends to bias the measured core radius towards large values. Hubble Space Telescope (HST) images, even with the aberrated pre-refurbishment point spread function (PSF), offer the advantage of resolving individual stars down to the main sequence turnoff even in the cluster core (Guhathakurta et al. 1992, hereafter GYSB). Star counts are more representative of the stellar mass density than the integrated light, and the effective Poisson errors are smaller (King 1966b ). It is preferable to work at short wavelengths (e.g., the $`U`$ band) where the brightest RGB stars are suppressed relative to the more numerous, bluer faint subgiants: this reduces sampling effects in integrated light measurements and increases the degree of faint star completeness (most notably in the vicinity of bright giants) for star count measurements (Calzetti et al. 1993; De Marchi et al. 1996, hereafter DPSGB).
This paper focuses on the question: What is the core radius of 47 Tuc as defined by the radial distribution of evolved stars? It examines recent measurements of the density profile of 47 Tuc that appear to be in disagreement with one another. In particular, the HST Wide Field/Planetary Camera 1 (WFPC1) data set analyzed by DPSGB is reanalyzed here using somewhat different photometric techniques; our results and those of DPSGB are compared to archival Wide Field Planetary Camera 2 (WFPC2) data. The background of the core radius controversy is given in § 2. In § 3, the available 47 Tuc data sets and the methodologies of this paper and DPSGB are discussed. In § 4 the photometric errors associated with each method are examined, demonstrating that the aperture photometry technique of DPSGB produces a radially-varying bias in the star counts used to measure $`r_{\mathrm{core}}`$. Core radius calculations are presented in § 5 along with tests which show that the core radius discrepancy can be explained in terms of the star count bias in the DPSGB study. The conclusions of this paper are presented in § 6.
## 2 The Controversy to Date
The first HST-based core radius for 47 Tuc was derived from F555W images obtained with the WFPC1 instrument operated in Planetary Camera mode<sup>4</sup><sup>4</sup>4All Wide Field/Planetary Camera 1 (WFPC1) data sets referred to in this paper were obtained in Planetary Camera mode. (GYSB): $`r_{\mathrm{core}}=23^{\prime \prime }\pm 2^{\prime \prime }`$, in good agreement with earlier ground-based measurements (cf. Harris & Racine (1979)). Shortly thereafter, Calzetti et al. (1993) analyzed an independent set of pre-refurbishment HST Faint Object Camera (FOC) ultraviolet data, and obtained results that were in conflict with previous work. They pointed out that the cluster center determined from the ground was significantly in error, biased by a few bright giant stars. More surprisingly, they derived a cluster density profile that did not resemble a King profile, but was instead fit by a superposition of two King profiles with core radii of $`25^{\prime \prime }`$ and $`8^{\prime \prime }`$, suggesting that 47 Tuc is on the verge of core collapse.
The cause of the discrepancy between these density profile measurements remained unclear for the next few years, and so did the true nature of 47 Tuc’s density profile. The center derived by Calzetti et al. (1993) is $`6^{\prime \prime }`$ from the earlier ground-based estimate (Aurière & Ortolani (1988)), but only $`1\stackrel{}{\mathrm{.}}4`$ from GYSB’s estimate which is well within the formal errors. When Calzetti et al.’s FOC ultraviolet image-based cluster centroid estimate is used in conjunction with the GYSB F555W data set, the resulting density profile is marginally smoother than the profile obtained using GYSB’s center estimate, possibly indicating that the former center estimate is more accurate (see § 1), but the best-fit core radius is $`20^{\prime \prime }`$ in both cases (Guhathakurta et al. (1993)). Two possible sources of discrepancy in density profile measurements are ruled out:
* Detailed image simulations show that the sample of faint RGB/subgiant stars used by GYSB is nearly complete and, in particular, the degree of incompleteness does not vary appreciably with radius.
* If there were a bluer-inward radial gradient in the mean color of faint RGB/subgiant stars, this would result in a difference between the density profiles derived from ultraviolet versus visual-band images. However, a star-by-star match between the Calzetti et al. and GYSB data sets reveals no such ($`\mathrm{F342W}\mathrm{F555W}`$) color gradient in 47 Tuc.
More recently, DPSGB used an F336W image from the deep WFPC1 data set of Gilliland et al. (1995) to revise the cluster center derived by Calzetti et al. (1993) and measure a core radius of $`r_{\mathrm{core}}=12^{\prime \prime }\pm 2^{\prime \prime }`$, arguing that neither the Calzetti et al. (1993) composite King profile nor the traditional large core radius fit the data. This $`12^{\prime \prime }`$ core radius is within a factor of two or three of the range of theoretical estimates for the core radius of a PCC cluster with 47 Tuc’s $`r_h`$. DPSGB argued that this is suggestive that 47 Tuc may be approaching core collapse as Calzetti et al. (1993) had proposed.
## 3 Data and Photometric Techniques
There are two HST imaging data sets analyzed in this paper: (1) pre-refurbishment WFPC1 data described by Gilliland et al. (1995), consisting of a very deep F336W image ($`99\times 1000`$ s) with excellent sub-pixel dithering and shorter F439W and F555W exposures (160 s and 60 s, respectively); and (2) archival WFPC2 data from program GO-6095 (PI: S. G. Djorgovski), consisting of F218W ($`4\times 800`$ s), F439W ($`2\times 50`$ s), and F555W (7 s) images. The deep WFPC1 data set combines the advantages of ultraviolet imaging (as used by Calzetti et al. 1993) with a wider field of view (equivalent to that used by GYSB), and is significantly deeper than all previous 47 Tuc data sets. The great improvement in PSF quality in the WFPC2 data set relative to the deep WFPC1 data set more than compensates for the shorter exposure times, resulting in a higher degree of completeness and about a factor of two smaller photometric errors (Guhathakurta et al., in preparation). Moreover, the larger field of view of WFPC2 ($`150^{\prime \prime }\times 150^{\prime \prime }`$ with the PC quadrant being partially filled, extending to $`r100^{\prime \prime }`$) compared to WFPC1 ($`68^{\prime \prime }\times 68^{\prime \prime }`$, extending to $`r55^{\prime \prime }`$) makes it better for density profile measurements.
Since this work and DPSGB use the same deep WFPC1 data set, it is important to examine and compare the stellar photometry techniques used in the two studies. This paper uses standard PSF-fitting photometry techniques based on daophot ii (Stetson (1992)). These techniques are described in detail in GYSB and Yanny et al. (1994), with minor modifications to adapt them to the doubly-oversampled combined F336W WFPC1 image constructed from the sub-pixel dithered exposures (Gilliland et al. 1995; Edmonds et al. (1996)). The main steps in the technique are outlined here. A preliminary star list is derived by applying daophot’s find routine to the deep WFPC1 F336W image. For each of the four CCDs, a set of bright, unsaturated, and relatively isolated stars are used to construct a PSF template, which is allowed to vary quadratically with position. This template is then fit to all stars on the image using the preliminary star list. The neighbors of the PSF template stars are subtracted, and the process is iterated a few times. Stars too faint to be detected in the initial find run are then added based on inspection of the PSF-subtracted frames. The process of PSF template building and fitting is iterated a few more times, and the final star list is fit. The magnitudes returned by daophot’s PSF-fitting routine allstar are converted to total instrumental magnitudes by using the PSF template stars to determine a spatially-dependent aperture correction. The resulting photometry in the F336W, F439W, and F555W bands is converted to the Johnson $`UBV`$ system through an empirical fit to ground-based data from Aurière et al. (1994). A similar set of techniques is applied to the WFPC2 data set (cf. Guhathakurta et al. (1996)) to derive an independent star list for which the instrumental F218W, F439W, and F555W magnitudes are converted to Johnson $`U`$, $`B`$, and $`V`$ magnitudes, respectively.
DPSGB uses a ‘core aperture photometry’ technique defined in De Marchi et al. (1993). PSF peaks are identified on the deep WFPC1 F336W image and aperture photometry is carried out using a small aperture comparable to the size of the PSF core (hence the term ‘core aperture photometry’) and a sky annulus of roughly twice that radius. An aperture correction factor is applied to convert the resulting aperture magnitudes to total magnitudes, compensating for the portion of the PSF that lies outside the aperture and for the fact that the sky annulus includes light from the star in question. This method is hereafter referred to as aperture photometry. The sample of Calzetti et al. (1993) is analyzed using the same aperture photometry technique on the same crowded field. While § 4 makes a specific comparison of our results to DPSGB because they share a common data set, the conclusions are expected to apply to Calzetti et al.’s $`r_{\mathrm{core}}`$ analysis as well. Neighbor contamination affects aperture photometry to a much greater degree than it affects PSF-fitting photometry. Since aperture photometry is based on the total flux within the aperture, the contribution of neighboring stars is a direct source of photometric bias. Neighbors have a somewhat smaller effect on the sky measurement since the latter is estimated from the mode.
There are three sets of stellar $`U`$ magnitudes discussed in this study. The photometry from this work using the deep WFPC1 data set is referred to as $`UU_{\mathrm{WFPC1}}`$ (this paper); the photometry presented by DPSGB based on the same data set is referred to as $`U_\mathrm{D}U_{\mathrm{WFPC1}}(\mathrm{DPSGB})0.39`$ (see § 4.1), where $`U_{\mathrm{WFPC1}}(\mathrm{DPSGB})`$ is identical to the quantity $`m_{336}`$ in that paper. Photometry from the archival WFPC2 data set is referred to as $`U_{\mathrm{WFPC2}}`$. The WFPC2 PSF is much sharper than that of WFPC1, resulting in more accurate photometry; thus the WFPC2 data set is used to define the ‘true’ stellar magnitudes of the stars identified on both the WFPC1 and WFPC2 images. For each star in the WFPC1 and WFPC2 data sets, the positional offset in arcseconds, $`\mathrm{\Delta }\alpha `$ and $`\mathrm{\Delta }\delta `$ (relative to reference star ‘E’ in the 47 Tuc core; GYSB), is determined using the IRAF task metric (Gilmozzi, Ewald, & Kinney (1995)) which uses pointing information based on the HST Guide Star Catalog.
## 4 Photometric Error and its Effect on Star Counts
### 4.1 Star-by-Star Comparison
In order to compare the performance of the two photometric techniques described above, the input star lists used in the two studies are examined. If, for example, one list systematically misses stars in the crowded central region of the cluster, then this will result in differences in the measured core radii. The star finding/PSF-fitting iterations described above yield a set of 14801 stars identified on the deep WFPC1 F336W image. A star list kindly provided by Guido De Marchi gives the positions and magnitudes of 14979 stars from the same image as reduced by DPSGB. The bright stars in these two star lists match with an rms positional difference of one-third of one PC pixel, $`\sigma =0\stackrel{}{\mathrm{.}}016`$. The entire matched star list contains 12372 stars with positions agreeing to within $`3\sigma `$ (1 PC pixel) between the two input star lists. It should be noted that 141 stars appear twice on the star list from DPSGB. The majority of the $`2500`$ unmatched stars fall below the main sequence turnoff and thus are not used to calculate $`r_{\mathrm{core}}`$. An inspection of the unmatched stars suggests that many such stars from the DPSGB list may be spurious, as they tend to lie in the wings of bright stars, while unmatched stars from this work tend to be faint stars barely visible on the image. The fraction of evolved stars (defined to be those with $`U<18.11`$) that remain unmatched in each sample is quite small: 217 from this paper and 246 from DPSGB out of a total of $`>4000`$ evolved stars. Comparison of our (DPSGB’s) WFPC1 photometry to WFPC2 photometry of these unmatched stars indicates that 10% (40%) of the unmatched stars are main sequence stars scattered to brighter magnitudes as a result of photometric error. As expected, the unmatched stars are such a small fraction of the total that they have a negligible effect on the $`r_{\mathrm{core}}`$ estimate (§ 5). We conclude that the set of stars used in the $`r_{\mathrm{core}}`$ analysis is basically the same for our study and that of DPSGB; the sample differences are too small to explain the vastly different conclusions about $`r_{\mathrm{core}}`$.
Establishing a common photometric system is the next step in comparing the photometry derived from the two data reduction techniques. This paper converts instrumental magnitudes to standard Johnson $`UBV`$ magnitudes empirically as described in § 3 above. DPSGB uses $`m_{336}U_{\mathrm{WFPC1}}`$ (DPSGB) throughout, referring to instrumental magnitudes “converted to the WF/PC ground system”. A star-by-star comparison between the two data sets indicates that a constant offset of $`0.39`$ mag applied to $`U_{\mathrm{WFPC1}}`$ (DPSGB) brings the majority of stars into good agreement with our measured $`U`$ magnitude. The quantity $`U_\mathrm{D}`$ will hereafter be used in place of $`U_{\mathrm{WFPC1}}(\mathrm{DPSGB})0.39`$.
To examine the variation of photometric scatter with radius, the WFPC1 magnitude ($`U`$ or $`U_\mathrm{D}`$) of each star is compared to its ‘true’ magnitude, $`U_{\mathrm{WFPC2}}`$. The top and bottom panels of Fig. 1 plot $`(U_{\mathrm{WFPC2}}U)`$ and $`(U_{\mathrm{WFPC2}}U_\mathrm{D})`$, respectively, versus $`U_{\mathrm{WFPC2}}`$. Only stars whose ($`\mathrm{\Delta }\alpha `$, $`\mathrm{\Delta }\delta `$) positions match to within $`0\stackrel{}{\mathrm{.}}1`$ between the WFPC1 and WFPC2 data sets are used. The matched WFPC1 sample is divided at the median radius for stars with $`U<18.11`$, $`r_{\mathrm{med}}=21\stackrel{}{\mathrm{.}}54`$. The $`r<r_{\mathrm{med}}`$ (inner) subsamples are plotted in the left panels and the $`r>r_{\mathrm{med}}`$ (outer) subsamples are plotted in the right panels. The small number of points in the outer subsamples results from the limited overlap between the WFPC1 and WFPC2 fields of view away from the core of the cluster. There are systematic differences between the WFPC1 photometry and the WFPC2 photometry: $`U`$ and $`U_\mathrm{D}`$ are systematically too faint for bright red giants and too bright for main sequence turnoff stars; these are caused by errors in the conversion to Johnson $`U`$ magnitudes for the WFPC1 data set (§ 4.2).
The $`U`$ vs. $`U_{\mathrm{WFPC2}}`$ scatter is similar for inner and outer samples while the scatter increases dramatically toward smaller radii for $`U_\mathrm{D}`$ vs. $`U_{\mathrm{WFPC2}}`$. Moreover, DPSGB’s photometry of turnoff stars has a larger scatter than that of this work at all radii. While the standard deviation and mean value of $`\mathrm{\Delta }U`$ $`[(U_{\mathrm{WFPC2}}U)`$ or $`(U_{\mathrm{WFPC2}}U_\mathrm{D})]`$ are indicative of the overall rms photometric scatter and mean bias, respectively, it is instructive to examine the full distribution of $`\mathrm{\Delta }U`$ values. Stars in the interval $`17.8<U_{\mathrm{WFPC2}}<19.8`$ are used to construct cumulative distributions of $`\mathrm{\Delta }U`$. The cumulative distributions using DPSGB’s photometry show a long, asymmetric tail of large $`(U_{\mathrm{WFPC2}}U_\mathrm{D})`$, while the distributions using the photometry from this study are quite symmetrical and narrow, implying lower scatter and little or no bias. More quantitatively, the $`\mathrm{\Delta }U`$ values corresponding to the 95th percentile of the distribution are: +0.5, +0.5 (this study: inner and outer subsamples, respectively); +1.2, +0.75 (DPSGB: inner and outer subsamples, respectively).
The conclusion to be drawn from this subsection is that the star lists used in the two studies are effectively identical, but photometric accuracy is not. DPSGB’s photometry shows a systematically greater photometric error, and a bias in the sense that stars are more likely to be measured as brighter than they truly are than fainter. More importantly, this error and bias are greatest near the cluster center.
### 4.2 Color-Magnitude Diagrams
Color-magnitude diagrams (CMDs; Fig. 2) provide an alternate perspective on the nature of photometric error in the WFPC1 $`U`$-band data set, complementary to the discussion in the previous subsection. As described below, the CMDs also help to highlight a discrepancy between the WFPC1 and WFPC2 $`U`$-band photometric calibration. The CMDs combine $`B`$ and $`V`$ magnitudes, obtained via PSF-fitting, with each of the WFPC1 $`U`$ magnitudes under examination: our PSF-fitting photometry and DPSGB’s aperture photometry. The $`B`$ vs. $`UV`$ and $`B`$ vs. $`U_\mathrm{D}V`$ CMDs (left and right, respectively) are presented for stars in the inner and outer halves of the matched star sample (upper and lower, respectively). The $`B`$ vs. $`UV`$ diagrams use $`U^{\mathrm{lim}}=18.11`$, while the $`B`$ vs. $`U_\mathrm{D}V`$ diagrams use $`U_\mathrm{D}^{\mathrm{lim}}=17.71`$. Limiting magnitudes will be discussed in detail in § 4.3.
The use of the CMD as a diagnostic tool is best illustrated through stars with the most extreme photometric errors. A plot of $`UU_\mathrm{D}`$ vs. $`U_{\mathrm{WFPC2}}`$ is used to identify ‘outliers’ with $`UU_\mathrm{D}>+3\sigma `$, where $`\sigma `$ is calculated in running bins of width 0.5 mag in $`U_{\mathrm{WFPC2}}`$ (open squares in Fig. 2); the $`UU_\mathrm{D}`$ distribution is skewed towards positive values such that very few stars have $`UU_\mathrm{D}<3\sigma `$. The open squares in $`B`$ vs. $`U_\mathrm{D}V`$ diagrams generally lie on or beyond the blue fringe of the subgiant branch and RGB defined by the rest of the points. By contrast, the vast majority of the open squares in the $`B`$ vs. $`UV`$ diagrams are distributed evenly about the underlying subgiant and red giant branches, with only a few points near the red edge. This indicates that outliers in the $`UU_\mathrm{D}`$ distribution are predominantly the result of the measured $`U_\mathrm{D}`$ value being systematically too bright. As expected, the greater degree of crowding causes outliers in the inner subsample to display larger deviations in the CMD than those in the outer subsample. Moreover, the width of the RGB appears to be slightly greater in $`U_\mathrm{D}V`$ than in $`UV`$.
The solid lines show the median color of stars from the WFPC2 sample in 0.1 mag bins. The discrepancy between WFPC2 and WFPC1 photometry for both this study and DPSGB is worst for the reddest stars (bright RGB) and for the bluest stars (turnoff), and in opposite senses. This results from error in the Johnson $`U`$ calibration for the WFPC1 data set, which has been improved for the WFPC2 data set. Both data sets have been empirically transformed to the Johnson system using the ground-based data of Aurière et al. (1994). However, the larger field of view of WFPC2 includes stars farther from the cluster center and thus less affected by crowding in Aurière et al.’s data. The WFPC2 ridge line is in good agreement with Aurière et al.’s photometry, which should be accurate outside the crowded cluster core.
### 4.3 Luminosity Function
In this subsection, we examine the stellar luminosity function of 47 Tuc. This is a key element in understanding the bias in star counts in a magnitude-limited sample that results from photometric error (scatter and bias), and therefore in understanding the difference in core radius estimates between DPSGB and GYSB. The $`U`$-band stellar luminosity functions for the inner and outer samples are shown in the left and right panels of Fig. 3, respectively. Data points represent the full $`U`$-band sample from this paper, the dot-dashed line is the DPSGB $`U_\mathrm{D}`$ sample excluding the 141 duplicate stars, and the solid line is the WFPC2 sample for $`U_{\mathrm{WFPC2}}19`$, normalized to match the WFPC1 luminosity functions. The WFPC2 luminosity functions shown in the two panels are identical except for the normalization, and are based on the full star list derived from the entire WFPC2 field of view. Since this study relies only on post–main-sequence stars which are not expected to be affected by mass segregation, we have chosen to combine the entire WFPC2 sample into a single luminosity function; a full study of mass segregation will be presented in Guhathakurta et al. (in preparation).
In order to isolate a sample of post–main-sequence stars, the $`U`$ magnitude of the main sequence turnoff is determined using a sample of turnoff stars identified in a WFPC2 $`V`$ vs. $`BV`$ CMD. The mean $`U_{\mathrm{WFPC2}}`$ magnitude of these stars is adopted as the limiting magnitude for the deep WFPC1 data set analyzed in this paper: $`U^{\mathrm{lim}}=18.11`$. On the other hand, DPSGB used $`m_{336}^{\mathrm{lim}}U_{\mathrm{WFPC1}}^{\mathrm{lim}}(\mathrm{DPSGB})=18.1`$ which corresponds to $`U_\mathrm{D}^{\mathrm{lim}}=17.71`$. Vertical lines are drawn at $`U=18.11`$ and $`U=17.71`$ in Fig. 3; these correspond to the diagonal lines in Fig. 1. Note that the magnitude cut used by DSPGB falls on a steep part of the luminosity function. The next subsection quantifies the consequences of choosing a limiting magnitude on a steep part of the luminosity function.
Of the three samples, WFPC2 photometry produces the sharpest luminosity function. The luminosity function derived from this study’s WFPC1 photometry has a somewhat more gradual rise from the RGB to the main sequence turnoff than the WFPC2 luminosity function as a result of its larger photometric errors. The DPSGB luminosity functions, particularly their inner sample, have the most gradual transition between the RGB and turnoff, indicating the largest photometric errors. There are significant differences between the WFPC2 luminosity function and those derived from the deep WFPC1 data (our study and DPSGB) most notably at the bright end and near the turnoff; these are a result of systematic error in the conversion to Johnson $`U`$ for the WFPC1 data set (§ 4.2).
### 4.4 Discussion of Errors and “Limiting Magnitude Bias”
The above diagnostics (§§ 4.1–4.3) complement one another by providing somewhat different perspectives on the photometric error in the various data sets. The $`\mathrm{\Delta }U`$ plots indicate the photometric scatter as a function of radius for each of the deep WFPC1 photometry sets, but only under the assumption that $`U_{\mathrm{WFPC2}}`$ is the true magnitude of the star. While it is possible in principle to compare the three data sets to empirically determine the photometric error in each, the strongly non-gaussian error distributions make this impractical. Instead, the CMD diagnostic described in § 4.2 clearly demonstrates that there are large photometric errors associated with the $`U_\mathrm{D}`$ photometry but not with $`U`$, and this conclusion is independent of the accuracy of $`U_{\mathrm{WFPC2}}`$.
A skeptic may wonder if there are systematic errors generic to the PSF-fitting method that ‘cancel’ each other in the $`B`$ vs. $`UV`$ CMD or in the $`U`$ vs. $`U_{\mathrm{WFPC2}}`$ comparison causing the PSF-fitting magnitudes to appear more accurate than they actually are. A related question is: How accurate is $`U_{\mathrm{WFPC2}}`$? The luminosity function diagnostic is useful in this regard, despite the mismatches caused by errors in the WFPC1 Johnson $`U`$ conversion. Figure 3 clearly shows that the most prominent features (e.g., the steep rise at the subgiant branch, the peak at the main sequence turnoff) are sharpest for WFPC2 photometry, slightly smoothed out by photometric errors for our deep WFPC1 photometry, and even more smoothed out (larger photometric errors) in the case of the DPSGB data.
Photometric errors associated with the aperture photometry technique result in a significant, radially-varying bias in the star counts; we hereafter refer to this as “limiting magnitude bias”. The limiting magnitude used by DPSGB is on a steep part of the luminosity function where many stars which are slightly fainter than the limiting magnitude scatter into the sample while relatively few stars scatter out of the sample. Limiting magnitude bias is exacerbated by the asymmetry in the distribution of photometric errors (photometric bias; see § 4.1) and by the fact that the errors (scatter and bias) tend to increase towards the fainter magnitudes. Since aperture photometry errors increase with stellar crowding towards the cluster center, so does the degree of limiting magnitude bias.
The differential limiting magnitude bias between inner and outer subsamples may be quantified as follows. Since the bounding radius of the inner sample, $`r_{\mathrm{med}}`$, was chosen to be the median radius of stars in the matched sample with $`U<18.11`$, the inner to outer star count ratio is expected to be unity. It is not surprising that $`N(r<r_{\mathrm{med}})/N(r>r_{\mathrm{med}})=1.02\pm 0.03`$ for stars from this study with $`U<18.11`$. By contrast, stars with $`U_\mathrm{D}<17.71`$ from the DPSGB sample have $`N(r<r_{\mathrm{med}})/N(r>r_{\mathrm{med}})=1.42\pm 0.06`$. The latter ratio is significantly greater than unity, indicating radially-varying limiting magnitude bias in the DPSGB sample.
Alternatively, the limiting magnitude bias for each WFPC1 sample can be quantified with respect to the ‘true’ WFPC2 photometry. The diagonal lines in Fig. 1 correspond to $`U=U^{\mathrm{lim}}`$ (upper panels) and $`U_\mathrm{D}=U_\mathrm{D}^{\mathrm{lim}}`$ (lower panels). The vertical dashed lines indicate $`U_{\mathrm{WFPC2}}=U^{\mathrm{lim}}`$ (upper panels) and $`U_{\mathrm{WFPC2}}=U_\mathrm{D}^{\mathrm{lim}}`$ (lower panels). The ratios $`N(U<U^{\mathrm{lim}})/N(U_{\mathrm{WFPC2}}<U^{\mathrm{lim}})`$ are $`1.54\pm 0.06`$ and $`1.34\pm 0.15`$ for inner and outer subsamples, respectively, while $`N(U_\mathrm{D}<U_\mathrm{D}^{\mathrm{lim}})/N(U_{\mathrm{WFPC2}}<U_\mathrm{D}^{\mathrm{lim}})`$ are $`2.32\pm 0.14`$ and $`1.09\pm 0.22`$ for inner and outer subsamples, respectively. These ratios are greater than unity because errors in the WFPC1 Johnson $`U`$ conversion cause the main sequence magnitudes to be systematically too bright (§ 4.2; Fig. 1). Each ratio indicates the number of stars selected in a sample relative to the corresponding true (WFPC2-based) number of stars satisfying the selection criterion. For the photometry from this work, the excess relative to unity is similar in the inner and outer subsamples, so the derived core radius should not be greatly affected. DPSGB’s photometry results in a large and significant excess of stars selected in the inner subsample, but not in the outer subsample. This radial variation has a substantial impact on the core radius measurement (§ 5). The large difference between the inner and outer DPSGB ratios is a direct result of the increased scatter and bias in their photometry at small radii.
## 5 Core Radius Calculations
The effect of limiting magnitude bias on the measurement of the core radius is explored in this section. A one-parameter surface density profile of the form:
$$\sigma (r)=\frac{\sigma _0}{[1+(r/r_{\mathrm{core}})^2]}$$
(1)
is used to make maximum likelihood fits to various magnitude-limited samples of evolved stars to determine $`r_{\mathrm{core}}`$, while the normalization constant $`\sigma _0`$ is constrained by the total number of stars in each sample. Calculations are performed using both GYSB’s estimate of the cluster center:
$$\alpha {}_{\mathrm{J2000}}{}^{}=00^\mathrm{h}24^\mathrm{m}05^\mathrm{s}.87$$
$$\delta {}_{\mathrm{J2000}}{}^{}=72^{}04^{}57\stackrel{}{\mathrm{.}}8$$
(2)
and DPSGB’s center estimate:
$$\alpha {}_{\mathrm{J2000}}{}^{}=00^\mathrm{h}24^\mathrm{m}05^\mathrm{s}.29$$
$$\delta {}_{\mathrm{J2000}}{}^{}=72^{}04^{}56\stackrel{}{\mathrm{.}}3.$$
(3)
The two center estimates agree to within the quoted uncertainty of $`1\stackrel{}{\mathrm{.}}0`$$`1\stackrel{}{\mathrm{.}}3`$ in each estimate. Unless otherwise stated, all core radius estimates in this paper are based on the GYSB center.
Table 5 contains a summary of the core radius measurements. The three letters identifying each calculation indicate the source of: (1) the star list (D for DPSGB, H for this work, M for the matched sample); (2) the photometry (D or H); and (3) the limiting magnitude (D or H). For direct comparison with DPSGB’s $`r_{\mathrm{core}}`$ measurement, the 141 duplicate stars are included in their star list. The result of the maximum likelihood fit (case DDD) is $`r_{\mathrm{core}}=13\stackrel{}{\mathrm{.}}7\pm 1\stackrel{}{\mathrm{.}}8`$, in agreement with the DPSGB value of $`r_{\mathrm{core}}=12\stackrel{}{\mathrm{.}}2\pm 2\stackrel{}{\mathrm{.}}1`$ derived using radial binning and a least-squares fit. A maximum likelihood fit to the evolved star sample from this paper (case HHH) gives $`r_{\mathrm{core}}=21\stackrel{}{\mathrm{.}}8\pm 2\stackrel{}{\mathrm{.}}0`$. Applying a limiting magnitude of $`U_\mathrm{D}^{\mathrm{lim}}=18.11`$ instead of $`U_\mathrm{D}^{\mathrm{lim}}=17.71`$ to DPSGB’s star list and photometry results in a core radius of $`17\stackrel{}{\mathrm{.}}4\pm 1\stackrel{}{\mathrm{.}}5`$ (case DDH). The change in derived core radius from $`13\stackrel{}{\mathrm{.}}7`$ to $`17\stackrel{}{\mathrm{.}}4`$ is a direct result of the change in limiting magnitude and is independent of whether the 141 duplicate stars are included. Using the star list and photometry from this study but choosing DPSGB’s limiting magnitude $`U^{\mathrm{lim}}=17.71`$ gives $`r_{\mathrm{core}}=21\stackrel{}{\mathrm{.}}0\pm 3\stackrel{}{\mathrm{.}}3`$ (case HHD). Thus, the large scatter/bias in DPSGB’s photometry and their increase towards the crowded cluster center, combined with a poor choice of limiting magnitude, are responsible for the spuriously low core radius estimate.
Table 5 also shows that the measured $`r_{\mathrm{core}}`$ is independent of whether the star list from this study, DPSGB’s star list, or the matched star list is used. The dependence of the derived $`r_{\mathrm{core}}`$ on limiting magnitude and photometric technique is roughly the same for all three lists. Likewise the choice of cluster center has little effect on the derived $`r_{\mathrm{core}}`$: although core radii based on the DPSGB center tend to be $`0\stackrel{}{\mathrm{.}}5`$ smaller than those based on GYSB’s center, this difference is less than the uncertainty in an individual measurement.
The WFPC2 data set provides a consistency check on the core radius of 47 Tuc. For all stars derived from the archival WFPC2 data set with $`U_{\mathrm{WFPC2}}<18.11`$, the same maximum likelihood algorithm used above yields $`r_{\mathrm{core}}=23\stackrel{}{\mathrm{.}}1\pm 1\stackrel{}{\mathrm{.}}7`$. An independent sample of 7044 turnoff stars with $`17V<18`$ drawn from the WFPC2 $`V`$ vs. $`BV`$ CMD (avoiding potential blend artifacts that may have scattered to brighter magnitudes) results in $`r_{\mathrm{core}}=23\stackrel{}{\mathrm{.}}3\pm 1\stackrel{}{\mathrm{.}}2`$. While these turnoff stars are fainter on average than the evolved stars used in earlier density profile studies, they are expected to have roughly the same mass as red giants so that mass segregation effects are unlikely to be important. A third calculation truncates the $`U_{\mathrm{WFPC2}}<18.11`$ sample at $`r85^{\prime \prime }`$ as a precaution against inter-CCD edge effects caused by vignetting near the ridges of the pyramid mirror. A core radius of $`r_{\mathrm{core}}=24\stackrel{}{\mathrm{.}}0\pm 1\stackrel{}{\mathrm{.}}9`$ (case WFPC2 in Table 5) is derived from this truncated sample. It is reassuring that these WFPC2 $`r_{\mathrm{core}}`$ measurements are consistent with GYSB, earlier ground-based measurements, and the HHH and MHH calculations above. Figure 4 shows the radial surface density profile of the WFPC2 truncated sample. Also shown are an $`r_{\mathrm{core}}=24\stackrel{}{\mathrm{.}}0`$ profile (solid line) and an $`r_{\mathrm{core}}=12\stackrel{}{\mathrm{.}}2`$ profile (dashed line); the latter is clearly inconsistent with the data. These profiles have been normalized to the total observed number of stars in the plot. Figure 5 shows the cumulative radial distributions of the MHH sample (left panel) and the truncated WFPC2 sample (right panel). Also shown are profiles for the best fitting core radii (dotted lines) and for core radii that differ from the best fit value by $`\pm 1\sigma `$ (dashed lines).
As an alternative to the maximum likelihood technique, we conduct Kolmogorov-Smirnov tests between the star count samples and the fitting function in Eq. (1). The MHH sample of 4504 stars (Table 5) and the truncated WFPC2 sample of 3963 stars are used to construct cumulative radial distributions (Fig. 5). One-sided Kolmogorov-Smirnov tests indicate that the MHH data differ from the $`r_{\mathrm{core}}=22^{\prime \prime }`$ fitting function by an amount that would be exceeded by chance 26% of the time. Similarly, the truncated WFPC2 data differ from an $`r_{\mathrm{core}}=24^{\prime \prime }`$ profile by an amount that would be exceeded by chance 33% of the time. Thus, the functional form adopted here provides an adequate fit to the star count data.
As a final check on the $`r_{\mathrm{core}}`$ calculations, one can test whether the degree of limiting magnitude bias (§ 4) for a given sample is sufficient to explain the core radius derived for it. Using the photometry of DPSGB and their sample of stars with $`U_\mathrm{D}<U_\mathrm{D}^{\mathrm{lim}}`$ (case DDD in Table 5), the inner to outer star count ratio is $`N(r<r_{\mathrm{med}})/N(r>r_{\mathrm{med}})=1.42\pm 0.06`$. Integrating a $`14^{\prime \prime }`$ profile (comparable to the maximum likelihood fit for this sample) over the area of the WFPC1 field of view predicts an inner to outer ratio of 1.49 consistent with the directly measured star count ratio. On the other hand, integration of the $`22^{\prime \prime }`$ profile that best fits the star list and PSF-fitting photometry used in this paper (case HHH in Table 5) yields an inner to outer ratio of 1.05; this is comparable to the directly measured star count ratio of $`1.02\pm 0.03`$.
## 6 Conclusions
This paper presents estimates of the density profile of the globular cluster 47 Tuc based on three samples of stars (star list and photometry) derived independently from Hubble Space Telescope WFPC1 and WFPC2 images. Apparent discrepancies amongst the core radius measurements published by Guhathakurta et al. (1992), De Marchi et al. (1996), and Calzetti et al. (1993) are investigated. Our conclusion is that there is severe, radially-varying bias in the magnitude-limited star counts used by De Marchi et al. and Calzetti et al., and this causes their core radius estimates to be spuriously low ($`r_{\mathrm{core}}14^{\prime \prime }`$) relative to other determinations ($`r_{\mathrm{core}}23^{\prime \prime }`$). This “limiting magnitude bias” is a result of large photometric scatter/bias associated with the application of their aperture photometry method to the crowded central regions of the cluster, coupled with a choice of limiting magnitude near the steep part of the stellar luminosity function. In general, such a choice of limiting magnitude is dangerous; even with symmetric errors the resulting sample will be contaminated by large numbers of stars just fainter than the cutoff. Any radial variation in the magnitude of the errors will cause the degree of this contamination to vary, resulting in an incorrect determination of the radial density profile.
Combining De Marchi et al.’s photometry with a limiting magnitude near the main sequence turnoff at the peak of the luminosity function reduces, but does not eliminate, the discrepancy ($`r_{\mathrm{core}}18^{\prime \prime }`$); the radial variations in DPSGB’s photometry (larger photometric scatter/bias at small radii) have a significant effect on the derived core radius even with an optimal choice of limiting magnitude. A more accurate PSF-fitting method is used in this paper to indepedently derive two sets of stellar photometry, one from the deep WFPC1 data analyzed by De Marchi et al. and the other from an archival WFPC2 data set. The core radii derived using these two photometry sets are independent of the choice of limiting magnitude and star list, and are consistent with each other and with previous ground-based and HST work: $`r_{\mathrm{core}}23^{\prime \prime }`$ (cf. Harris & Racine (1979); Djorgovski & King (1984); Guhathakurta et al. 1992).
The best fit core radius for the surface density distribution of evolved stars in 47 Tuc is about 15% of the cluster half-mass radius ($`r_\mathrm{h}=174^{\prime \prime }`$). This is significantly larger than the range of $`(r_{\mathrm{core}}/r_\mathrm{h})`$ values found in numerical simulations of post–core-collapse clusters: 0.01–0.04 (Cohn (1980); Goodman (1987); Gao et al. (1991)). It should be noted however that the surface brightness profile is not a perfect discriminant between a relaxed and post–core-collapse cluster, and it is advisable to combine it with velocity dispersion data (Gebhardt & Fischer 1995).
###### Acknowledgements.
We are grateful to Guido De Marchi for making stellar photometry tables available to us in electronic form, and to Randi Cohen for help in the early phase of this project. We would like to thank Fernando Camilo, Paulo Freire, Karl Gebhardt, and Fred Rasio for helpful discussions, and the referees, especially Tad Pryor, for several useful suggestions.
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# References
CHARGE ASYMMETRY IN 1-1000 GEV ELECTROMAGNETIC SHOWERS AND POSSIBILITY OF ITS MEASUREMENT
R.O.Avakian, K.A.Ispirian<sup>∗)</sup>, R.K.Ispirian, V.G.Khachatryan, P.Sona <sup>a)</sup> and E.Uggerhoj <sup>b)</sup>
Yerevan Physics Institute, Brothers Alikhanian 2, Yerevan, 375036, Armenia
a) INFN and University of Firenze, Firenze, Italy
b) Institute for Storage Ring Facilities, University of Aarhus, Denmark
Abstract
For the high energy electromagnetic showers the thickness dependence of a) the development of electron and positron components, b) the difference between the secondary electron and positron numbers, c) the charge asymmetry of high energy electromagnetic showers, as well as d) the spectral distributions of the components at the shower maxima for various energies of primary particle energies,1 - 1000 GeV were investigated employing GEANT Monte Carlo simulation package. Using these simulation results it is discussed the possibility of observation and study of the charge asymmetry with the help of a magnetic spectrometer which is important for the current and future experiments on the detection of radiowaves produced by high energy neutrinos.
PACS: 13.40.-f; 96.40.Pq
Keywords: Electromagnetic Showers; Cherenkov Radiation; Magnetic Spectrometers
I.Introduction
At present all the elementary electromagnetic processes taking place when high energy electrons and photons pass through the matter are well known. Therefore, the formulation of the correct theory for high energy electromagnetic showers (EMS) is possible in principle.However, due to mathematical difficulties the construction of the EMS theory
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\*) e-mail: karo@robert.yerphi.am
is realized in various approximations, and almost always the study of the EMS theory and comparison with the experimental data are carried out with the help of Monte Carlo simulalations using the existing codes as EGS, GEANT and others developed at SLAC, CERN and other laboratories. In all the available EMS codes the ”fate” of electrons and photons is followed down to certain lowest energies, cut energies, in order to escape very large program complication and computation times. Usually these cut energies are less than 1 MeV. In some cases as in the study of biological processes and some processes considered below when high density energy depositions, say, by $`\delta `$\- or Compton electrons are essential one should lower these threshold energies. Nevertheless, there is a satisfactory agreement between the experimental observations and Monte Carlo simulation results on EMS, and one can use the the laters to study some processes which are not observed and studied yet.
In 1961 G.Askarian predicted an excess of electrons over positrons in high energy EMS due to positron annihilation in flight, Compton and $`\delta `$\- electrons and estimated the intensity of coherent Cherenkov and transition radiation radiowaves produced by this moving negative charge excess. The estimates carried out in neglecting the contribution of Compton-, $`\delta `$\- and not mentioned in photo-electrons show that the number of electrons in EMS can exceed the number of positrons by more than 10 $`\%`$ at energies higher than hundreds MeV. The more accurate calculations (see and references therein) of such electron surplus with the help Monte Carlo simulations confirm the predictions of . In the latest calculations it has been shown that the following processes give main contribution in the production of the EMS charge excess: Compton scattering (50-60 $`\%`$) , Bhaba scattering (30-35 $`\%`$), positron annihilation in flight (5-20 $`\%`$), and their contribution depends weakly on the primary particle energy and relatively strongly on the component energy. Due to the low cross section the number of MeV photons is larger than the combined number of electrons and positrons in EMS, and Compton effect gives the largest part of the excess in MeV energy region.However, at present there is no direct experimental results on the EMS charge excess, while the existing indirect data (see below) are ambigeous and need correct interpretation.
In the work it has been also estimated and shown that the intensity of the coherent Cherenkov radiation radiowaves produced by this moving negative charge excess is sufficient to be used for detection of high energy EMS on the earth and moon. Following it was suggested to use this radiowave production for the undergroud detection of high energy neutrinos in the salt mines, while the mechanisms of radiowave production has been considered in details by the authors of the work assuming various time and space distributions of the charge excess. After these and later theoretical and experimental investigations carried out in sixties devoted to EMS charge excess and radiowave production (see, ) many works have been published by the cosmic ray physicists because the method promised to be very convinient for the very high energy neutrino astrophysics and neutrino oscillation problems.
Despite the achievements in this field after more than 35 years of theoretical and experimental investigations and many interesting projects under construction (see and Proceedings of last International Cosmic Ray Conferences), the technique of detection of EMS with the help on radio antennas has not yet been proven and difficulties are anticipated . In this connection it seems reasonable to study experimentally the various characteristics of the EMS charge excess and of the coherent Cherenkov as well as transition radiation produced by the available high energy electron beams at various accelerators.
When this paper was ready for submitting an electronic preprint has appeared in which the authors in addition to the existing experimental studies devoted to the far infrared and submillimeter coherent Cherenkov and transition radiation have investigated the polarization, angular, coherence and other properties of the same radiations in GHz radio region using 15.2 MeV electron bunches. The authors conclude that it is necessary to carry out more accurate measurement for various applications.
Taking into account the above said, the actuality of the problem and the available contemporary computational possibilities in this work we study the processes connected with the charge excess at primary energies 0.5 - 256 GeV with the help of the GEANT code package. It is shown the possibility of the observation and experimental study of this processes at YerPhI, SPS, CERN, and FERMILAB.
II.Results of Simulations
We have chosen GEANT to carry out the necessary Monte Carlo simulations on EMS for two reasons. First, long term practice indicates that GEANT handles in a proper way. Second, GEANT is disigned to to simulate the geometry of the experimental setup, which is essential fo our purpose. The agreement between our calculations and published results, in particular, on the depth dependence at higher energies, witnesses the correctness of our calculations. Calculations have been performed for various kinetic energy cuts for electrons and photons, $`T_{cut}=T_{cut}^e=E_{cut}^\gamma `$ from $`T_{cut}=50keV`$ up to $`T=12MeV`$ and when the primary particles were photons (the results for electrons do not differ significantly from those for photons) with total number $`N_\gamma `$ and various primary photon energies $`E_\gamma `$ from 0.5 GeV up to 1000 GeV. Each element of the calculation array with fixed $`E_\gamma `$ and $`E_{cut}`$ containes information on a) the dependence of the electron and positron numbers $`N_{e^{},e^+}`$ upon the depth t in radiation length units; b) the dependence of the excess $`\nu =N_e^{}N_{e^+})`$ on t; c) the dependence of the charge assymetry $`A=(N_e^{}N_{e^+})/(N_e^{}+N_{e^+})`$ on t and d) the energy spectrum of the electrons at the depth where the maximum of the charge excess for the given parameters takes place, $`t=t_{max}`$. All the calculations presented in this work have been carried out for BGO because it is a diamagnetic insulator, has a small radiation length unit (useful properties for radiowave detection), has sufficient scintillation yield which can be useful in some cases and is available. Fig.1 shows the information of one array element when $`E_\gamma =128`$ GeV and $`T_{cut}=0.4`$ MeV. As it is seen from Fig.1 a the showering behaviors for electrons and positrons are similar, but they differ significantly in magnitudes. The behavior of the excess t-dependence (see Fig.1b) reminds the usual behavior of shower curves with tails of the form $`exp(\alpha t)`$ and its more intense part around the maxima can be approximated roughly by the symmetric function $`exp(\omega _0\tau ^2)`$ where $`\tau =tt_{max}`$ and $`\alpha `$ and $`\omega _0`$ are constants as it is suggested in to calculate the radiation intensity. As it follows from Fig.1c for the given $`E_{cut}=0.4`$ MeV the assymetry exceeds the value given in because of the contribution of low energy electrons produced due to Compton effect. However, as it will be shown below for energies of the electron component higher then few MeV the asymmetry becomes less than it is predicted in . The results given in Fig.1d show that indeed one can measure the assymetry, and such measurement is easier for lower energies of electrons and positrons. As it is seen from Fig.2 the charge asymmetry virtually disappeares above 20 MeV.
Using many such simulation results as ones presented in Fig.1 one can reveal the characteristic properties of the EMS charge asymmetry necessary for the future employment. In Fig.3 a and b it is given the dependence of the charge excess on $`E_\gamma `$ (for fixed $`E_{cut}`$) and $`E_{cut}`$ (for fixed $`E_\gamma `$), respectively. $`\nu `$ increases almost linearly with the increase of $`E_\gamma `$ and decreases with the decrease of $`E_{cut}`$. In Fig.4 a and b it is given given the dependence of the charge asymmetry on $`E_\gamma `$ (for fixed $`E_{cut}`$) and $`E_{cut}`$ (for fixed $`E_\gamma `$), respectively. It is seen that $`A`$ almost does not depend on $`E_\gamma `$, decreases with the decrease of $`E_{cut}`$ and almost vanishes when $`E_{cut}>10`$ MeV. Therefore, it is advantageous to study the charge asymmetry at possible higher primary particle energies and lower cut energies.
As expected the calculations show that at the shower maxima the contributions from various processes resulting in charge excess depend on the component energy, and in the energy region below few MeV where the number of the electrons and the charge asymmetry are larger the proportion of the contribution fron various processes coincides with that given in for higher energies, $`E1TeV`$.
III. Asymmetry Meassurement Using Magnetic Spectrometers
Various characteristics of EMS have been investigated experimentally for a wide energy region of electrons and photons from 50 MeV up to few TeV with accelerator and cosmic ray particles using various methods and detectors. The authors of the work used streamer chambers. In all these works the measurements have been carried out for secondary particle energies not less than 1 MeV. This is not because the corresponding Monte Carlo calculations of that time were available for $`E_{cut}>1`$ MeV, but because the applied methods did not allow to decrease further the energies of the detected secondary particles because of larger energy measurement errors due to multiple scattering.
The use of streamer chambers in magnetic field with insulator layers (BGO) in which the EMS are developed seems more suitable for the EMS charge excess investigations because they give the possibility to carry out the measurements at various depths simultaneously. With low Z gas filling and magnetic fields B = 0.03 T the expected accuracy for the energy measurements are about 16, 14, 7 and less than 5 $`\%`$ for electron energies 0.075, 0.15, 0.3 and 1.0 MeV, respectively. The use of streamer as well as time projection chambers is connected with technical difficulties, and it will be much easier to perform such studies with the help of low energy magnetic pair spectrometers.
It will be convinient to carry out such an experiment with the arrangement NA59, SPS, CERN (see Fig.6 of ) proposed for other purpose and which will be ready in spring 1999. Since as it has been mentioned above the charge excess characteristics do not depend whether the primary particles are electron or photons, the 150-180 GeV electrons of the H2 beamline or the gamma quanta produced by these electrons must be focused (The beam angular and energy spread are not important, while its cross section radius must be decreased to 1-2 cm, since the Mollier radius for BGO is 2.4 cm) on 5-15 cm thick BGO slabs replacing the berillium target in the experiment NA59, and the pair spectrometer magnet with B.l = 0.52 Tm (l is the length of the magnetic field) must be replaced with a weak magnet with B.dl = 0.0033 Tm. Such weak fields and few meter distance provide the deflection of particles with energies higher than few MeV under angles greater than the angles under which the particles leave the BGO radiator and detect them at transversal distances larger than few Moller radius. The higher energy components do not touch the sensitive parts of the detectors downstream the magnet. The energy of the photons in the region 96-144 GeV is determined by the tagging system.The energy of the negative and positive shower particles coming out from BGO is measured with the help of two or three drift chambers. Since there are no polarization measurements the charge excess measurements at some depth will be much easier than the shower measurements on the arrangement NA43 using polarized photon beams. Again the multiple scattering in various thin windows alowes to determine the charge for secondary particle energies higher than few MeV. Since the primary electron beam intensity is $`4.10^5min^1`$ and the expected number of photons with energies 96-144 GeV will be only less by one order, the measurement time at one depth estimated with the above given curves is about 1 and 10 hours in the cases of primary electrons and photons, respectively .
Though the number of shower components is much lower at GeV energies (see Fig. 2a), nevertheless, EMS charge excess measurements are also possible at such energies using the ejected electron and photon beams, say, at Yerevan Synchrotron because of their higher intensity, about $`10^9`$ electrons or photons per second.
IV.Discussion
In this work it is reported the results of the Monte Carlo simulations on the negative charge excess and differential energy spectra of secondary particles in EMS taking into account all the processes in the primary particle energy interval 1-256 GeV. For the energy region of the shower components below few MeV (above ten MeVs) the presented results predict much larger (smaller) excess than the estimates . Nevertheless, as it has been shown in this work this charge excess can be measured with the low intensity, but high energy ($``$ 100 GeV ) electron and photon beams at SPS, CERN and Fermilab or using the high intensity but relatively low energy ( $``$ 1 GeV) beams at YerPhI. At present since the expected intensities of the atmospheric EMS coherent radiation of very high energy particles in the radio diapason is higher, than the intensity of the EMS transition radiation in the clouds and the sensitivity threshold of antennas, the formers are detected unambiguously in coincidence with other extended shower detectors. However many problems concerning the correct mechanisms of radio wave production, spectral and angular distribution etc. remain unsolved before wide application in very high energy neutrino astrophysics. The results of the excess measurements proposed in this work can shed light on the problems of EMS radiowave detection in dense (ice or salt) and air media.
Figure Captions
Fig.1. Characteristic charge excess dependences in BGO for primary photons when $`E_\gamma =128`$ GeV and $`E_{cut}=0.4`$ MeV. a) Shower curves separately for secondary electrons (solid histograms) and positrons (dashed histogram); b) Dependence of the excess $`\nu =N_e^{}N_{e^+}`$ upon t (in radiation length units): c) Dependence of the asymmetry $`A=(N_e^{}N_{e^+})/(N_e^{}+N_{e^+})`$ upon depth, and d) Differential spectrum of the electrons at the maximum of the charge excess (100 events are simulated).
Fig.2. The dependence of the charge asymmetry on the energy of electrons and positrons at the shower maximum in the energy intervals a)1-10 MeV and b)10- 50 MeV.
Fig.3. The dependence of the charge excess on a) $`E_{cut}`$ for the fixed $`E_\gamma =128`$ GeV and on b) $`E_\gamma `$ for the fixed $`E_{cut}=0.4`$ MeV.
Fig.4. The dependence of the charge asymmetry on a) $`E_{cut}`$ for the fixed $`E_\gamma =128`$ GeV and on b) $`E_\gamma `$ for the fixed $`E_{cut}=0.4`$ MeV.
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# Why odd-space and odd-time dimensions in even-dimensional spaces?*footnote **footnote *NBI-HE-00-26
## I Acknowledgement.
This work was supported by the Ministry of Science and Technology of Slovenia as well as by funds CHRX - CT - 94 - 0621, INTAS 93 - 3316, INTAS - RFBR 95 \- 0567, SCI-0430-C (TSTS). One of us (H.B.N.) wishes to thank many people for numerous discussions on the related subject of why we have 3+1 dimensions, and in most recently especially Bo Sture Skagerstam, S. E. Rugh, Y. Takanishi, D. Bennett, C. D. Froggatt and N. Stillits; N.M.B. would like to thank A. Borštnik and B. Gornik for discussions helpful for this letter.
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# FUSE Observations of the Galactic and Intergalactic Medium Towards H1821+643
## 1 Introduction
In this paper, we present intermediate resolution ($`20`$ km s<sup>-1</sup>) spectra of the quasar H1821+643 ($`V=14.2,z=0.297`$) which is visually the brightest known object in the sky at $`z>0.2`$. H1821+643 is an X-ray selected, radio-quiet quasar (Pravdo & Marshall, 1984), and resides in a cD galaxy at the center of a rich cluster of galaxies (Schneider et al., 1992). The optical properties of the QSO have been reported by Hutchings & Neff (1991) and Kolman et al. (1993). The latter authors find evidence for a strong optical/UV bump and a soft x-ray excess. The cluster of galaxies hosting the QSO is one of the most luminous X-ray clusters known (Hall, Ellingson & Green, 1997). Spectroscopic studies of H1821+643 have been carried out in the UV at low resolution (Bahcall et al., 1992; Kolman et al., 1993; Tripp, Lu & Savage, 1998) and intermediate resolution (Savage, Sembach & Lu, 1995; Savage, Tripp & Lu, 1998; Penton, Stocke & Shull, 2000), in order to study the Galactic and intergalactic medium along the line of sight.
## 2 Observations
H1821+643 was observed by FUSE on Oct 10 and 13, 1999 for a total integration time of 48.8 ksec. Moos et al. (2000) and Sahnow et al. (2000) provide an overview of the FUSE instrument and its performance. Data in the wavelength interval $`9901185`$ Å covered by both LiF channels were obtained. Processing steps included orbital doppler compensation, background subtraction, flux and wavelength calibration. No flat-fielding of the data was performed. The dispersion solution was derived from pre-launch calibration spectra, and adjusted based on comparison of FUSE absorption line velocities to velocities of appropriate absorption lines in the STIS/HST echelle spectrum of H1821+643 obtained by Tripp, Savage & Jenkins (2000). In general, the velocities are accurate to $`5`$ km s<sup>-1</sup>. The measured spectral resolution was $`2025`$ km s<sup>-1</sup> depending on wavelength. The detector background and scattered light in the instrument are extremely low, leading to accurate knowledge of the residual intensities in absorption lines. Equivalent widths and column densities of selected absorption lines were measured using the techniques described by Sembach & Savage (1992) and are reported in Table 1. Spectra shown in Figures 1-3 have been smoothed by 5 pixels (0.03Å).
## 3 Galactic Absorption
The H1821+643 sight line $`(l,b)=(94^{},27^{})`$ passes over 3 spiral arms seen in H i emission in a survey above the galactic plane by Kepner (1970): (1) the intermediate arm ($`v_{LSR}42`$ km s<sup>-1</sup>), the Perseus arm ($`v_{LSR}75`$ km s<sup>-1</sup>) and the outer arm ($`v_{LSR}93,120`$ km s<sup>-1</sup>). For the observations reported here $`v_{LSR}=v_{hel}+16`$ km s<sup>-1</sup>. The sight line is also in a direction where the warp of the outer galaxy extends to large Galactic latitudes. The interstellar spectrum of H1821+643 is dominated by lines of H<sub>2</sub> and Fe ii in the FUSE bandpass. Also present are Ar i, C ii, Si ii, P ii, and O vi. Strong O i and N i lines were also detected in spectra obtained while the satellite was in the earth’s shadow, where the terrestrial day glow emission lines in these species are absent. The $`10301040`$Å region of the spectrum displaying O vi $`\lambda \lambda `$1032,1038, C ii $`\lambda `$1036 and several H<sub>2</sub> lines is shown in Figure 1.
The Fe ii lines are strong and have negative velocity wings extending out to as much as $`170`$ km s<sup>-1</sup> for Fe ii $`\lambda `$1145. These negative velocity components are formed above the outer arm, and are also seen in Mg ii and Si ii lines in GHRS data (Savage, Sembach & Lu, 1995). The H<sub>2</sub> lines are local and are not seen at large negative velocities.
It is interesting to compare the strengths of the Ar i $`\lambda \lambda `$1048,1066 lines to those in the N i multiplet at 1134Å. In a fully neutral medium, these lines should have similar equivalent widths since they have nearly similar values of $`P\mathrm{log}(A)+\mathrm{log}(f\lambda )`$ where $`A=`$ cosmic abundance, i.e., $`P8.59.0`$. However, the Ar i lines are noticeably weaker than the strongly saturated N i lines in the spectra of H1821+643. Sofia & Jenkins (1998) have pointed out that in diffuse clouds Ar i is unlikely to be depleted onto dust grains. Nevertheless, they argue that in regions that are partially ionized by EUV radiation, Ar i may appear to be deficient relative to H i (or N i) because it has a substantially larger ionization cross section and thus might be more ionized. Note that Ar i is only clearly detected in the intermediate arm. It is not detected in the Perseus or outer arms, probably because the Ar is more highly ionized at greater distances from the Galactic plane.
The O vi absorption at $`120`$ to $`150`$ km s<sup>-1</sup> is associated with the outer arm and distant warp of our Galaxy. If we assume that the halo corotates with the underlying disk, the implied Galactocentric distance of the absorbing gas is $`2550`$ kpc, and the distance above the plane is $`1020`$ kpc. Independent evidence for the existence of O vi absorption in the outer halo of our galaxy is supplied by the absence of O vi absorption at velocities exceeding $`70`$ km s<sup>-1</sup> in the FUSE spectrum of K1-16, which is only $`85^{\prime \prime }`$ away from H1821+643 on the sky, and is at a distance of 1.6 kpc (Kruk et al., 2000). Consequently, the high negative-velocity O vi absorption seen in H1821+643 must be formed in the Galactic halo beyond K1-16.
A high negative-velocity ($`215`$ km s<sup>-1</sup>) component of O vi $`\lambda `$1031.93 is coincident with a low velocity H<sub>2</sub> line from the $`(60)P(3)`$ rotational level at 1031.19Å. We have modeled the $`J=3`$ H<sub>2</sub> lines in the spectrum, and find a good fit for a total column density of N(H<sub>2</sub>) $`2\times 10^{15}`$ cm<sup>-2</sup> and a cloud excitation temperature of $`T_{ex}500`$ K. Based on this model, the expected $`H_2`$ absorption at 1031.19Å cannot account for the depth or width of the observed line, indicating significant absorption by O vi $`\lambda `$1032 at $`215`$ km s<sup>-1</sup>. This confirms the tentative identification of this component in the C iv $`\lambda `$1549 line made with GHRS/HST by Savage, Sembach & Lu (1995). This detection is quite interesting, since the line velocity corresponds to the limiting velocity of $`v_{LSR}190\pm 20`$ km s<sup>-1</sup> for Milky Way gas assuming corotation and a flat rotation curve for the outer galaxy. Collisional ionization seems the most likely source of O vi ionization in this component. Photoionization would require a very hard radiation field ($`E>114`$ eV) to ionize O v to O vi, and the large photoionization parameter, $`U`$, requires a very low value of $`n_H`$ and an extremely long pathlength, $`l=N_H/n_H>100`$ kpc, to produce the absorption (see Sembach et al. (2000)).
## 4 Intervening Absorption
Tripp, Lu & Savage (1998) have detected strong (W$`{}_{\lambda }{}^{}>200`$ mÅ) Ly$`\alpha `$ towards H1821+643 at redshifts of 0.1213, 0.1476, 0.1699, 0.2132 and 0.2249. Savage, Sembach & Lu (1995) also reported intervening Ly$`\alpha `$ absorption at $`z=0.02454`$. We have detected absorption from four of these absorbers in the FUSE spectrum, and the column densities are reported in Table 1. The Lyman series lines detected with FUSE in these systems provide important constraints on N(H i) since the stronger lines in the HST bandpass are saturated. The reader is also referred to Shull et al. (2000) for discussion of intervening Ly$`\beta `$ absorbers. One must be aware of the potential for ejected material from the QSO to masquerade as intervening gas, as shown for C iv absorption line systems by Richards et al. (1999). We believe that there is good evidence that the absorbers in H1821+643 are truly intervening, based on the proximity of intervening galaxies to the line of sight (Tripp, Lu & Savage, 1998; Bowen & Tripp, 2000).
Tripp, Lu & Savage (1998) report a strong, blended line of Ly$`\alpha `$ at $`z=0.12123,0.12157`$. In addition to detecting the the Ly$`\beta `$ line in this system, we report the detection of a line at 1157.17Å which we have identified as O vi $`\lambda `$1032 at $`z=0.12137`$ (see Figure 2). The line is clearly seen in both the LiF1 and LiF2 spectra. The O vi $`\lambda `$1038 line is blended with the Ly$`\delta `$ line in the $`z=0.225`$ absorber and could not be measured. This is the sixth intervening O vi absorption system observed towards H1821+643, and provides further evidence that low-$`z`$ O vi systems contain a large fraction of the baryons at the present epoch (Tripp, Savage & Jenkins, 2000; Tripp & Savage, 2000). These observations are in accord with cosmological simulations by Cen & Ostriker (1999) that predict a substantial fraction of present-day baryons are in a shock-heated phase at $`10^510^7`$K.
We have identified an observed line at 1143.6Å as C iii $`\lambda `$977 at $`z=0.1705`$. This identification was made possible by comparison with the wavelength of a Ly$`\alpha `$ absorber recently observed by Tripp, Savage & Jenkins (2000) with STIS/HST. The Ly$`\gamma `$ line for this absorber is not conclusively detected in the FUSE spectrum, although a weak, rather broad feature is at the predicted wavelength. The Ly$`\beta `$ line in this absorber is redshifted into the strong Galactic N i triplet at 1200Å.
Intervening absorption at $`z=0.225`$ has been detected in the Ly$`\alpha `$ and Ly$`\beta `$ lines and in O vi with HST (Savage, Tripp & Lu, 1998). Recently, Tripp, Savage & Jenkins (2000) have also clearly detected Si iii in this system, and the component structure establishes that this is a multiphase absorber. However, this absorber is not detected by HST in low ionization lines such as Si ii or the high ionization lines of Si iv and C iv. Savage, Tripp & Lu (1998) showed that the absorption could occur in low density, extended gas photoionized by the UV background or in hot collisionally ionized gas in an intervening galaxy or galaxy group. Evidence for the presence of a galaxy group at $`z=0.225`$ has been provided by Schneider et al. (1992) and Tripp, Lu & Savage (1998).
The Ly$`\delta `$, Ly$`ϵ`$ and Ly$`\zeta `$ lines at $`z=0.22491`$ are detected in the FUSE spectrum. Two of these lines are shown in Figure 2. The lines are clearly resolved into 2 components with velocity separation of $`70`$ km s<sup>-1</sup>. We have not identified any metal lines arising in the $`z=0.225`$ system in the FUSE spectrum. It would be interesting to obtain FUSE short-wavelength (SiC) spectra covering $`900995`$Å, where we have the possibility of detecting the redshifted Ne viii $`\lambda \lambda `$770,780 doublet. Ne viii has a higher ionization potential than O vi (207 eV vs. 114 eV). If detected, it would indicate the presence of collisionally ionized gas at a temperature of $`T10^{5.5}`$K.
## 5 Associated Absorption
We report the detection of several “associated” ($`z_{abs}z_{em}`$) absorption lines in the spectrum of H1821+643. We identify the observed line at $`1021.45`$Å as absorption by the rest-frame EUV line O iv $`\lambda `$787.7 (see Figure 3). This line has a FWHM $`150`$ km s<sup>-1</sup>, which places it in the category of narrow absorption line (NAL) absorbers (Weymann et al., 1979; Hamann & Ferland, 2000). We have also identified O iii $`\lambda `$832.93 with an observed line at $`\lambda _{obs}1080.05`$Å. Finally, we have tentatively identified a weak line at $`\lambda _{obs}1019.85`$Å as S v $`\lambda `$786, although this line is much narrower than the O iv $`\lambda `$787 line. These lines cannot be redshifted, higher order Lyman lines, because the corresponding Ly$`\alpha `$ lines are not detected in the longer wavelength UV spectra obtained with GHRS/STIS.
Associated absorption by C iii $`\lambda `$977, Si iv $`\lambda `$1394, C iv $`\lambda \lambda `$1548,1550, O vi $`\lambda \lambda `$1032,1038 as well as several members of the Lyman series of hydrogen have been detected in H1821+643 by Savage, Tripp & Lu (1998) and Penton, Stocke & Shull (2000). The FUSE and HST observations combined show a broad range of ionization in the associated absorber(s) – including O iii, O iv, and O vi. However, no low ionization species such as Mg ii or Si ii has been observed. We see no evidence for Ne viii $`\lambda \lambda `$770,780 absorption in the FUSE data at $`z_{abs}=0.29673`$.
There are 3 possible sites of associated absorption in the H1821+643 spectrum: (1) the host galaxy of the QSO, (2) the intracluster medium (ICM), or (3) a cluster galaxy or galaxies along the line of sight to the QSO. Absorbing gas which is near the QSO central engine is often called “intrinsic” absorption. In many specific cases, there is strong spectroscopic evidence that associated NALs are intrinsic. The evidence includes: (1) time variable line strengths, (2) smooth absorption that is broad compared to the thermal line width, (3) partial covering of the continuum source, and (4) presence of excited-state absorption (Hamann & Ferland, 2000). None of these properties convincingly describes the H1821+643 associated absorber. Time-variability of absorption lines has not been reported, and cannot be addressed with this dataset. The O iv line is broader than its thermal line width, but is narrower than most intrinsic systems, which typically have widths of $`500`$ km s<sup>-1</sup>. The associated Ly$`\alpha `$ line observed in the HST/STIS spectrum obtained by Tripp, Savage & Jenkins (2000) is black at the line center (T. M. Tripp, private communication). This requires full coverage of the continuum source. Finally, excited-state absorption is not observed. All of this evidence points to the associated absorber in H1821+643 being unlike normal “intrinsic” absorbers.
Nevertheless, we still think that a likely origin of the associated absorption lies in the nearby environs of the H1821+643 host galaxy. Halos of cluster galaxies will be ram-pressure stripped by their passage through the intracluster medium. Gas in the ICM is extremely hot ($`T10^7`$K), and will have negligible ionic fractions of O iii and O iv. If a cooling flow exists in the cluster, then the O vi and less ionized species formed in this cooling gas would exist close to the cD galaxy at the cluster center – the host of H1821+643.
The photoionization models presented by Hamann (1997) indicate that an ionization parameter of $`U0.1`$ (for a wide range of continuum shapes) is needed to simultaneously produce O iv and O vi, but then the fractional ionization of O iii will be much too low to explain the absorption we detect. Hence, a multiphase model is required to explain the broad range of ionization present in the H1821+643 absorber. Multiple sites for the formation of the absorption is also consistent with the O iv $`\lambda `$787 line width.
This work is based on data obtained for the Guaranteed Time Team by the NASA-CNES-CSA FUSE mission operated by the Johns Hopkins University. Financial support to U. S. participants has been provided by NASA contract NAS5-32985.
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# Anti-holomorphic twistor and symplectic structure 2000 Mathematics Subject Classification. 53D05. Key words and phases. twistor space, pure spinor, symplectic structure
## 1 Introduction
Recently, the interest of symplectic manifolds has been growing in a perspective of Mathematical Physics related field, for example, Quantum cohomology theory, Seiberg-Witten theory etc. By definition, a manifold having a non-degenerate closed two form $`\omega `$ is called a symplectic manifold. This category of manifolds was firstly understood as that of Kähler manifolds, which has even odd betti number, for example, later on some mathematician like B. Thurston and R. Gompf constructed examples of symplectic manifolds which cannot have Kähler structure. Moreover R. Gompf \[Go\] find a systematic way of constructing symplectic manifolds and show that every finitely presented group can be realized as a fundamental group of a symplectic 4-manifold. It reveals that the symplectic category is much more bigger than Kähler one and expected to be characterized as cohomology condition of given manifold such as $`aH^2(X,𝐑)`$ and $`0a\mathrm{}aH^{2m}(X,𝐑)`$. This expectation has been broken in advent of Seiberg-Witten theory for the 4-dimensional topology. It has been known that every symplectic 4-manifold has non-zero Seiberg-Witten invariants \[T1, T2\], which indicates that condition of having symplectic structure on 4-manifolds is quite subtle. Taking closer look at the Taubes’s paper \[T1\], we can find that he was making use of the characterization of symplectic form, which is there are canonical $`Spin^c`$ structure associated almost complex structure $`J`$ and naturally induced a nowhere vanishing positive spinor $`u`$ which is harmonic ,i.e., $`/𝐃u=0`$. In this paper, we are going to show that this characterization is equivalent to the existence of the symplectic form on a given manifold. First of all, notice that symplectic form $`\omega `$ on a given manifold realized as an imaginary part of Hermitian metric for some almost complex structure $`J`$ on $`TX`$. Hence the existence of almost complex structure is a necessary condition for that of symplectic structure. Given a Riemannian even dimensional manifold, the orthogonal almost complex structure is equivalent to a section of the twistor space which is a canonical fiber bundle of $`SO(2m)/U(m)`$. We will discuss on this in the section 2. After choosing a twistor $`u`$ , equivalently having an almost complex structure $`J`$), there is the naturally associated $`Spin_2m`$ equivariant Hermitian metric on $`(TX,J)`$ and a canonical $`Spin^C`$ representation. The imaginary part of the Hermitian metric $`\omega `$ is our candidate for the symplectic form. It is easily derived that the condition for $`d\omega =0`$ is equivalent to the section $`u`$ is anti-holomorphic and harmonic ($`/𝐃u=0`$), where $`u`$ can be understood as a positive spinor of the canonical positive spinor bundle. To prove this theorem is main purpose of this paper. It also gives a simple characterization of symplectic structure on smooth 4-manifolds, which is the same as the Taubes’ analysis of symplectic form. Conclusively, the number of equations for $`\omega `$ being a symplectic form is $`m(m1)/2+m(m1)(m2)/6`$ which is bigger than that for integrability condition which is $`m(m1)/2`$. That is a little bit odd because the space of symplectic form (deformation space; it is open in the space of two form $`\mathrm{\Omega }^2(X)`$) is rather larger than the integrable complex structure which is finite dimension. On the other hand, since the deformation space is a kind of big, there are a lots of such an anti-holomorphic and harmonic twistor for some Riemannian metric on $`TX`$ once $`X`$ has a symplectic structure. It gives rise a question whether the condition we have found is “generic”, which means in symplectic manifold, the generic Riemannian metric can be induced by a symplectic form $`\omega `$ and an almost-complex structure $`J`$ associated to it.
## 2 Pure Spinor and twistor
Fix $`𝐑^n`$ be the standard inner product( $`<,>`$) real vector space and extend this metric $`𝐂`$-linearly to $`𝐂^n=R^n𝐂`$. Let $`\mathrm{𝐂𝐥}_n=Cl_n𝐂`$ be the associated complexified Clifford algebra. Let $`/𝐒_𝐂`$ be the fundamental $`\mathrm{𝐂𝐥}_n`$-module which defines the irreducible complex spinor space. For each spinor $`\sigma /𝐒_𝐂`$, we can consider the $`𝐂`$-linear map
$$j_\sigma :𝐂^n/𝐒_𝐂\text{ given by}j_\sigma (v)v\sigma $$
Generically, this map is injective. However, there are interesting spinors for which $`dim(\mathrm{ker}j_\sigma )>0`$.
###### Definition 2.1
A complex subspace $`V𝐂^n`$ is said to be isotropic (with respect to the bilinear form $`<,>`$) if $`<v,w>=0`$ for all $`v,wV`$.
We define a hermitian inner product $`(,)`$ on $`𝐂^n`$ by setting $`(v,w)=<v,\overline{w}>`$. Clearly, if $`V𝐂^n`$ is an isotropic subspace, then $`V\overline{V}`$ in this hermitian inner product. In particular, therefore, we have
$$2dim_𝐂Vn.$$
###### Definition 2.2
A spinor $`\sigma `$ is pure if $`\mathrm{ker}j_\sigma `$ is a maximal isotropc subspace, i.e., if $`dim(\mathrm{ker}j_\sigma )=\text{[n/2]}`$.
Denote by $`P/𝐒`$ the subset of pure spinors in $`/𝐒_C`$, and denote by $`_n`$ the set of maximal isotropic subspaces of $`𝐂^n`$ Both $`P/𝐒`$ and $`I_n`$ are naturally acted upon by the group $`Pin_n`$, and the assignment $`\sigma \mathrm{ker}j_\sigma `$ gives a $`Pin_n`$-equivariant map
$$K:P/𝐒_n.$$
From this point on we shall assume that $`n=2m`$ is an even integer, and furthermore that $`𝐑^{2m}`$ is oriented.
###### Definition 2.3
An orthogonal almost complex structure on $`𝐑^{2m}`$ is an orthogonal transformation $`J:𝐑^{2m}R^{2m}`$ which satisfies $`J^2=\text{Id}`$. For any such $`J`$, an associated unitary basis of $`𝐑^{2m}`$ is an ordered orthonormal basis of the form $`\{e_1,Je_1,\mathrm{}e_m,Je_m\}`$. Any two unitary bases for a given $`J`$ determine the same orientation. This is called the canonical orientation associated $`J`$.
Let $`𝒞_m`$ denote the set of all orthogonal almost complex structures on $`𝐑^{2m}`$. It is easy to see that $`𝒞_m`$ is a homogeneous space for the group $`O_{2m}`$. It falls into two connected components $`𝒞_m^+`$ and $`𝒞_m^{}`$ where $`𝒞_m^+SO_{2m}/U_m`$ consists of those almost complex structures whose canonical is positive ( i.e. agrees with given one on $`𝐑^{2m}`$). Associated to any $`J𝒞_m`$ there is a decomposition
$$𝐂^{2m}=V(J)\overline{V(J)},\text{where}$$
$$V(J)\{v𝐂^{2m}:Jv=iv\}=\{v_0+iJv_0:v_0𝐑^{2m}\}$$
There is an $`O_{2m}`$-equivalent bijection
$$𝒞_m\stackrel{V}{}_{2m}$$
which associates to $`J`$ the isotropic subspace $`V(J)`$ Let $`_{2m}^+`$ denote the component corresponding to $`𝒞_m^+`$. Using the complex volume element $`\omega _𝐂=i^me_1\mathrm{}e_{2m}`$, we have a decomposition $`/𝐒_𝐂=/𝐒_𝐂^+/𝐒_𝐂^{}`$ into $`+1`$ and $`1`$ eigenspace respectively. Easy calculation gives a decomposition $`P/𝐒=P/𝐒^+P/𝐒^{}`$ of the pure spinor space into positive and negative types. Let $`𝐏(P/𝐒^+)`$ denote the projectivization of the pure spinor space, i.e., $`𝐏(P/𝐒^+)=P/𝐒/`$ where we say that $`\sigma \sigma ^{}`$ if $`\sigma =t\sigma ^{}`$ for some $`t𝐂`$. Each of the space $`𝐏(P/𝐒^\pm ),𝒞_m^\pm `$ and $`_{2m}^\pm `$ are acted upon by $`Spin_{2m}`$, in fact by $`SO_{2m}`$.
###### Proposition 2.4
The maps $`\sigma K(\sigma )`$ and $`JV(J)`$ induce $`SO_{2m}`$-equivariant diffeomorphisms
$$𝐏(P/𝐒^+)\stackrel{K}{}_{2m}^+\stackrel{V}{}𝒞_m^+\text{ and}𝐏(P/𝐒^{})\stackrel{K}{}_{2m}^{}\stackrel{V}{}𝒞_m^{}$$
We refer to the original book \[LM\] for details.
For the sake of further discussion, we will fix $`VI_{2m}^+`$ and let $`J𝒞_m^+`$ be the associated complex structure. Choose a unitary basis $`\{e_1,Je_1,\mathrm{},e_m,Je_m\}`$ of $`𝐑^{2m}`$ and set
$$\epsilon _j=\frac{1}{\sqrt{2}}(e_jiJe_j)\overline{\epsilon }_j=\frac{1}{\sqrt{2}}(e_j+iJe_j).$$
Define
$$\omega _j=\epsilon _j\overline{\epsilon }_j\overline{\omega }_j=\overline{\epsilon }_j\epsilon _j$$
(1)
Let $`W`$ be a linear subspace invariant under multiplication by $`e_j`$ and $`Je_j`$. Then there is a hermitian orthogonal direct sum decomposition
$$W=W_jW_j^{}$$
where
$$W_j=\overline{\omega }_jW=\mathrm{ker}(\mu _{\overline{\epsilon }_j}|_W)\text{and}W_j=\omega _jW=\mathrm{ker}(\mu _{\epsilon _j}|_W)$$
and where $`\mu _{\epsilon _j}:WW`$ is defined by $`\mu _{\epsilon _j}(w)=\epsilon _jw.`$ By direct inductive calculation, we can construct
$$/𝐒_m=\mathrm{ker}(\mu _{\overline{\epsilon }_1})\mathrm{}\mathrm{ker}(\mu _{\overline{\epsilon }_m})dim_𝐂/𝐒_m=1$$
The complex volume form $`\omega _𝐂=i^me_1Je_1\mathrm{}e_mJe_m`$ has the value $`+1`$ on $`/𝐒_m`$ because $`\overline{\epsilon }_j\sigma =0ie_jJe_j\sigma =\sigma .`$ Therefore, $`/𝐒_m/𝐒_𝐂^+`$.
We clearly have that $`V(J)=\mathrm{ker}j_\sigma `$ for $`\sigma /𝐒_m`$. Hence $`/𝐒_m`$ is independent of the choice of unitary basis and the map $`V[/𝐒_m]`$ gives the desired map $`K^1`$ for the above proposition.
###### Definition 2.5
The bundle $`\tau (X)𝐏(P/𝐒^+)`$ is called the twistor space of $`X`$.
Note that $`𝐏(P/𝐒^+)`$ is an $`SO_{2m}`$-bundle and is globally defined whether or not $`X`$ is a spin manifold.
The total space of $`\tau (X)`$ carries a canonical almost complex structure defined by using the canonical decomposition of tangent space of $`\tau (X)`$, which is induced by the Riemannian connection of $`X`$.
$$T(\tau (X))=𝒱$$
where $``$ is a field of horizontal planes and $`𝒱`$ is the field of tangent planes to the fibers. As noted, $`𝒱`$ has an almost complex structure integrable on the fibers since the fiber is naturally homogeneous complex manifold ($`SO_{2m}/U(m)`$). The bundle $``$ has a “tautological” almost complex structure defined, via the identification $`\pi _{}:_JTX`$, to be the structure $`J`$ itself.
The question of integrability of $`J`$ already accomplished by M. Michelsohn.
###### Theorem 2.6
\[LM, M\] Let $`X`$ be an oriented(even-dimensional) riemannian manifold with an almost complex structure determined by a projective spinor field $`u\mathrm{\Gamma }(\tau (X))`$. Then this almost complex structure is integrable if and only if $`u`$ is holomorphic.
This will be proved in Remark 3.3. As mentioned above, $`\tau (X)`$ carries a canonical almost complex structure. Now a $`C^1`$-map between almost complex manifolds $`f:(X,J_X)(Y,J_Y)`$ will be called holomorphic (resp. anti-holomorphic) if its differential $`f_{}`$ is everywhere $`J`$-linear(resp. anti-$`J`$-linear) i.e., if $`f_{}J_X=\pm J_Yf_{}`$ respectively.
###### Remark 2.7
More succinctly one could say that cross-section of $`\tau (X)`$ induce almost complex structure, and holomorphic cross-section induce the integrable ones However, the condition that a cross-section $`u`$ be holomorphic is not linear since the complex structure on $`X`$ depends itself on $`u`$.
We will prove that the complimentary condition for the holomorphicity, which is anti-holomorphic and harmonic is equivalent to that $`u`$ induce a symplectic structure on $`X`$
###### Definition 2.8
$`\omega \mathrm{\Omega }^2(X)`$ is a symplectic form if it is nondegenerate closed form. Moreover, $`(X,\omega )`$ is called a symplectic structure on $`X`$.
Given a twistor $`u𝐏(P/𝐒)`$, there is naturally associated nodegenerate differential 2-form. It is induced by the hermitian metric with respect to the almost complex structure $`J`$ and Riemannian metric $`g`$ on $`TX`$ i.e.,
$$\omega (v,w)g(Jv,w)$$
where $`J`$ is the almost complex structure corresponding to $`s𝐏(P/𝐒)`$. Moreover it can be written as in terms of unitary basis, in other words, $`\omega =_{i=0}^me_i^{}(Je_i)^{}`$ where $`e^{}T^{}X`$ such that $`e^{}(v)=g(e,v)𝐑`$. Recall that $`\omega _j=\epsilon _j\overline{\epsilon }_j`$ for complex unitary basis $`\{\epsilon _1,\mathrm{},\epsilon _m,\overline{\epsilon }_1\mathrm{}\overline{\epsilon }_m\}`$ of $`(TX𝐂)`$. Since $`\omega _j=\epsilon _j\overline{\epsilon }_j=1ie_jJe_j`$, $`i\omega =m_j\omega _j`$.
$`\omega _1\mathrm{}\omega _m`$ $`=`$ $`{\displaystyle (1ie_jJe_j)}`$
$`=`$ $`1i{\displaystyle e_jJe_j}{\displaystyle \underset{jk}{}}(e_jJe_j)(e_kJe_k)+\mathrm{}`$
$`=`$ $`1i\omega +(1/2)(1)^2i\omega i\omega +\mathrm{}+(1/m!)(1)^mi\omega \mathrm{}i\omega `$
$`=`$ $`1i\omega +(1/2!)(i)^2\omega ^2+\mathrm{}+(1/m!)(i)^m\omega ^m`$
where $`\omega ^k=\stackrel{ktimes}{\stackrel{}{\omega \mathrm{}\omega }}\mathrm{\Omega }^{2k}(X)`$.
###### Remark 2.9
The above equality comes from the identification between $`TX`$ and $`TX^{}`$ via Riemannian metric. Note that $`(1/m!)i^m\omega ^m=i^me_1Je_1\mathrm{}e_mJe_m=\omega _𝐂`$
Note that $`_𝐂\omega ^k=k!/(mk)!\omega ^{mk}`$ i.e.,
$$d\omega =0d\omega =d^{}\omega =0_g(\omega _1+\mathrm{}+\omega _m)=0$$
where $`_g`$ is the Laplacian operator with respect to metric $`g`$. Hence we have that $`\omega `$ defines a symplectic form if and only if $`\stackrel{~}{\omega }=\omega _1+\mathrm{}+\omega _m`$ is harmonic. Our goal is to prove the following theorem.
###### Theorem 2.10
Let $`X`$ be an oriented(even-dimensional) riemannian manifold with an almost complex structure determined by a projective spinor field $`u\mathrm{\Gamma }(\tau (X))`$. Then this almost complex structure carries symplectic structure if and only if $`u`$ is harmonic and anti-holomorphic.
The product element , $`q=\overline{\omega }_1\mathrm{}\overline{\omega }_m`$ (conjugate of the above product ), of the complexified Clifford algebra $`\mathrm{𝐂𝐥}_{2m}(X)`$ can be be characterized at least locally by an element of $`qEnd(/𝐒^+)`$ such that
$$q(\sigma )=\{\begin{array}{cccc}0& & \text{ if }\hfill & \sigma s^{}/𝐒_𝐂\hfill \\ k\sigma & k𝐂^{}\hfill & \text{ and if }\hfill & [\sigma ]=s\hfill \end{array}$$
Note that we have not defined a complex spin representation $`/𝐒_𝐂`$ globally over $`X`$. Without any specification of the complex spinor bundle, the $`\overline{\omega }_1\mathrm{}\overline{\omega }_m`$ is well-defined as an element of $`\mathrm{𝐂𝐥}_{2m}(X)`$. Using the almost complex structure associated with the twistor $`u`$, we can define canonical $`spin^c`$ structure and canonical complex spin representation. Given the canonical $`Spin^c`$ representation, the product element $`q=2^muu^{}`$, which is an element of $`q\text{End}_𝐂(/𝐒^+)`$ in a way of that $`q(\alpha )=<\alpha ,u>u`$. In the next section, we will prove that $`(d\omega )u=0`$ if and only if $`/𝐃u=0`$ by using the action of $`q`$.
## 3 $`Spin^c`$ representation and proof of Theorem 2.10
Since $`Spin_n^cSpin_n\times _{𝐙_2}U(1)`$, we have a short exact sequence
$$0𝐙_2Spin_n^c\stackrel{\xi }{}SO_n\times U(1)1.$$
A principal $`SO_n`$-bundle $`P`$ carries a $`Spin^C`$ structure if any only if the $`w_2(P)`$ is the mod 2 reduction of an integral class. Given a twistor $`u𝐏(P/𝐒^+)`$, there is the canonical orthogonal almost complex structure structure $`J`$ on $`TX`$ associated with $`u`$. This $`J`$ defines a canonical $`Spin^c`$ structure $`det_𝐂TX=K_X^1`$ since the first Chern class of $`K_X^1`$ is an integral lift of the second Stiefel Whitney class, i.e., $`c_1(K_X^1)w_2(X)\text{mod 2}`$. Let $`/𝐒_𝐂`$ be the associated spinor bundle. Using the complex volume form $`\sqrt{1}^me_1Je_1\mathrm{}e_m\mathrm{}Je_m`$, we have the decomposition of $`/𝐒_𝐂`$ by the $`\pm `$-eigenspace of the complex volume element, where $`/𝐒^\pm =(1\pm \omega _𝐂)/𝐒_𝐂`$ Set
$$\epsilon _j=\frac{1}{\sqrt{2}}(e_jiJe_j)\overline{\epsilon }_j=\frac{1}{\sqrt{2}}(e_j+iJe_j).$$
be an unitary basis for $`TX`$ as above. Define
$$/𝐒_𝐂/𝐒_{i_1,\mathrm{},i_m}\mathrm{ker}(\mu _{\epsilon _{i_1}})\mathrm{}\mathrm{ker}(\mu _{\epsilon _{i_m}})$$
where $`\mu _{\epsilon _{i_k}}=\{\begin{array}{cc}\mu _{\epsilon _k}& i_k=k\hfill \\ \mu _{\overline{\epsilon _k}}& i_k=\overline{k}\hfill \end{array}`$ Let $`\sigma =\{i_1,\mathrm{}i_m\}`$ be the complex index used as above, define $`|\sigma |`$ be the number of elements of the subset $`\{i_k=k\}`$. Then we have
$$/𝐒_𝐂^+\underset{|\sigma |=2i}{}/𝐒_{ı_1\mathrm{}i_m}/𝐒_𝐂^{}\underset{|\sigma |=2i1}{}/𝐒_{ı_1\mathrm{}i_m}$$
Especially, the twistor $`u`$ is contained in $`/𝐒_{\overline{1},\mathrm{}\overline{m}}`$ which is characterized as $`\overline{\epsilon }_ju=0`$ for all $`j`$. We can express the Dirac operator in terms of the unitary basis, which follows that
$`/𝐃`$ $`=`$ $`e_j_{e_j}+Je_j_{Je_j}`$
$`=`$ $`{\displaystyle \frac{1}{2}}(\epsilon _j+\overline{\epsilon }_j)_{\epsilon _j+\overline{\epsilon }_j}{\displaystyle \frac{1}{2}}(\epsilon _j\overline{\epsilon }_j)_{\epsilon _j\overline{\epsilon }_j}`$
$`=`$ $`\overline{\epsilon }_j_{\epsilon _j}+\epsilon _j_{\overline{\epsilon }_j}`$
###### Remark 3.1
Note that the covariant derivative $``$ is $`Spin^c`$ connection which is induced from both the Levi-Civita connection and the $`U(1)`$ connection on $`K_X^1`$. It should be well-noticed that our theorem is nothing to do with a $`U(1)`$ connection. Even though the condition we have imposed is related to simply “local” question, the $`spin^c`$ structure enable us to work with globally. Furthermore, the following argument we will present below works finely without any $`spin^c`$ structure.
To define a Dirac operator on the spinors, we should specify a $`U(1)`$ connection on $`K_X^1`$. There is a canonical $`U(1)`$ connection unique up to gauge transformation $`A_0`$ such that $`<u,u>=0`$. We will abuse the notation $`/𝐃`$ for the Dirac operator, $`/𝐃_{A_0}`$, which is induced by the Levi-Civita connection and the canonical $`U(1)`$ connection $`A_0`$. Our index notation convention indicates that $`_{\stackrel{~}{e}_j}\stackrel{~}{e}_k=_l\stackrel{~}{\omega }_k^l(\stackrel{~}{e}_j)e_l`$ and $`\mathrm{\Gamma }_{j,k}^l=\stackrel{~}{\omega }_k^l(\stackrel{~}{e}_j)`$ is the Chistoffel symbol. Let $`e_j=\stackrel{~}{e}_{2j1},Je_j=\stackrel{~}{e}_{2j}`$ then $`\epsilon _j=\frac{1}{\sqrt{2}}(\stackrel{~}{e}_{2j1}i\stackrel{~}{e}_{2j})`$ where $`i=\sqrt{1}`$. Then we have
$`_{\epsilon _j}\overline{\epsilon }_k=`$ $`a_{j,k}^l\epsilon _l+c_{j,k}^l\overline{\epsilon }_l,_{\overline{\epsilon }_j}\epsilon _k=`$ $`\overline{a}_{j,k}^l\overline{\epsilon }_l+\overline{c}_{j,k}^l\epsilon _l`$
$`_{\overline{\epsilon }_j}\overline{\epsilon }_k=`$ $`b_{j,k}^l\epsilon _l+d_{j,k}^l\overline{\epsilon }_l,_{\epsilon _j}\epsilon _k=`$ $`\overline{b}_{j,k}^l\overline{\epsilon }_l+\overline{d}_{j,k}^l\epsilon _l`$
Since the Levi-Civita connection is naturally compatible with the Hermitian metric on $`TX𝐂`$, we have
$$a_{j,k}^l=<_{\epsilon _j}\overline{\epsilon }_k,\epsilon _l>=<\overline{\epsilon }_k,_{\overline{\epsilon }_j}\epsilon _l>=a_{j,l}^k$$
By the same manner, we have
$$b_{j,k}^l=b_{j,l}^k\text{ and}c_{j,k}^l=\overline{d}_{j,l}^k$$
###### Lemma 3.2
Let $`u`$ be a section of twistor space and $`J`$ be the associated orthogonal almost complex structure. Then $`u`$ is anti-holomorphic if and only if $`a_{j,k}^l=0`$ for all $`j,k,l`$ and $`u`$ is holomophic section if and only $`b_{j,k}^l=0`$ for all $`j,k,l`$.
First of all, we have to find the covariant derivative of $`u`$ which is
$$u=\frac{1}{2}\underset{k<l}{}\stackrel{~}{\omega }_k^l\stackrel{~}{e}_l\stackrel{~}{e}_ku$$
where $`\stackrel{~}{\omega }`$ is the $`so(2m)`$ connection 1-form( Levi-Civita connection with respect to $`g`$) associated with orthonomal basis $`\{\stackrel{~}{e}_1,\mathrm{}\stackrel{~}{e}_{2m}\}`$. Let
$$_{\epsilon _j}u\frac{1}{2}\underset{k<l}{}\stackrel{~}{a}_{j,k}^l\epsilon _l\epsilon _ku\text{mod }<u>$$
The coefficient $`\stackrel{~}{a}_{j,k}^l`$ can be derived as follows,
$$\overline{\epsilon }_tu=0\text{for all }t$$
By taking covariant derivative $`_{\epsilon _j}`$, we have
$$(_{\epsilon _j}\overline{\epsilon }_t)u+\overline{\epsilon }_t_{\epsilon _j}u=0$$
Hence
$`(_{\epsilon _j}\overline{\epsilon }_t)u`$ $`=`$ $`\overline{\epsilon }_t{\displaystyle \underset{k<l}{}}{\displaystyle \frac{1}{2}}\stackrel{~}{a}_{j,k}^l\epsilon _l\epsilon _ku`$
$`=`$ $`{\displaystyle \frac{1}{2}\stackrel{~}{a}_{j,k}^l\overline{\epsilon }_t\epsilon _l\epsilon _ku}`$
$`=`$ $`\{\begin{array}{cc}\stackrel{~}{a}_{j,k}^l\epsilon _lu& \text{ for }k=t<l\hfill \\ \stackrel{~}{a}_{j,k}^t\epsilon _ku& \text{ for }l=t>k\hfill \end{array}`$
Since $`\overline{\omega }_j=\overline{\epsilon }_j\epsilon _ju=2u`$. Therefore
$$<_{\epsilon _j}\overline{\epsilon }_t,\epsilon _s>=a_{j,t}^s=\stackrel{~}{a}_{j,t}^s$$
We get $`\stackrel{~}{a}_{j,k}^l=<_{\epsilon _j}\epsilon _k,\epsilon _l>.`$ By the analogous method, we can get
$$_{\overline{\epsilon }_j}u\frac{1}{2}\underset{k<l}{}b_{j,k}^l\epsilon _l\epsilon _ku\text{mod }<u>.$$
With this understood, it can be rephrased that $`u`$ is anti-holomorphic $``$ $`\overline{\epsilon }_t_{\epsilon _j}u=0`$ $``$ $`a_{j,k}^l=0<_{\epsilon _j}\overline{\epsilon }_k,\epsilon _l>=0`$ for all $`j,k,l`$. Also $`u`$ is holomorphic $``$ $`\overline{\epsilon _t}_{\overline{\epsilon }_j}u=0`$ $``$ $`a_{j,k}^l=0<_{\epsilon _j}\overline{\epsilon }_k,\epsilon _l>=0`$ for all $`j,k,l`$.
###### Remark 3.3
From the torsion free condition of Levi-Civita connection, we have
$$b_{j,k}^lb_{k,j}^l=<_{\overline{\epsilon }_j}\overline{\epsilon }_k_{\overline{\epsilon }_k}\overline{\epsilon }_j,\epsilon _l>=<[\overline{\epsilon }_j,\overline{\epsilon }_k],\epsilon _l>.$$
Since the anti-commutativity between upper index and right lower index $`b_{j,k}^l=b_{j,l}^k`$, we can get an equivalent condition which says that $`b_{j,k}^lb_{k,j}^l=0`$ if and only if $`b_{j,k}^l=0.`$ Hence it is easy to prove the Theorem 2.6 from the above equation.
We want to find an equivalent condition for the harmonic two form $`\omega `$ i.e. $`/𝐃\omega =(d+d^{})\omega =0`$, where $`d^{}`$ is the formal adjoint of $`d`$ with respect to $`g`$. The following lemma is about it.
###### Proposition 3.4
Let $`\omega =m+_k\omega _k`$ be the purely imaginary part of the Hermitian metric. Then $`/𝐃\omega =0`$ if and only if $`a_{j,k}^l=0`$ and $`b_{j,k}^l+b_{l,j}^k+b_{k,l}^j=0`$ for all $`j,k,l`$.
Proof: Since $`\omega `$ is purely imaginary two form, we have
$`/𝐃\omega `$ $`=`$ $`{\displaystyle \underset{j}{}}\epsilon _j_{\overline{\epsilon }_j}\omega +\overline{\epsilon }_j_{\epsilon _j}\omega `$
$`=`$ $`{\displaystyle \underset{j}{}}\epsilon _j_{\overline{\epsilon }_j}\omega \overline{{\displaystyle \underset{j}{}}\epsilon _j_{\overline{\epsilon }_j}\omega }`$
$`=`$ $`2i\text{ Im }/𝐃^{\frac{1}{2}}\omega `$
It suffices to consider the half part of the Dirac operator, it reads
$`/𝐃^{\frac{1}{2}}\omega `$ $`=`$ $`{\displaystyle \underset{j}{}}\epsilon _j_{\overline{\epsilon }_j}\omega `$
$`=`$ $`{\displaystyle \underset{j,k}{}}(\epsilon _j(_{\overline{\epsilon }_j}\epsilon _k)\overline{\epsilon }_k+\epsilon _j\epsilon _k_{\overline{\epsilon }_j}\overline{\epsilon }_k)`$
$`=`$ $`{\displaystyle \underset{j,k,l}{}}(\overline{a}_{j,k}^l\epsilon _j\overline{\epsilon }_l\overline{\epsilon }_k+\overline{c}_{j,k}^l\epsilon _j\epsilon _l\overline{\epsilon }_k+b_{j,k}^l\epsilon _j\epsilon _k\epsilon _l+d_{j,k}^l\epsilon _j\epsilon _k\overline{\epsilon }_k)`$
$`=`$ $`{\displaystyle \underset{j,k,l}{}}(\overline{a}_{j,k}^l\epsilon _j\overline{\epsilon }_l\overline{\epsilon }_k+b_{j,k}^l\epsilon _j\epsilon _k\epsilon _l)+{\displaystyle \underset{j,k,l}{}}(\overline{c}_{j,l}^k+d_{j,k}^l)\epsilon _j\epsilon _k\overline{\epsilon }_l`$
$`=`$ $`{\displaystyle \underset{i,j,l}{}}(\overline{a}_{j,k}^l\epsilon _j\overline{\epsilon }_l\overline{\epsilon }_k+b_{j,k}^l\epsilon _j\epsilon _k\epsilon _l)`$
Hence $`/𝐃\omega =0`$ if and only if $`/𝐃^{\frac{1}{2}}\omega =0a_{j,k}^l=0`$ for all $`i,j,k`$ and $`{\displaystyle \underset{\sigma }{}}b_{\sigma (j),\sigma (k)}^{\sigma (l)}`$ where $`\sigma `$ is the permutation of $`i,j,k`$. The relation $`b_{j,k}^l=b_{j,l}^k`$ completes the proposition.
Proof of Theorem 2.10 Since $`<u,u>=0`$, we have
$`/𝐃u`$ $`=`$ $`{\displaystyle \underset{j}{}}(\overline{\epsilon }_j_{\epsilon _j}u+\epsilon _j_{\overline{\epsilon }_j}u)`$
$`=`$ $`{\displaystyle \underset{j,k,l}{}}({\displaystyle \frac{1}{4}}a_{j,k}^l\overline{\epsilon }_j\epsilon _k\epsilon _lu+{\displaystyle \frac{1}{4}}b_{j,k}^l\epsilon _j\epsilon _k\epsilon _lu)`$
$`=`$ $`{\displaystyle \underset{j,k}{}}a_{j,j}^k\epsilon _ku+{\displaystyle \underset{j<k<l}{}}{\displaystyle \frac{1}{2}}(b_{j,k}^l+b_{l,j}^k+b_{k,l}^j)\epsilon _j\epsilon _k\epsilon _lu`$
Hence $`u`$ is anti-holomorphic pure spinor ($`\overline{\epsilon }_t_{\epsilon _j}=0`$ for all $`t,j`$) and harmonic ($`/𝐃u=0`$ ) gives an equivalent condition for $`\omega `$ being a symplectic form. Note that given symplectic manifold $`(X,\omega )`$ has such a anti-holomorphic and harmonic twistor $`u`$ by choosing any almost complex structure which calibrate $`\omega `$.
###### Corollary 3.5
$`(d\omega )u=0`$ if and only if $`u`$ is harmonic , i.e., $`/𝐃u=0`$.
Proof: Let $`q=_{j=1}^m\overline{\omega _j}=_j(1+ie_jJe_j)`$. Using the action $`q`$ on $`u`$, $`qu=2^mu`$, and taking Dirac operator on the both side, we can have
$`/𝐃qu`$ $`=`$ $`(/𝐃q)u+{\displaystyle }\stackrel{~}{e}_jq_{\stackrel{~}{e}_j}u`$
$`=`$ $`/𝐃qu(<u,u>=0q_{\stackrel{~}{e}_j}u=0)`$
$`=`$ $`/𝐃u`$
Thus $`/𝐃u=0`$ if and only if $`(/𝐃q)u=0`$. Moreover since $`(3i)^m(1)^{\frac{1}{2}p(p+1)}\phi \omega _𝐂=\phi `$ for $`\phi \mathrm{\Omega }^p(X)`$ and $`1/k!\omega ^k=1/(mk)!\omega ^{mk}`$, we have
$`/𝐃q`$ $`=`$ $`(d+d^{})q=id\omega +i^2/2!d\omega ^2+\mathrm{}+i^{m1}/(m1)!d\omega ^{m1}`$
$`i/(m1)!d\omega ^{m1}i^2/(m2)!d\omega ^{m2}\mathrm{}i^{m1}d\omega `$
$`=`$ $`id\omega +i^2/2!d\omega ^2+\mathrm{}+i^{m1}/(m1)!d\omega ^{m1}`$
$`i^{4m+1}d\omega \omega _𝐂+i^{4m+2}/2!d\omega ^2\omega _𝐂\mathrm{}+i^{m1}/(m1)!d\omega ^{m1}\omega _𝐂`$
Since $`\omega _𝐂u=u,\omega u=(mi)u`$, we have $`d\omega ^ku=k(d\omega )\omega ^{k1}u=k(mi)^{k1}d\omega u`$. Thus
$$(/𝐃q)u=2i(1+m+m^2/2!+\mathrm{}+m^{m2}/(m2)!)d\omega u.$$
This completes the proof.
###### Remark 3.6
In dimension $`2m6`$ every non-zero positive( or negative ) spinor is pure, i.e., $`P/𝐒^\pm =/𝐒_𝐂^\pm 0`$. This is simply because the group $`Spin_{2m}`$ acts transitively on the unit sphere in $`/𝐒_𝐂^\pm `$ in these dimensions.
In dimension 4, since $`\phi \mathrm{\Omega }^3(X,𝐑)`$ acts on $`u`$ injectively, we get $`(d\omega )u=0`$ if and only if $`d\omega =0`$ Hence the harmonic spinor $`u`$, equivalently anti-holomorphic twistor, gives a sufficient condition to induce a symplectic structure. The next corollary follows from it.
###### Corollary 3.7
In dimension 4, Let $`u`$ be a nowhere vanishing section of positive complex spinor bundle. Then $`/𝐃u=0`$ and $`<u,u>=0`$ then $`X`$ is symplectic 4-manifold.
Finally, suppose $`u=0`$, then $`u`$ is then both holomorphic and anti-holomorphic twistor. We have following corollary, which is proposition 9.8 in \[LM\].
###### Corollary 3.8
Suppose $`u`$ is parallel, then $`(X,g,J)`$ becomes a Kähler manifold.
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# Short-range spin correlations and induced local spin-singlet amplitude in the Hubbard model
\[
## Abstract
In this paper, from the microscopic Hubbard Hamiltonian we extract the local spin-singlet amplitude due to short-range spin correlations, and quantify its strength near half-filling. As a first application of the present approach, we study a problem of the energy dispersion and its d-wave modulation in the insulating cuprates, Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> and Ca<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub>. Without any adjustable parameters, most puzzling issues are naturally and quantitatively explained within the present approach.
\]
Recent discovery of a normal state pseudogap in underdoped copper oxides has attracted considerable attention both from experimentalists and theoretical physicists for many years. For these materials, the low frequency spectral weight begins to be strongly suppressed below some characteristic temperature $`T^{}`$ higher than $`T_c`$. This anomalous behavior has been observed through various experimental probes such as angle resolved photoemission spectroscopy (ARPES), specific heat, tunneling, NMR, and optical conductivity. One of the most puzzling questions in pseudogap phenomena is why $`T^{}`$ has a completely different doping dependence from $`T_c`$, in spite of possibly their close relation.
There are many theoretical scenarios to attempt to understand the pseudogap phenomena. These include the spinon pair formation without the Bose-Einstein condensation of holons in the slave boson (or fermion) mean-field theory, strong superconducting (SC) phase fluctuations, strong magnetic fluctuations near the antiferromagnetic (AF) instability, and so on. At present there is no consensus on the origin of the pseudogap and its relationship with the SC long-range order. Among several scenarios, the slave boson (or fermion) mean-field theory may shed some insight into the problem. This is because the pseudogap is closely related to a spin gap, the predicted phase diagram is, at least, qualitatively consistent with experiments, and furthermore it starts from the microscopic model ($`tJ`$ Hamiltonian) as opposed to other phenomenological models. On the other hand, the recombination of a spinon and a holon into a physical electron is nontrivial, and also the constraint of no-double-occupancy at each site is difficult to impose at a microscopic level. In deriving a local spin-singlet amplitude which may be responsible for the pseudogap behavior, in this paper we use the Hubbard Hamiltonian instead of the $`tJ`$ Hamiltonian. Local spin-singlet amplitude is induced directly from short-range spin correlations in the normal state and its strength is quantified near half-filling. In this paper short-range spin correlation means that when site $`i`$ is occupied by an electron with up-spin (or down-spin), then the nearest sites predominantly by electrons with the opposite spin.
We start by defining the one-band Hubbard model proposed by Anderson as the simplest model which might capture the correct low energy physics of copper oxides. The Hubbard model is described by the Hamiltonian where $`c_{i,\sigma }`$ destroys an electron at site $`i`$ with spin $`\sigma `$ on a two-dimensional square lattice
$`H=t{\displaystyle \underset{<i,j>,\sigma }{}}c_{i,\sigma }^{}c_{j,\sigma }+U{\displaystyle \underset{i}{}}c_{i,}^{}c_{i,}c_{i,}^{}c_{i,}.`$ (1)
$`t`$ is a hopping parameter between nearest neighbors $`<i,j>`$ and $`U`$ denotes local Coulomb repulsion. It is believed that the realistic strength of the Coulomb repulsion lies in between the weak and strong coupling regimes, namely, $`UW2W`$ where $`W`$ is the bandwidth of $`8t`$ in two dimensions.
As a first step to the microscopic understanding of the pseudogap, it is important to answer a more fundamental question: How can spin-singlet tendency appear directly from the Coulomb repulsion, presumably without the exchange of bosonic degrees of freedom such as spin fluctuations? As a simple example to illustrate this point qualitatively, let us consider the $`Ut`$ limit at half-filling where only an up-spin or down-spin electron is allowed at each site. The typical spin-singlet structure (with $`d`$-wave form factor) in which we are interested in this paper is $`\mathrm{\Delta }_g(i)=\frac{1}{2}_\delta g(\delta )(c_{i+\delta ,}c_{i,}c_{i+\delta ,}c_{i,})`$, where $`g(\delta )`$ is an appropriate structure factor. For a d-wave symmetry, for instance,
$`g(\delta )=\{\begin{array}{ccc}1/2\hfill & \text{if }\delta =(\pm 1,0),\hfill & \\ 1/2\hfill & \text{if }\delta =(0,\pm 1),\hfill & \\ 0\hfill & \text{if otherwise}.\hfill & \end{array}`$ (5)
Then the local spin-singlet amplitude for strongly correlated electrons ($`|\mathrm{\Delta }_g(i)|=(1\times 1+0\times 0)/2=1/2`$) is increased over its noninteracting value ($`|\mathrm{\Delta }_g(i)|_0=(1/2\times 1/2+1/2\times 1/2)/2=1/4`$ for both spins).
To make this argument more quantitative, we introduce a spin-singlet correlation function
$`\chi _g(i,\tau )=T_\tau \mathrm{\Delta }_g(i,\tau )\mathrm{\Delta }_g^{}(0,0),`$ (6)
where $`T_\tau `$ is the imaginary time ordering operator. In this paper we consider only the local spin-singlet amplitude $`|\mathrm{\Delta }_g(0)|^2`$, which may be obtained by $`\chi _g(i0,\tau 0^{})`$. Now we examine whether there is an increase in $`|\mathrm{\Delta }_g(0)|^2`$ for strongly correlated electrons with respect to $`|\mathrm{\Delta }_g(0)|^2_0`$ for the noninteracting electrons, which is similar in spirit to a renormalization group (RG) approach. As already illustrated in the previous paragraph, the increase in $`|\mathrm{\Delta }_g(0)|^2`$ for strongly correlated electrons crucially depends on the short-range spin correlations. Although there exists no controlled way of obtaining the wave function in general, certain local or short-range static quantities such as double occupancy (or equivalently local spin amplitude) or the nearest neighbor correlations are reasonably well captured by the mean-field state with AF order. Before going further, it is important to establish the validity of the current approximation by explicitly comparing the above quantities calculated in the mean-field state (with AF order) with available quantum Monte Carlo (QMC) results.
The double occupancy $`D=n_{i,}n_{i,}`$ plays an important role in gauging the degree of strong correlations in the Hubbard-type model. For instance, $`D0`$ at half-filling for $`U\mathrm{}`$, while $`D_0`$ evaluated in the noninteracting state is $`(n/2)^2=0.25`$. In Fig. 1(a) our calculations (solid curve) are compared with (virtually exact) QMC results (open circles) by White et al. for $`n=1`$, $`N=4\times 4`$, and $`T=t/16`$. For purely interaction induced effect, $`D_0`$ is subtracted from $`D`$. Above $`U3t`$ the agreement with QMC data is increasingly better and for $`U8t`$ the two results are almost indistinguishable. For the nearest neighbor correlations between $`i`$ and $`j`$, we calculate $`t_{eff}/t=c_{i,\sigma }^{}c_{j,\sigma }_U/c_{i,\sigma }^{}c_{j,\sigma }_{U=0}`$ for comparison with available QMC results. Again $`t_{eff}/t`$ evaluated in the noninteracting state, namely, unity is subtracted from $`t_{eff}/t`$ for purely interaction induced effect. Figure 1(b) shows our results evaluated in the mean-field state (solid curve) together with QMC results (open circles) by White et al. From intermediate to very strong coupling, the deviation with QMC results is less than $`10\%`$ level. In fact the interaction induced local spin-singlet amplitude $`|\mathrm{\Delta }_g(0)|^2`$ -$`|\mathrm{\Delta }_g(0)|^2_0`$ involves mainly nearest neighbor static correlations between electrons. Note that the spin density wave (SDW) approximation will be used below only to capture the reasonable local and short-range correlations between electrons.
We should also point out some limitations in this approach. The physical reason that the nearest neighbor correlation is overestimated (by up to $`10\%`$) compared with QMC data is as follows. Suppose site $`i`$ is occupied by an electron with up-spin in the strong coupling limit. Then the nearest sites are occupied by down-spin electrons as majority and also by up-spin electrons as minority. As the distance from site $`i`$ increases, the effective polarization of electron spins decreases in magnitude. In the mean-field state with AF order, however, the effective polarization is the same for all the distances from site $`i`$. This means that physical quantities evaluated in the mean-field state become progressively overestimated for increasingly distant correlations. Another limitation is that the mean-field critical doping $`x_c`$ where AF mean-field order vanishes increases with increasing $`U`$. This is inconsistent with more advanced treatment of the Hubbard model where with increasing $`U`$, $`x_c`$ decreases after reaching its maximum around at $`x0.21`$. Hence our approximation based on the Hubbard model is valid near half-filling (probably up to $`x0.10.15`$). For the $`tJ`$ Hamiltonian the current approximation will be valid even far away from half-filling, because this strong coupling feature is already taken into account there.
Then for $`d`$-wave type symmetries, the local spin-singlet amplitude $`|\mathrm{\Delta }_g(0)|^2`$ evaluated in the mean-state with AF order becomes
$`{\displaystyle \frac{n}{4N}}{\displaystyle \underset{\stackrel{}{k}}{\overset{^{}}{}}}\varphi _g^2(\stackrel{}{k})[f(E_{}(\stackrel{}{k}))+f(E_+(\stackrel{}{k}))]sign[{\displaystyle \frac{\varphi _g(\stackrel{}{k}+\stackrel{}{Q})}{\varphi _g(\stackrel{}{k})}}]`$ (7)
$`\times `$ $`{\displaystyle \frac{\mathrm{\Delta }_{sdw}}{2UN}}{\displaystyle \underset{\stackrel{}{k}}{\overset{^{}}{}}}\varphi _g^2(\stackrel{}{k}){\displaystyle \frac{\mathrm{\Delta }_{sdw}}{\lambda (\stackrel{}{k})}}[f(E_{}(\stackrel{}{k}))f(E_+(\stackrel{}{k}))],`$ (8)
where
$`\lambda (\stackrel{}{k})`$ $`=`$ $`\sqrt{((\epsilon (\stackrel{}{k})\epsilon (\stackrel{}{k}+\stackrel{}{Q}))/2)^2+\mathrm{\Delta }_{sdw}^2},`$ (9)
$`E_\pm (\stackrel{}{k})`$ $`=`$ $`(\epsilon (\stackrel{}{k})+\epsilon (\stackrel{}{k}+\stackrel{}{Q}))/2\pm \lambda (\stackrel{}{k}).`$ (10)
$`\epsilon (\stackrel{}{k})=2t(\mathrm{cos}k_x+\mathrm{cos}k_y)\mu `$ for nearest neighbor hopping, $`\mu `$ is the chemical potential controlling the particle density $`n`$, $`N`$ the total number of lattice sites, $`f(E)`$ the Fermi-Dirac distribution function, $`\stackrel{}{Q}`$ the AF wave vector $`(\pi ,\pi )`$, $`\varphi _g(\stackrel{}{k})`$ the Fourier transform of $`g_\delta (i)`$, and the summation accompanied by the prime symbol is over wave vectors in half of the first Brillouin zone. The SDW gap $`\mathrm{\Delta }_{sdw}`$ and the chemical potential $`\mu `$ are self-consistently determined through the gap and number equations for given $`U`$, $`T`$ and $`n`$. Figure 2(a) shows the local spin-singlet amplitude $`|\mathrm{\Delta }_d|^2`$ subtracted by that evaluated in the noninteracting state for purely interaction induced effect. The most interesting finding is that the local spin-singlet amplitude is induced directly by short-range spin correlations without any explicit driving (attractive) interaction. The short-range spin correlations are ultimately ascribed to the strong local Coulomb repulsion $`U`$. The interaction induced $`|\mathrm{\Delta }_d|^2`$ rapidly decreases with increasing doping, and if the second limitation in the current approximation mentioned above is properly corrected, it will decrease even faster with doping.
Because the value of the first term is nearly the same as $`|\mathrm{\Delta }_g^+(0)|^2_0`$ for the noninteracting electrons, the positivity of the second term leads to an increase in the local spin-singlet amplitude. For $`d_{xy}`$ ($`\varphi (\stackrel{}{k})=2\mathrm{sin}k_x\mathrm{sin}k_y`$) symmetry, for example, the local spin-singlet amplitude is even suppressed due to $`\varphi (\stackrel{}{k}+\stackrel{}{Q})=\varphi (\stackrel{}{k})`$. We found that among several symmetries including local s-wave, extended s-wave, $`d_{xy}`$, and $`d_{x^2y^2}`$, the $`d_{x^2y^2}`$ ($`\varphi _d(\stackrel{}{k})=\mathrm{cos}k_x\mathrm{cos}k_y`$) symmetry shows the largest increase in the local spin-singlet amplitude. Likewise the local spin amplitude $`(S_z(i))^2`$ is also induced from the same Coulomb repulsion.
In order to determine the effective strength of the induced local spin-singlet amplitude (with $`d`$-wave symmetry) corresponding to Fig. 2(a), we consider a Hamiltonian
$`H=t{\displaystyle \underset{<i,j>,\sigma }{}}c_{i,\sigma }^+c_{j,\sigma }+V_{ind}{\displaystyle \underset{i}{}}\mathrm{\Delta }_d^+(i)\mathrm{\Delta }_d(i)`$ (11)
with the same parameters ($`t`$, $`T`$, and $`n`$) as before. Now the strategy is to find $`V_{ind}`$ in such a manner that $`V_{ind}`$ gives the same interaction induced local spin-singlet amplitude $`|\mathrm{\Delta }_d|^2`$ -$`|\mathrm{\Delta }_d|^2_0`$ as in Fig. 2(a). Please note that $`V_{ind}`$ is the effective strength of the induced local spin-singlet amplitude, but not attractive interaction yielding a local SC pair in the Bose-Einstein limit. As before, we may gauge the validity of a given approximation (BCS approximation in this case) by explicitly comparing some local and short-range static quantities calculated in that approximation with available numerical data for the above Hamiltonian. Unfortunately there exist no such numerical results. Then one may use general knowledge which has been obtained through study on the repulsive and attractive Hubbard models in the SDW and BCS approximations, respectively. The SDW approximation for the repulsive Hubbard model which was discussed in details in the previous paragraphs, suggests that for intermediate to strong coupling the mean-field approximation reasonably well captures both local and short-range static quantities. For the BCS approximation of the attractive Hubbard model, it has also been shown that the local static quantity (double occupancy) in that approximation is in good agreement with QMC data for $`|U|3t`$. However, the numerical results of nearest neighbor correlations in the attractive Hubbard model are not available in the literature to our knowledge. In spite of lack of complete numerical data for the attractive Hubbard model and for the Hamiltonian Eq. 11, one may safely expect that the BCS mean-field approximation to Eq. 11 will give at least reasonable local and short-range static quantities for intermediate to strong coupling. Note that the mean-field (BCS) approximation will be used below only to capture the reasonable local and short-range correlations between electrons, as before in the SDW approximation.
In the BCS approximation for the above Hamiltonian, $`|\mathrm{\Delta }_d(0)|^2`$ is given as
$`{\displaystyle \frac{n}{4N}}{\displaystyle \underset{\stackrel{}{k}}{}}\varphi _d^2(\stackrel{}{k})[{\displaystyle \frac{1}{2}}{\displaystyle \frac{\epsilon (\stackrel{}{k})}{2E(\stackrel{}{k})}}\mathrm{tanh}(E(\stackrel{}{k})/2T)]`$ (12)
$`+`$ $`[{\displaystyle \frac{1}{N}}{\displaystyle \underset{\stackrel{}{k}}{}}\varphi _d(\stackrel{}{k}){\displaystyle \frac{\mathrm{\Delta }_d(\stackrel{}{k})}{2E(\stackrel{}{k})}}\mathrm{tanh}(E(\stackrel{}{k})/2T)]^2,`$ (13)
where $`E(\stackrel{}{k})=\sqrt{\epsilon ^2(\stackrel{}{k})+\mathrm{\Delta }_d^2(\stackrel{}{k})}`$. As before, the d-wave gap $`\mathrm{\Delta }_d(\stackrel{}{k})=\mathrm{\Delta }\varphi _d(\stackrel{}{k})`$ and the chemical potential $`\mu `$ are self-consistently determined through the gap and number equations for given $`V_{ind}`$, $`T`$ and $`n`$. In Fig. 2(b) $`V_{ind}`$ is plotted as a function of doping $`x=1n`$ for $`T=0`$ and $`U=1.5W=12t`$ near half-filling. The BCS approximation ($`|V_{ind}|3t`$) is expected to be accurate for $`x0.1`$ and to be less accurate beyond it. The strength of the induced local spin-singlet amplitude rapidly decreases with doping, just as the induced AF correlations do. Beyond $`x0.10.15`$, $`|V_{ind}|`$ will decrease even faster when the strong coupling effect is correctly included, as was mentioned in Fig. 2(a). For $`Ut`$ at $`n=1`$, $`V_{ind}`$ saturates to be $`7.64t`$ and $`|\mathrm{\Delta }_d(0)|^20.25`$, consistent with our previous result ($`|\mathrm{\Delta }_d(0)|=0.5`$) based on qualitative argument. $`V_{ind}`$ and its doping dependence have some interesting consequences which may answer some of the puzzling issues in the high temperature superconductors.
As a first application of the present approach, we attack a problem of the energy dispersion and its d-wave modulation in the insulating cuprates, Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> and Ca<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub>. Recent ARPES experiments for an insulating cuprate Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> clearly show that the near isotropy and the overall band dispersion along $`(\pi /2,\pi /2)(\pi ,0)`$ and $`(\pi /2,\pi /2)(0,0)`$ cannot be explained by considering only AF order or its fluctuations, unless some adjustable fitting parameters such as $`t^{}`$ and $`t^{\prime \prime }`$ are introduced. Furthermore a d-wave-like modulation of the insulating gap in Ca<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> is totally mysterious from that point of view. Because the induced local spin-singlet and spin amplitudes increase with decreasing doping, near half-filling and at low temperatures, both short-range spin-singlet and AF fluctuations are strong and coexist, leading to a resonating valence bond (RVB) type state. Our general experience tells that as long as the energy dispersion is concerned, it is basically identical in a strongly fluctuating state (or in the pseudogap state) and in a long-range ordered state. The only difference is that in the former state the gap is filled with some spectral weight (pseudogap) and the spectral function is spilt into two (relatively broad) peaks instead of two delta functions. Thus at half-filling and at low temperatures, a hole strongly interacts with both short-range spin-singlet and AF fluctuations, yielding an energy dispersion similar to that in the coexistence state of the SC and AF long-range order. The energy dispersion is then given by
$`\sqrt{\epsilon ^2(\stackrel{}{k})+\mathrm{\Delta }_{sdw}^2+\mathrm{\Delta }^2(\mathrm{cos}k_x\mathrm{cos}k_y)^2}.`$ (14)
From the overall bandwidth ($``$ 320 meV) experimentally measured from $`(\pi /2,\pi /2)`$ to $`(0,0)`$ along which $`\varphi _d(\stackrel{}{k})`$ vanishes, the Coulomb repulsion $`U`$ (and $`\mathrm{\Delta }_{sdw}U/2`$) is determined to be $`11.4t`$ for $`t=`$250 meV. For $`U=11.4t`$, $`n=1`$ and $`T=0`$, the effective strength $`V_{ind}`$ of the induced spin-singlet amplitude is completely determined to give $`V_{ind}=6.09t`$ and $`\mathrm{\Delta }=2.07t`$. Figure 3 clearly demonstrates that the energy dispersion and its near isotropy along $`(\pi /2,\pi /2)(\pi ,0)`$ and $`(\pi /2,\pi /2)(0,0)`$ are in excellent agreement with experiments without any adjustable parameters. Our energy dispersion along $`(\pi /2,\pi /2)(\pi ,0)`$ is proportional to $`(\mathrm{cos}k_x\mathrm{cos}k_y)^2`$ (solid curve in the inset) at low energies as opposed to $`|\mathrm{cos}k_x\mathrm{cos}k_y|`$ (dashed curve) predicted from the flux-phase in the $`tJ`$ model. Recent experiments show a significant deviation of the energy dispersion from the cusp-like form along $`(\pi /2,\pi /2)(\pi ,0)`$. A similar result to ours was recently obtained in the context of SO(5) symmetry. The state with a spin gap (or pseudogap) of $`(\mathrm{\Delta }\varphi _d)^2/UJ\varphi _d^2`$ and also with a Mott-Hubbard gap of order $`U`$ at half-filling is expected to continuously evolve into a state with a relatively smaller spin gap (or pseudogap) but without a charge gap away from half-filling.
In this paper we have considered only the local spin-singlet amplitude induced by short-range spin correlations and its consequences. The long-range $`d`$-wave superconductivity is beyond the scope of the present approach. How local spin-singlets acquire local SC phases and eventually establish their long-range phase coherence is a challenging problem to the theory of high temperature superconductivity. Quantitative analytical study of the Hubbard model in the physically relevant regime and of the interplay between AF and $`d`$-wave pairing correlations beyond a mean-field level is not available at present, mainly due to the absence of a small parameter. Certainly it is a subject for future study. It is also desirable to perform numerical simulations for a Hamiltonian with $`d`$-wave symmetry such as Eq. 11 if it is possible. Then the mean-field BCS approximation used here to obtain reasonable local and short-range static correlations between electrons will be tested for its range of validity.
In summary, from the microscopic Hubbard Hamiltonian, the local spin-singlet amplitude due to short-range spin correlations was explicitly extracted and its strength was quantified near half-filling. As a first application of the present approach, a problem of the energy dispersion and its d-wave modulation in the insulating cuprates, Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> and Ca<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> was studied. Without any adjustable parameters, most puzzling issues on the energy dispersion were naturally and quantitatively explained within the present approach.
The author would like to thank A. M. Tremblay for numerous help and discussions throughout this work. The author also thanks C. Bourbonnais, E. Dagotto, R. Gooding, M. Norman, D. Sénéchal, and W. Stephan for stimulating discussions. The present work was supported by a grant from the Natural Sciences and Engineering Research Council (NSERC) of Canada and the Fonds pour la formation de Chercheurs et d’Aide à la Recherche (FCAR) of the Québec government.
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# 1 Introduction
## 1 Introduction
Let $`YX`$ be a smooth fibre bundle of some classical field model. We study cohomology of exterior forms on the infinite-order jet space $`J^{\mathrm{}}Y`$ of $`YX`$. This cohomology plays an important role in the field-antifield BRST formalism for constructing the descent equations \[1-4\].
In the framework of this BRST formalism, one considers the so-called horizontal complex
$$0R𝒪_{\mathrm{}}^0\stackrel{d_H}{}𝒪_{\mathrm{}}^{0,1}\stackrel{d_H}{}\mathrm{}\stackrel{d_H}{}𝒪_{\mathrm{}}^{0,n},$$
(1)
where $`𝒪_{\mathrm{}}^{0,}`$ is a subalgebra of horizontal (semibasic) exterior forms on $`J^{\mathrm{}}Y`$, and $`d_H`$ is the horizontal (total) differential. Extended to the jet space of ghosts and antifields, these forms and their cohomology are called local forms and local cohomology. Given the BRST operator $`𝐬`$, one defines the total BRST operator $`𝐬+d_H`$ and examines the BRST cohomology modulo $`d_H`$.
It should be emphasized that, the above mentioned BRST formalism is formulated on a contractible fibre bundle $`Y=R^{n+m}R^n`$. Of course, this is not the generic case of gauge theory and its outcomes to topological field models and anomalies. The key point is that, in this case, the horizontal complex (1) is exact. This fact called the algebraic Poincaré lemma is crucial for constructing the (local) descent equations in BRST theory. To write global descent equations in a non-trivial topological context, one should study (global) cohomology of the complex (1). The question is that there are (at least) two classes of exterior forms on the infinite-order jet space $`J^{\mathrm{}}Y`$, and both of them fail to be differential forms on $`J^{\mathrm{}}Y`$ in a rigorous sense because $`J^{\mathrm{}}Y`$ is not a Banach manifold \[5-7\].
Recall that the infinite-order jet space of a smooth fibre bundle $`YX`$ is defined as a projective limit $`(J^{\mathrm{}}Y,\pi _j^{\mathrm{}})`$ of the surjective inverse system
$$X\stackrel{\pi }{}Y\stackrel{\pi _0^1}{}\mathrm{}J^{r1}Y\stackrel{\pi _{r1}^r}{}J^rY\mathrm{}$$
(2)
of finite-order jet manifolds $`J^rY`$ \[5-8\]. Provided with the projective limit topology, $`J^{\mathrm{}}Y`$ is a paracompact Fréchet (but not Banach) manifold . Its paracompactness will be an essential tool for cohomological calculations below. Given a bundle coordinate chart $`(\pi ^1(U_X);x^\lambda ,y^i)`$ on the fibre bundle $`YX`$, we have the coordinate chart $`((\pi ^{\mathrm{}})^1(U_X);x^\lambda ,y_\mathrm{\Lambda }^i)`$, $`0|\mathrm{\Lambda }|,`$ on $`J^{\mathrm{}}Y`$, together with the transition functions
$$y_{}^{}{}_{\lambda +\mathrm{\Lambda }}{}^{i}=\frac{x^\mu }{x^\lambda }d_\mu y_\mathrm{\Lambda }^i,$$
(3)
where $`\mathrm{\Lambda }=(\lambda _k\mathrm{}\lambda _1)`$, $`|\mathrm{\Lambda }|=k`$, is a multi-index, $`\lambda +\mathrm{\Lambda }`$ is the multi-index $`(\lambda \lambda _k\mathrm{}\lambda _1)`$ and $`d_\lambda `$ are the total derivatives
$`d_\lambda =_\lambda +{\displaystyle \underset{|\mathrm{\Lambda }|=0}{}}y_{\lambda +\mathrm{\Lambda }}^i_i^\mathrm{\Lambda }.`$
The differential calculus on $`J^{\mathrm{}}Y`$ can be introduced as operations on the $`R`$-ring $`𝒬_{\mathrm{}}^0`$ of locally pull-back functions on $`J^{\mathrm{}}Y`$. A real function $`f`$ on $`J^{\mathrm{}}Y`$ is called so if, for each point $`qJ^{\mathrm{}}Y`$, there is a neighbourhood $`U_q`$ such that $`f|_{U_q}`$ is the pull-back of a smooth function on some finite-order jet manifold $`J^kY`$ with respect to the surjection $`\pi _k^{\mathrm{}}`$. It should be emphasized that the paracompact space $`J^{\mathrm{}}Y`$ admits the partition of unity performed by elements of $`𝒬_{\mathrm{}}^0`$ . The difficulty lies in the geometric interpretation of derivations of the $`R`$-ring $`𝒬_{\mathrm{}}^0`$ as vector fields on the Fréchet manifold $`J^{\mathrm{}}Y`$ and their dual as differential forms on $`J^{\mathrm{}}Y`$ .
Therefore, one usually considers the subring $`𝒪_{\mathrm{}}^0`$ of the ring $`𝒬_{\mathrm{}}^0`$ which consists of the pull-back onto $`J^{\mathrm{}}Y`$ of smooth functions on finite-order jet spaces. The Lie algebra of derivations of $`𝒪_{\mathrm{}}^0`$ is isomorphic to the projective limit onto $`J^{\mathrm{}}Y`$ of the Lie algebras of projectable vector fields on finite-order jet manifolds. The associated algebra of differential forms is introduced as the direct limit $`(𝒪_{\mathrm{}}^{},\pi _k^{\mathrm{}})`$ of the direct system
$$𝒪^{}(X)\stackrel{\pi ^{}}{}𝒪^{}(Y)\stackrel{\pi _0^1^{}}{}𝒪_1^{}\stackrel{\pi _1^2^{}}{}\mathrm{}\stackrel{\pi _{r1}^r^{}}{}𝒪_r^{}\mathrm{}$$
(4)
of differential $`R`$-algebras $`𝒪_r^{}`$ of exterior forms on finite-order jet manifolds $`J^rY`$. This direct limit exists in the category of $`R`$-modules, and the direct limits of the familiar operations on exterior forms make $`𝒪_{\mathrm{}}^{}`$ a differential exterior $`R`$-algebra. This algebra consists of all exterior forms on finite-order jet manifolds modulo the pull-back identification. Therefore, one usually thinks of elements of $`𝒪_{\mathrm{}}^{}`$ as being the pull-back onto $`J^{\mathrm{}}Y`$ of exterior forms on finite-order jet manifolds. Being restricted to a coordinate chart $`(\pi ^{\mathrm{}})^1(U_X)`$ on $`J^{\mathrm{}}Y`$, elements of $`𝒪_{\mathrm{}}^{}`$ can be written in the familiar coordinate form, where basic forms $`\{dx^\lambda \}`$ and contact 1-forms $`\{\theta _\mathrm{\Lambda }^i=dy_\mathrm{\Lambda }^iy_{\lambda +\mathrm{\Lambda }}^idx^\lambda \}`$ provide the local generators of the algebra $`𝒪_{\mathrm{}}^{}`$. There is the canonical splitting of the space of $`m`$-forms
$`𝒪_{\mathrm{}}^m=𝒪_{\mathrm{}}^{0,m}𝒪_{\mathrm{}}^{1,m1}\mathrm{}𝒪_{\mathrm{}}^{m,0}`$
into spaces $`𝒪_{\mathrm{}}^{k,mk}`$ of $`k`$-contact forms. Accordingly, the exterior differential on $`𝒪_{\mathrm{}}^{}`$ is decomposed into the sum $`d=d_H+d_V`$ of horizontal and vertical differentials
$`d_H:𝒪_{\mathrm{}}^{k,s}𝒪_{\mathrm{}}^{k,s+1},d_H(\varphi )=dx^\lambda d_\lambda (\varphi ),\varphi 𝒪_{\mathrm{}}^{},`$
$`d_V:𝒪_{\mathrm{}}^{k,s}𝒪_{\mathrm{}}^{k+1,s},d_V(\varphi )=\theta _\mathrm{\Lambda }^i_\mathrm{\Lambda }^i\varphi ,`$
which obey the nilpotency rule
$$d_Hd_H=0,d_Vd_V=0,d_Vd_H+d_Hd_V=0.$$
(5)
In studying the algebra $`𝒪_{\mathrm{}}^{}`$ of pull-back exterior forms on $`J^{\mathrm{}}Y`$, the key point of lies in the fact that the infinite-order De Rham complex of these forms
$$0R𝒪_{\mathrm{}}^0\stackrel{d}{}𝒪_{\mathrm{}}^1\stackrel{d}{}\mathrm{}$$
(6)
is the direct limit of the De Rham complexes of exterior forms on finite-order jet manifolds. Then, as was repeatedly proved, the cohomology groups $`H^{}(𝒪_{\mathrm{}}^{})`$ of the complex (6) are equal to the De Rham cohomology groups $`H^{}(Y)`$ of the fibre bundle $`Y`$ . This fact enables one to say something on the topological obstruction to the exactness of the (infinite-order) variational complex in the calculus of variations in field theory \[6,7,10-12\]. At the same time, the $`d_H`$-cohomology of the horizontal complex (1) of pull-back exterior forms on $`J^{\mathrm{}}Y`$ remains unknown. The local exactness of this complex only has been repeatedly proved (see, e.g., ). If a fibre bundle $`YX`$ admits a global section, there is also a monomorphism of the De Rham cohomology groups $`H^{}(X)`$ of the base $`X`$ to the cohomology groups of the complex (1) .
Here, we show that the problem of cohomology of the horizontal complex has a comprehensive solution by enlarging the algebra $`𝒪_{\mathrm{}}^{}`$ to the algebra $`𝒬_{\mathrm{}}^{}`$ of the above mentioned locally pull-back exterior forms on $`J^{\mathrm{}}Y`$. We introduce these forms in an algebraic way as global sections of the sheaf $`𝔔_{\mathrm{}}^{}`$ of differential algebras on $`J^{\mathrm{}}Y`$ which is the direct limit of the direct system
$$𝔒_X^{}\stackrel{\pi ^{}}{}𝔒_Y^{}\stackrel{\pi _0^1^{}}{}𝔒_1^{}\stackrel{\pi _1^2^{}}{}\mathrm{}\stackrel{\pi _{r1}^r^{}}{}𝔒_r^{}\mathrm{}$$
(7)
of sheaves of exterior forms on finite-order jet manifolds $`J^rY`$. As a consequence, we have the exact sequence of sheaves
$$0R𝔔_{\mathrm{}}^0\stackrel{d_H}{}𝔔_{\mathrm{}}^{0,1}\stackrel{d_H}{}\mathrm{}\stackrel{d_H}{}𝔔_{\mathrm{}}^{0,n}$$
(8)
of horizontal forms on $`J^{\mathrm{}}Y`$ and the corresponding complex of $`𝒬_{\mathrm{}}^0`$-modules of their global sections
$$0R𝒬_{\mathrm{}}^0\stackrel{d_H}{}𝒬_{\mathrm{}}^{0,1}\stackrel{d_H}{}\mathrm{}\stackrel{d_H}{}𝒬_{\mathrm{}}^{0,n}.$$
(9)
Since $`J^{\mathrm{}}Y`$ is paracompact and admits a partition of unity by elements of $`𝒬_{\mathrm{}}^0`$, all sheaves $`𝔔_{\mathrm{}}^{0,m}`$ in the exact sequence (8) are fine and, consequently, acyclic. Therefore, the well-known theorem on a resolution of a sheaf can be applied in order to obtain the cohomology groups of the horizontal complex (9).
From the physical viewpoint, an extension of the class of exterior forms to $`𝒬_{\mathrm{}}^{}`$ enables us to concern effective field theories whose Lagrangians involve derivatives of arbitrary high order
Here, we study $`d`$-, $`d_V`$\- and $`d_H`$-cohomology of the horizontal complex (9) on the infinite-order jet space $`J^{\mathrm{}}Y`$ of an affine bundle $`YX`$. Note that affine bundles provide a standard framework in quantum field theory because almost all existent quantization schemes deal with linear and affine quantities. Moreover, the De Rham cohomology groups of an affine bundle $`YX`$ are equal to those of its base $`X`$. Therefore, as we will see, the obstruction to the exactness of the horizontal complex (9) lies only in exterior forms on $`X`$. Since the BRST operator $`𝐬`$ eliminate these forms, the global descent equations can be constructed though their right-hand sides are not equal to zero.
Moreover, we can restrict our consideration to vector bundles $`YX`$ without loss of generality as follows. Let $`YX`$ be a smooth affine bundle modelled over a smooth vector bundle $`\overline{Y}X`$. A glance at the transformation law (3) shows that $`J^{\mathrm{}}YX`$ is an affine topological bundle modelled on the vector bundle $`J^{\mathrm{}}\overline{Y}X`$. This affine bundle admits a global section $`J^{\mathrm{}}s`$ which is the infinite-order jet prolongation of a global section $`s`$ of $`YX`$. With $`J^{\mathrm{}}s`$, we have a homeomorphism
$`\widehat{s}_{\mathrm{}}:J^{\mathrm{}}Yqq(J^{\mathrm{}}s)(\pi ^{\mathrm{}}(q))J^{\mathrm{}}\overline{Y}`$
of the topological spaces $`J^{\mathrm{}}Y`$ and $`J^{\mathrm{}}\overline{Y}`$, together with an exterior algebra isomorphism $`\widehat{s}_{\mathrm{}}^{}:\overline{𝒪}_{\mathrm{}}^{}𝒪_{\mathrm{}}^{}`$. Moreover, it is readily observed that the pull-back morphism $`\widehat{s}_{\mathrm{}}^{}`$ commutes with the differentials $`d`$, $`d_V`$ and $`d_H`$. Therefore, the differential algebras $`𝒬_{\mathrm{}}^{}`$ and $`\overline{𝒬}_{\mathrm{}}^{}`$ have the same $`d`$-, $`d_V`$\- and $`d_H`$-cohomology.
Given a smooth vector bundle $`YX`$, we will show the following.
* 1. The De Rham cohomology groups of the differential algebra $`𝒬_{\mathrm{}}^{}`$ are isomorphic to those of the base $`X`$.
2. Its $`d_V`$-cohomology groups are trivial.
3. The $`d_H`$-cohomology groups of contact elements of the algebra $`𝒬_{\mathrm{}}^{}`$ are trivial.
4. The $`d_H`$-cohomology groups of its horizontal elements (i.e., cohomology of the horizontal complex (9)) coincide with the De Rham cohomology groups of the base $`X`$.
Note that the results (i) and (ii) are also true for the differential algebra $`𝒪_{\mathrm{}}^{}`$. The result (iii) takes place for an arbitrary smooth fibre bundle $`YX`$, and recovers that in Refs. , obtained by means of the Mayer-Vietoris sequence.
## 2 Differential algebra $`𝒬_{\mathrm{}}^{}`$
Throughout the paper, smooth manifolds are real, finite-dimensional, Hausdorff, paracompact, and connected.
Given the surjective inverse system (2), we have the direct system (7) of ringed spaces $`(J^kY,𝔒_k^{})`$ whose structure sheaves $`𝔒_k^{}`$ are sheaves of differential $`R`$-algebras of exterior forms on finite-order jet manifolds $`J^kY`$, and $`\pi _{r1}^r^{}`$ are the pull-back morphisms. Throughout, we follow the terminology of Ref. where by a sheaf is meant a sheaf bundle. The direct system (7) admits a direct limit $`𝔔_{\mathrm{}}^{}`$ which is a sheaf of differential exterior $`R`$-algebras on the infinite-order jet space $`J^{\mathrm{}}Y`$. This direct limit exists in the category of sheaves of $`R`$-modules, and the direct limits of the familiar operations on exterior forms provide $`𝔔_{\mathrm{}}^{}`$ with a differential exterior algebra structure .
Accordingly, we have the direct system (4) of the structure algebras $`𝒪_k^{}=\mathrm{\Gamma }(J^kY,𝔒_k^{})`$ of global sections of sheaves $`𝔒_k^{}`$, i.e., $`𝒪_k^{}`$ are differential $`R`$-algebras of (global) exterior forms on finite-order jet manifolds $`J^kY`$. As was mentioned above, the direct limit of (4) is a differential exterior $`R`$-algebra $`(𝒪_{\mathrm{}}^{},\pi _k^{\mathrm{}})`$, isomorphic to the algebra of all exterior forms on finite-order jet manifolds modulo the pull-back identification.
The crucial point is that the limit $`𝒪_{\mathrm{}}^{}`$ of the direct system (4) of structure algebras of sheaves $`𝔒_k^{}`$ fails to coincide with the structure algebra $`𝒬_{\mathrm{}}^{}=\mathrm{\Gamma }(J^{\mathrm{}}Y,𝔔_{\mathrm{}}^{})`$ of the limit $`𝔔_{\mathrm{}}^{}`$ of the direct system (7) of these sheaves. The sheaf $`𝔔_{\mathrm{}}^{}`$, by definition, is the sheaf of germs of local exterior forms on finite-order jet manifolds. These local forms constitute a presheaf $`𝔒_{\mathrm{}}^{}`$ from which the sheaf $`𝔔_{\mathrm{}}^{}`$ is constructed. It means that, given a section $`\varphi \mathrm{\Gamma }(𝔔_{\mathrm{}}^{})`$ of $`𝔔_{\mathrm{}}^{}`$ over an open subset $`UJ^{\mathrm{}}Y`$ and any point $`qU`$, there exists a neighbourhood $`U_q`$ of $`q`$ such that $`\varphi |_{U_q}`$ is the pull-back of a local exterior form on some finite-order jet manifold. However, $`𝔒_{\mathrm{}}^{}`$ does not coincide with the canonical presheaf $`\mathrm{\Gamma }(𝔔_{\mathrm{}}^{})`$ of sections of the sheaf $`𝔔_{\mathrm{}}^{}`$.
In particular, the $`R`$-ring $`𝒬_{\mathrm{}}^0`$ is isomorphic to the above mentioned ring of real locally pull-back functions on $`J^{\mathrm{}}Y`$. Indeed, any element of $`𝒬_{\mathrm{}}^0`$ defines obviously such a function on $`J^{\mathrm{}}Y`$. Conversely, the germs of any locally pull-back function $`f`$ on $`J^{\mathrm{}}Y`$ belong to the sheaf $`𝔔_{\mathrm{}}^{}`$, i.e., $`f`$ is a section of $`𝔔_{\mathrm{}}^{}`$, and different such functions $`f`$ and $`f^{}`$ are different sections of $`𝔔_{\mathrm{}}^{}`$.
There are obvious monomorphisms of algebras $`𝒪_{\mathrm{}}^{}𝒬_{\mathrm{}}^{}`$ and presheaves $`𝔒_{\mathrm{}}^{}\mathrm{\Gamma }(𝔔_{\mathrm{}}^{})`$.
For short, we agree to call $`𝔔_{\mathrm{}}^{}`$ (resp. $`𝒬_{\mathrm{}}^{}`$) the sheaf (resp. algebra) of locally pull-back exterior forms on $`J^{\mathrm{}}Y`$. The exterior algebra operations and differentials $`d`$, $`d_V`$, $`d_H`$ are defined on $`𝒬_{\mathrm{}}^{}`$ just as on $`𝒪_{\mathrm{}}^{}`$. At the same time, it should be emphasized again that elements of the differential algebras $`𝒪_{\mathrm{}}^{}`$ and $`𝒬_{\mathrm{}}^{}`$ are not differential forms on $`J^{\mathrm{}}Y`$ in a rigorous sense. Therefore, the standard theorems, e.g., the well-known De Rham theorem (, Theorem 2.12.3) can not be applied automatically to these differential algebras.
## 3 De Rham cohomology
Let $`YX`$ be an arbitrary smooth fibre bundle. We consider the complex of sheaves of $`𝒬_{\mathrm{}}^0`$-modules
$$0R𝔔_{\mathrm{}}^0\stackrel{d}{}𝔔_{\mathrm{}}^1\stackrel{d}{}\mathrm{}$$
(10)
on the infinite-order jet space $`J^{\mathrm{}}Y`$ and the corresponding infinite-order De Rham complex
$$0R𝒬_{\mathrm{}}^0\stackrel{d}{}𝒬_{\mathrm{}}^1\stackrel{d}{}\mathrm{}$$
(11)
of locally pull-back exterior forms on $`J^{\mathrm{}}Y`$.
Since locally pull-back exterior forms fulfill the Poincaré lemma, the complex of sheaves (10) is exact. Since the paracompact space $`J^{\mathrm{}}Y`$ admits a partition of unity performed by elements of $`𝒬_{\mathrm{}}^0`$ , the sheaves $`𝔔_{\mathrm{}}^r`$ of $`𝒬_{\mathrm{}}^0`$-modules are fine for all $`r0`$ . Then they are acyclic, i.e., the cohomology groups $`H^{>0}(J^{\mathrm{}}Y,𝔔_{\mathrm{}}^r)`$ of the paracompact space $`J^{\mathrm{}}Y`$ with coefficients in the sheaf $`𝔔_{\mathrm{}}^r`$ vanish . Consequently, the exact sequence (10) is a fine resolution of the constant sheaf $`R`$ of germs of local constant real functions on $`J^{\mathrm{}}Y`$. Then the well-known (generalized De Rham) theorem on a resolution of a sheaf on a paracompact space (, Theorem 2.12.1) can be called into play in order to find the cohomology groups of the infinite-order De Rham complex complex (6).
In accordance with this theorem, we have an isomorphism
$$H^{}(𝒬_{\mathrm{}}^{})=H^{}(J^{\mathrm{}}Y,R)$$
(12)
of the De Rham cohomology groups $`H^{}(𝒬_{\mathrm{}}^{})`$ of the differential algebra $`𝒬_{\mathrm{}}^{}`$ and the cohomology groups $`H^{}(J^{\mathrm{}}Y,R)`$ of the infinite-order jet space $`J^{\mathrm{}}Y`$ with coefficients in the constant sheaf $`R`$. In the case of a vector bundle $`YX`$, we can say something more.
LEMMA 1. If $`YX`$ is a vector bundle, there is an isomorphism
$$H^{}(J^{\mathrm{}}Y,R)=H^{}(X,R)$$
(13)
of cohomology groups of the infinite-order jet space $`J^{\mathrm{}}Y`$ with coefficients in the constant sheaf $`R`$ and those $`H^{}(X,R)`$ of the base $`X`$.
Proof. The cohomology groups with coefficient in the constant sheaf $`R`$ on homotopic paracompact topological spaces are isomorphic . If $`YX`$ is a vector bundle, its base $`X`$ is a strong deformation retract of the infinite-order jet space $`J^{\mathrm{}}Y`$. To show this, let us consider the map
$`[0,1]\times J^{\mathrm{}}Y(t;x^\lambda ,y_\mathrm{\Lambda }^i)(x^\lambda ,ty_\mathrm{\Lambda }^i)J^{\mathrm{}}Y.`$
A glance at the transition functions (3) shows that, given in the coordinate form, this map is well-defined if $`YX`$ is a vector bundle. It is a desired homotopy from $`J^{\mathrm{}}Y`$ to the base $`X`$ which is identified with its image under the global zero section of the vector bundle $`J^{\mathrm{}}YX`$. $`\mathrm{}`$
Combining the isomorphisms of cohomology groups (12), (13), and the well-known isomorphism $`H^{}(X,R)=H^{}(X)`$, we come to the manifested isomorphism
$`H^{}(𝒬_{\mathrm{}}^{})=H^{}(X)`$
of the De Rham cohomology groups $`H^{}(𝒬_{\mathrm{}}^{})`$ of the differential algebra $`𝒬_{\mathrm{}}^{}`$ of locally pull-back forms on the infinite-order jet space $`J^{\mathrm{}}Y`$ of a vector bundle $`YX`$ to the De Rham cohomology groups $`H^{}(X)`$ of the base $`X`$.
It follows that any closed form $`\varphi 𝒬_{\mathrm{}}^{}`$ on $`J^{\mathrm{}}Y`$ is decomposed into the sum $`\varphi =\phi +d\xi `$ where $`\phi 𝒪^{}(X)`$ is a closed form on $`X`$.
## 4 Cohomology of $`d_V`$
Due to the nilpotency rule (5), the vertical and horizontal differentials $`d_V`$ and $`d_H`$ on the differential exterior algebra $`𝒬_{\mathrm{}}^{}`$ define the bicomplex
$$\begin{array}{ccccccccccccccccccc}& & & & \hfill _{d_V}& \text{}\hfill & & \hfill _{d_V}& \text{}\hfill & & & \hfill _{d_V}& \text{}\hfill & & & \hfill _{d_V}& \text{}\hfill & & \\ & & 0& & & 𝒪_{\mathrm{}}^{k,0}\hfill & \stackrel{d_H}{}& & 𝒬_{\mathrm{}}^{k,1}\hfill & \stackrel{d_H}{}& \mathrm{}& & 𝒬_{\mathrm{}}^{k,m}\hfill & \stackrel{d_H}{}& \mathrm{}& & 𝒬_{\mathrm{}}^{k,n}\hfill & & \\ & & & & & \mathrm{}\hfill & & & \mathrm{}\hfill & & & & \mathrm{}\hfill & & & & \mathrm{}\hfill & & \\ & & & & \hfill _{d_V}& \text{}\hfill & & \hfill _{d_V}& \text{}\hfill & & & \hfill _{d_V}& \text{}\hfill & & & \hfill _{d_V}& \text{}\hfill & & \\ 0& & R& & & 𝒬_{\mathrm{}}^0\hfill & \stackrel{d_H}{}& & 𝒬_{\mathrm{}}^{0,1}\hfill & \stackrel{d_H}{}& \mathrm{}& & 𝒬_{\mathrm{}}^{0,m}\hfill & \stackrel{d_H}{}& \mathrm{}& & 𝒬_{\mathrm{}}^{0,n}\hfill & & \\ & & & & \hfill _\pi ^{\mathrm{}}& \text{}\hfill & & \hfill _\pi ^{\mathrm{}}& \text{}\hfill & & & \hfill _\pi ^{\mathrm{}}& \text{}\hfill & & & \hfill _\pi ^{\mathrm{}}& \text{}\hfill & & \\ 0& & R& & & 𝒬^0(X)\hfill & \stackrel{d}{}& & 𝒬^1(X)\hfill & \stackrel{d}{}& \mathrm{}& & 𝒬^m(X)\hfill & \stackrel{d}{}& \mathrm{}& & 𝒬^n(X)\hfill & \stackrel{d}{}& 0\\ & & & & & \text{}\hfill & & & \text{}\hfill & & & & \text{}\hfill & & & & \text{}\hfill & & \\ & & & & & 0\hfill & & & 0\hfill & & & & 0\hfill & & & & 0\hfill & & \end{array}$$
(14)
\[5-8,10,11,13,14\]. The rows and columns of these bicomplex are horizontal and vertical complexes. Let us consider a vertical one
$$0𝒬^m(X)\stackrel{\pi ^{\mathrm{}}}{}𝒬_{\mathrm{}}^{0,m}\stackrel{d_V}{}\mathrm{}\stackrel{d_V}{}𝒬_{\mathrm{}}^{k,m}\stackrel{d_V}{}\mathrm{},mn.$$
(15)
PROPOSITION 2. If $`YX`$ is a vector bundle, then the vertical complex (15) is exact.
Proof. Local exactness of a vertical complex on a coordinate chart $`((\pi ^{\mathrm{}})^1(U_X);x^\lambda ,y_\mathrm{\Lambda }^i)`$, $`0|\mathrm{\Lambda }|,`$ on $`J^{\mathrm{}}Y`$ follows from a version of the Poincaré lemma with parameters (see, e.g., ). We have the the corresponding homotopy operator
$`\sigma ={\displaystyle _0^1}t^k[\overline{y}\varphi (x^\lambda ,ty_\lambda ^i)]dt,\varphi 𝒬^{k,m}_{\mathrm{}},`$
where $`\overline{y}=y_\mathrm{\Lambda }^i_i^\mathrm{\Lambda }`$. Since $`YX`$ is a vector bundle, it is readily observed that this homotopy operator is globally defined on $`J^{\mathrm{}}Y`$, and so is the exterior form $`\sigma `$. $`\mathrm{}`$
It means that any $`d_V`$-closed form $`\varphi 𝒬_{\mathrm{}}^{}`$ is the sum $`\varphi =\phi +d_V\xi `$ of a $`d_V`$-exact form and an exterior form $`\phi `$ on $`X`$.
Of course, $`d_V`$-cohomology of the bicomplex (14) is not trivial if the typical fibre of the fibre bundle $`YX`$ is not contractible.
## 5 Cohomology of $`d_H`$
Turn now to the rows of the bicomplex (14) (excluding the bottom one which is obviously the De Rham complex on the base $`X`$). The algebraic Poincaré lemma (see, e.g., ) is obviously extended to elements of $`𝒬_{\mathrm{}}^{}`$.
PROPOSITION 3. If $`YX`$ is a contractible fibre bundle $`R^{n+m}R^n`$, the rows of the bicomplex (14) are exact, i.e., they are always locally exact.
It follows that the the corresponding complexes of sheaves of contact forms
$$0𝔔_{\mathrm{}}^{k,0}\stackrel{d_H}{}𝔔_{\mathrm{}}^{k,1}\stackrel{d_H}{}\mathrm{}\stackrel{d_H}{}𝔔_{\mathrm{}}^{k,n},k>0,$$
(16)
and the above mentioned horizontal complex
$$0R𝔔_{\mathrm{}}^0\stackrel{d_H}{}𝔔_{\mathrm{}}^{0,1}\stackrel{d_H}{}\mathrm{}\stackrel{d_H}{}𝔔_{\mathrm{}}^{0,n}$$
(17)
are exact. Recall that, since $`J^{\mathrm{}}Y`$ is paracompact and admits a partition of unity by elements of $`𝒬_{\mathrm{}}^0`$, all sheaves except the constant sheaf $`R`$ in the complexes (16), (17) are fine. However, the exact sequences (16) and (17) fail to be fine resolutions of the sheaves $`𝔔_{\mathrm{}}^{k,0}`$ and $`R`$, respectively, because of their last terms. At the same time, following directly the proof of the above mentioned generalized De Rham theorem (, Theorem 2.12.1) till these terms, one can show the following.
PROPOSITION 4. If $`YX`$ is an arbitrary smooth fibre bundle, then the cohomology groups $`H^r(k,d_H)`$, $`r<n`$, of the complex of $`𝒬_{\mathrm{}}^0`$-modules
$`0𝒬_{\mathrm{}}^{k,0}\stackrel{d_H}{}𝒬_{\mathrm{}}^{k,1}\stackrel{d_H}{}\mathrm{}\stackrel{d_H}{}𝒬_{\mathrm{}}^{k,n}`$
are isomorphic to the cohomology groups $`H^r(J^{\mathrm{}}Y,𝔔^{k,0})`$ of $`J^{\mathrm{}}Y`$ with coefficients in the sheaf $`𝔔^{k,0}`$ and, consequently, are trivial because the sheaf $`𝔔^{k,0}`$ is fine.
PROPOSITION 5. If $`YX`$ is an arbitrary smooth fibre bundle, the cohomology groups $`H^r(d_H)`$, $`r<n`$, of the horizontal complex (9) are isomorphic to the cohomology groups $`H^r(J^{\mathrm{}}Y,R)`$ of $`J^{\mathrm{}}Y`$ with coefficients in the constant sheaf $`R`$.
Note that one can also study the exact sequence of presheaves
$`0R_{\mathrm{}}𝔒_{\mathrm{}}^0\stackrel{d}{}𝔒_{\mathrm{}}^1\stackrel{d}{}\mathrm{},`$
but comes again to the results of Propositions 5, 5. Because $`J^{\mathrm{}}Y`$ is paracompact, the cohomology groups $`H^{}(J^{\mathrm{}}Y,𝔔_{\mathrm{}}^{})`$ of $`J^{\mathrm{}}Y`$ with coefficients in the sheaf $`𝔔_{\mathrm{}}^{}`$ and those $`H^{}(J^{\mathrm{}}Y,𝔒_{\mathrm{}}^{})`$ with coefficients in the presheaf $`𝔒_{\mathrm{}}^{}`$ are isomorphic. It follows that the cohomology group $`H^0(J^{\mathrm{}}Y,𝔒_{\mathrm{}}^{})`$ of the presheaf $`𝒪_{\mathrm{}}^{}`$ is isomorphic to the $`R`$-module $`𝒬_{\mathrm{}}^{}=H^0(J^{\mathrm{}}Y,𝔔_{\mathrm{}}^{})`$, but not $`𝒪_{\mathrm{}}^{}`$.
If $`YX`$ is a vector bundle, Lemma 3 and Proposition 5 lead to the manifested isomorphism
$`H^r(d_H)=H^r(X,R)=H^r(X),r<n,`$
of the $`d_H`$-cohomology groups $`H^{<n}(d_H)`$ of the horizontal complex (9) to the De Rham cohomology groups $`H^{<n}(X)`$ of the base $`X`$. Then, combining Propositions 5 and 5, we conclude that any $`d_H`$-closed form $`\varphi 𝒬_{\mathrm{}}^{}`$ on the infinite-order jet manifold $`J^{\mathrm{}}Y`$ is decomposed into the sum
$$\varphi =\phi +d_H\xi $$
(18)
where $`\phi 𝒪^{}(X)`$ is a closed form on $`X`$.
Turn to outcomes of this result to BRST theory. Since $`𝐬\phi =0`$, the decomposition (18) prevents one from the topological obstruction to definition of global descent equations in BRST theory on vector (and affine) bundles.
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# Fluctuations and the existence of potential in dissipative semiclassical systems
## I Introduction
One of the important issues in nonequilibrium phenomena in the macroscopic nonlinear systems is to understand the interplay of nonlinearity of the system and the fluctuations of its environment. The problem is fairly general in the context of chemical reactions , breakdown of electronic devices , phase transitions etc. The essential description of the physical situation rests on the Fokker-Planck equations for the probability distribution functions of the relevant variables of the dynamics. In the weak noise limit the fluctuations have been described by appropriate auxiliary Hamiltonian or path integral methods. The theoretical results have been corroborated by remarkable experiments on fluctuations using analogue electronic circuits , which allow the phase space trajectories of fluctuations to be observed directly in a precise manner. These studies have enriched our understanding in several theoretical issues, e. g. , the symmetry between the growth and the decay of classical fluctuations in equilibrium and its breakdown under nonequilibrium conditions , the existence of a nonequilibrium potential of a dissipative system etc.
It is the purpose of this paper to extend the theory to the semiclassical context. The quantization of the system itself adds a new dimension to the interplay of nonlinearity and stochasticity in a dissipative system. To make a fair comparison with classical theory we adopt the Wigner’s phase space distribution function of c-number variables. The weak noise limit can then be appropriately employed to develop an auxiliary Hamiltonian formulation at the semiclassical level in terms of these phase space variables. This allows us to realize the existence of an optimal force of purely quantum origin derivable from the fluctuating field and relate it to the momentum of the auxiliary Hamiltonian. The quantum correction also makes its presence felt in the growth and decay of fluctuations of thermally equilibrated semiclassical systems keeping the symmetry preserved.
The outline of the paper is as follows: In Sec. II we introduce the general aspects of dynamics of dissipative quantum system in terms of the Wigner phase space function. In Sec. III we consider the weak noise and semiclassical limit under overdamped condition and take resort to the well-known auxiliary Hamiltonian description. The quantum part of the optimal force is then explicitly derived. The symmetry between the growth and decay of fluctuations in a thermalized quantum system is discussed in Sec. IV. The existence of a semiclassical contribution to the potential in a dissipative system is then shown in Sec. V. The paper is concluded in Sec. VI.
## II Quantum dynamics of a dissipative system
We consider a dynamical system characterized by a potential $`V(x)`$ coupled to an environment. Evolution of such an open quantum system has been studied over the last several decades under a variety of reasonable assumptions . Specifically interesting is the semiclassical limit of an Ohmic environment. The dissipative time evolution of the Wigner distribution function $`W(x,p,t)`$ for the system with unit mass ($`m=1`$) under the potential $`V(x)`$ can be described by
$$\frac{W}{t}=p\frac{W}{x}+\frac{V}{x}\frac{W}{p}+ϵ\underset{n1}{}\frac{\mathrm{}^{2n}(1)^n}{2^{2n}(2n+1)!}\frac{^{2n+1}V}{x^{2n+1}}\frac{^{2n+1}W}{p^{2n+1}}+\gamma \frac{pW}{p}+D\frac{^2W}{p^2},$$
(1)
where $`\gamma `$ and $`D`$ are the dissipation constant and the diffusion coefficient, respectively. $`x`$ and $`p`$ are c-number co-ordinate and momentum variables. The drift term is a direct consequence of the existence of $`\gamma `$-dependent term in the imaginary part of the exponent in the expression for the propagator for the density operator in the Feynman-Veron theory and has been shown to be responsible for appearance of a damping force in the classical equation of motion for the Brownian particle to ensure quantum-classical correspondence. $`\gamma `$ and D are related by the fluctuation-dissipation relation, $`D=\frac{\gamma }{2}\mathrm{}\omega \mathrm{coth}\frac{\mathrm{}\omega }{2k_bT}`$ (in the semiclassical limit $`D=\gamma kT`$ ). $`\omega `$ is the renormalised linear frequency of the nonlinear system. The quantum correction to classical Liouville motion is contained in the $`\mathrm{}`$-containing terms in the sum. $`ϵ`$ is a parameter (whose value is 1) which is kept in the equation for bookkeeping the Wigner correction term in our further analysis. We put $`ϵ=1`$ at the end of calculation.
Eq.(1) had been used earlier in several occasions. For example, Zurek and Paz and others have studied some interesting aspects of quantum-classical correspondence in relation to decoherence and chaos. Based on this equation and its variant chaotic dissipative systems has been studied. . The equation also yields the simplest leading order quantum correction term to classical Kramers’ rate . The primary reason for choosing Eq.(3) as our starting point is that it reaches the correct classical limit when $`\mathrm{}0`$ so that $`D`$ becomes a thermal diffusion coefficient ($`\gamma kT`$) in the high temperature limit and the Wigner function reduces to the corresponding classical phase space distribution function and we recover the Kramers’ equation which describes classical Brownian motion of a particle in phase space.
## III Weak noise and semiclassical limit of quantum dissipative dynamics:
The weak noise limit of a dissipative system within a semiclassical description can be conveniently described by a “WKB-like” ansatz (we refer to “WKB-like” since we are considering more than one dimension. Traditionally WKB refer to one dimension only) of the Eq.(1) for the Wigner function of the form
$$W(x,p,t)=Z(x,t)\mathrm{exp}(\frac{xp}{\mathrm{}})\mathrm{exp}(\frac{s}{D_1}).$$
(2)
where $`D_1=\frac{D}{\omega ^2}`$. The weak noise limit is defined as $`D_10`$ and semiclassical limit refers to ($`\mathrm{}0`$). $`Z(x,t)`$ is a prefactor and $`s(x,p,t)`$ is a classical action which is a function of c-number variables $`x`$ and $`p`$ , satisfying the following Hamilton-Jacobi equation,
$`{\displaystyle \frac{s}{t}}+p{\displaystyle \frac{s}{x}}V^{}{\displaystyle \frac{s}{p}}\gamma p{\displaystyle \frac{s}{p}}+\omega ^2({\displaystyle \frac{s}{p}})^2`$ (3)
$`+ϵ{\displaystyle \underset{n1}{}}{\displaystyle \frac{x^{2n}(1)^{3n+1}}{2^{2n}(2n)!}}{\displaystyle \frac{^{2n+1}V}{x^{2n+1}}}{\displaystyle \frac{s}{p}}=0.`$ (4)
The derivation of Eq.(3) is based on the following consideration. Since in the weak noise limit $`D_1`$ is the relevant small parameter one obtains with ansatz (2) in leading order a term proportional to $`\mathrm{}^{2n}\left(\frac{1}{D_1}\frac{s}{p}\right)^{2n+1}`$ which is not balanced by any other term of the same order $`D_{1}^{}{}_{}{}^{(2n+1)}`$. This is because the highest derivative in Eq.(1) does not have a factor scaling with the corresponding power of $`D_1`$. The successive terms next to the leading order are also $`\mathrm{}`$-containing terms. All of these terms vanishes in the semiclassical limit ($`\mathrm{}0`$). Therefore the leading order term that remains gives rise to Eq.(3). It is thus obvious that an ansatz (2) with $`\mathrm{}`$ finite is not feasible. The semiclassical limit $`\mathrm{}0`$ is a necessary requirement for the validity of ansatz (2). The above equation can be solved by integrating the Hamiltonian equations of motion,
$`\dot{x}`$ $`=`$ $`p`$ (5)
$`\dot{X}`$ $`=`$ $`P\gamma X`$ (6)
$`\dot{p}`$ $`=`$ $`V^{}+\gamma p2\omega ^2Xϵ{\displaystyle \underset{n1}{}}{\displaystyle \frac{(1)^{3n+1}x^{2n}}{2^{2n}(2n)!}}{\displaystyle \frac{^{2n+1}V}{x^{2n+1}}}`$ (7)
$`\dot{P}`$ $`=`$ $`\left[V^{\prime \prime }ϵ{\displaystyle \underset{n1}{}}{\displaystyle \frac{(1)^{3n+1}}{2^{2n}(2n)!}}{\displaystyle \frac{}{x}}\left(x^{2n}{\displaystyle \frac{^{2n+1}V}{x^{2n+1}}}\right)\right]X`$ (8)
which are derived from the following effective Hamiltonian $`H_{eff}`$
$`H_{eff}`$ $`=`$ $`pPV^{}X\gamma Xp+\omega ^2X^2`$ (10)
$`+ϵ{\displaystyle \underset{n1}{}}{\displaystyle \frac{x^{2n}(1)^{3n+1}}{2^{2n}(2n)!}}{\displaystyle \frac{^{2n+1}V}{x^{2n+1}}}X.`$
Here we have put $`\frac{s}{x}=P`$ and $`\frac{s}{p}=X`$. The introduction of additional degree-of-freedom by incorporating the auxiliary momentum (P) and co-ordinate (X) makes the system an effectively two-degree-of-freedom system. The origin of these two variables is the thermal fluctuations of the environment . The auxiliary Hamiltonian $`H_{eff}`$ is not to be confused with the microscopic Hamiltonian comprising the system, the bath and their coupling. Thus the phase space trajectories concern fluctuations of the c-number variables.
Under overdamped condition ($`\ddot{x}\gamma \dot{x};\ddot{X}\gamma \dot{X}`$) Eqs. (4) can be easily reduced to the following form ;
$`\dot{x}`$ $`=`$ $`K(x)+{\displaystyle \frac{2\omega ^2X}{\gamma }}`$ (11)
$`\dot{X}`$ $`=`$ $`{\displaystyle \frac{K(x)}{x}}X`$ (12)
where
$$K(x)=\frac{1}{\gamma }\left[V^{}+ϵ\underset{n1}{}\frac{(1)^{3n+1}x^{2n}}{2^{2n}(2n)!}\frac{^{2n+1}V}{x^{2n+1}}\right]$$
(13)
It is easy to recognize the quantity $`\frac{2\omega ^2X}{\gamma }`$ as a momentum and redefine it as $`p_r`$. Therefore Eqs(6) may be rewritten as
$`\dot{x}`$ $`=`$ $`K(x)+p_r`$ (14)
$`\dot{p_r}`$ $`=`$ $`{\displaystyle \frac{K(x)}{x}}p_r.`$ (15)
So the overdamped motion is described by the following effective Hamiltonian,
$$H_{od}=\frac{p_r^2}{2}+K(x)p_r.$$
(16)
The above auxiliary Hamiltonian description (8, 9) is isomorphic in form to that of Luchinsky and McClintock , who had considered an overdamped classical Brownian motion in a force field $`K(x)`$, driven by a weak white noise $`\zeta (t)`$ whose intensity $`D_11`$ as
$$\dot{x}=K(x)+\zeta (t),\zeta =0,\zeta (t)\zeta (0)=D_1\delta (t).$$
(17)
Equivalently the corresponding Fokker-Planck equation for the probability density $`P(x,t)`$ is
$$\frac{P(x,t)}{t}=\frac{K(x)}{x}P(x,t)+\frac{D_1}{2}\frac{^2P(x,t)}{x^2}.$$
(18)
The large fluctuations of scale $`\sqrt{D_1}`$ can therefore be treated in the weak noise limit $`D_10`$ by “WKB like” approximation of the Fokker-Planck equation (11) in the form
$$P(x,t)=y(x,t)\mathrm{exp}[\frac{\varphi (x,t)}{D_1}].$$
(19)
Here y(x, t) is the prefactor and $`\varphi (x,t)`$ is a “classical” action describing a Hamiltonian-Jacobi equation which can be solved by solving Hamilton’ s equation (8) with $`p_r=\frac{\varphi }{x}`$ as the momentum for auxiliary system.
The important distinguishing feature of the above description in which the system is treated semiclassically is the structure of $`K(x)`$ which is given by equation(7) and comprises of two terms ;
$$K(x)=K_{cl}+K_{semi}(x)$$
(20)
where
$$K_{cl}=\frac{V^{}(x)}{\gamma }$$
(21)
is derivable from purely classical potential $`V(x)`$ and $`K_{semi}(x)`$ does not explicitly involve $`\mathrm{}`$,
$$K_{semi}(x)=\frac{ϵ}{\gamma }\underset{n1}{}\frac{x^{2n}(1)^{3n+1}}{2^{2n}(2n)!}\frac{^{2n+1}V}{x^{2n+1}}$$
(22)
originates from the nonlinearity of the potential $`V(x)`$ and quantum nature of the system. The quantum contribution to $`K(x)`$ is therefore likely to influence both the fluctuational and the relaxational paths of the dynamics and is also responsible for the existence of a potential. Our objective is now to explore these aspects in the following two sections.
## IV Large fluctuations in equilibrated semiclassical systems
In the thermally equilibrated systems a typical large fluctuation of the variable $`x`$ implies a temporary departure from its stable state, $`x_s`$ to some remote state $`x_f`$. This is followed by a return to $`x_s`$ as a result of relaxation in the absence of fluctuations $`p_r`$. A nonzero value of $`p_r`$ which results from fluctuations due to surrounding drives the system away from $`x_s`$ along a set of trajectories which form the unstable invariant manifold and define the so called fluctuational paths. On the other hand the system relaxes along the relaxational return path to $`x_s`$ under the condition $`p_r=0`$, which form stable invariant manifold. The latter condition implies $`\dot{x}=K(x)`$. In each case the trajectory represents the optimal paths along which the system is expected to move with overwhelming probability. Luchinsky and McClintock have studied these paths in analog electronic circuits and demonstrated the growth and the decay of classical fluctuations in equilibrium. We extend this analysis to semiclassical domain using the same model potential,
$$V(x)=\frac{1}{4}x^4\frac{1}{2}x^2.$$
(23)
The quantum contribution to the growth and the decay of fluctuations can be understood by recognizing the $`K_{semi}(x)`$ term in the dynamics (8). In Fig.1 we compare both the deterministic fluctuational and relaxational (optimal) paths for quantum and classical thermally equilibrated systems. It is important to note that the maximum possible amplitude of large fluctuations is almost double for the quantum system compared to that for the corresponding the classical system. This is due to the addition of the nonlinear force term of quantum origin, $`K_{semi}(x)`$ in the c-number equation (8), which is shown to be derivable from the fluctuating field, and is related to the momentum of the auxiliary Hamiltonian.
Before leaving this section we make a brief remark on the thermally nonequilibrated systems. Since the detailed balance is not operative here, the optimal path to a given state not just the time-reversed dynamical path along which the system moves from this state to the stable state in absence of fluctuations $`p_r`$. Thus for the driven system, for example, the pattern of optimal path is generically different from that for the thermally equilibrated systems. It may display singularities whose topological manifestations as caustics, switching line, cusps etc have been thoroughly studied for classical systems. We believe that $`K_{semi}(x,t)`$ where t signifies the driving by a periodic force in $`V(x,t)`$ is likely to play an important role in their quantum counterparts.
## V Existence of a potential for dissipative semiclassical system
In a significant analysis Graham and coworkers had examined the condition for existence of potential for classical dissipative systems. We now extend this analysis to the present semiclassical context.
The general criterion for a dissipative dynamical system described by autonomous equations of the form
$$\dot{x}^\nu =K^\nu (x)$$
(24)
to have a potential $`\varphi (x)`$ with respect to $`Q^{\nu \mu }`$ (positive, semidefinite symmetric matrix, which are considered to be the matrix of transport coefficients) if there exists a single-valued continuously differentiable and globally defined function $`\varphi (x)`$, bounded from below which is stationary in the limit sets of Eq.(17) and which satisfies
$$K^\nu (x)=\frac{1}{2}Q^{\nu \mu }\frac{\varphi (x)}{x^\mu }+r^\nu (x)$$
(25)
with
$$r^\nu \frac{\varphi (x)}{x^\nu }=0.$$
(26)
Here for simplicity $`Q^{\nu \mu }`$ is assumed to be independent of $`x`$. The first and the second terms of Eq.(18) correspond to irreversible and reversible part, respectively.
The stochastic process $`x(D_1,t)`$ which involve Eq.(17) and a symmetric non-negative matrix $`Q^{\nu \mu }`$ is governed by the Fokker-Planck equation for probability distribution function $`P(x,t)`$
$$\frac{P(x,t)}{t}=\frac{}{x^\nu }K^\nu (x)P+\frac{D_1}{2}\frac{^2}{x^\nu x^\mu }Q^{\nu \mu }p.$$
(27)
For $`D_1=0`$ the above description reduces to (17). For $`D_10`$ the steady state distribution defines the function $`\varphi (x,t)`$ by
$$P(x,D_1,t\mathrm{})=N(D_1)\mathrm{exp}[\frac{\varphi (x,D_1)}{D_1}]$$
(28)
N is the normalization constant. If $`\varphi (x)=lim_{D_10}\varphi (x,D_1)`$ is a single-valued, continuously differentiable function bounded from below it satisfies
$$K^\nu (x)\frac{\varphi (x)}{x^\nu }+\frac{1}{2}Q^{\nu \mu }\frac{\varphi (x)}{x^\nu }\frac{\varphi (x)}{x^\mu }=0$$
(29)
Eq. (22) is equivalent to Eqs. (18, 19). Interpreting Eq.(22) as usual as a Hamilton-Jacobi equation by defining $`\varphi (x)`$ as action and $`\frac{\varphi }{x^\nu }=P_\nu `$ as the momentum conjugate to $`x_\nu `$, one can construct the auxiliary Hamiltonian
$$H(x,p)=\frac{1}{2}Q^{\nu \mu }p_\nu p_\mu +K^\nu (x)p_\nu .$$
(30)
Graham and co-workers have argued that a potential can exist with equation (17) if the above Hamiltonian is integrable for $`H=0`$, because the condition implies that there exist a smooth separatrix which connect smoothly the stable and unstable manifolds emanating from the hyperbolic fixed points of the dynamical system.
We now turn back to our dissipative semiclassical system described by Eqs(8) and (9) where $`K(x)`$ is defined by Eq.(7). Recognizing Eq.(9) as Eq.(23) for a one-degree-of-freedom system we identify
$`Q`$ $`=`$ $`1`$ (31)
$`K(x)`$ $`=`$ $`K_{cl}+K_{semi}.`$ (32)
The potential function $`\varphi (x)`$ can therefore be calculated from Eq.(9) with $`H=0`$ as (since $`p_r=\frac{\varphi (x)}{x}`$)
$$\varphi (x)=2K(x)𝑑x.$$
(33)
The above expression can be made more explicit if we make use of Eq.(7) in (25). We obtain
$$\varphi (x)=\varphi _{cl}(x)+\varphi _{semi}(x)$$
(34)
where
$$\varphi _{cl}=\frac{2}{\gamma }V^{}(x)𝑑x=\frac{2}{\gamma }V(x)$$
(35)
and
$$\varphi _{semi}(x)=\frac{2}{\gamma }\underset{n1}{}\frac{x^{2n}(1)^{3n+1}}{2^{2n}(2n)!}\frac{^{2n+1}V}{x^{2n+1}}𝑑x$$
(36)
The existence of a potential for a dissipative, semiclassical dynamical system is thus ascertained. The method essentially relies on a dynamical definition $`p_r`$ as a derivative of the potential $`\varphi (x)`$ in a system described by an overdamped quantum Markov process in the weak noise and semiclassical limit. As elaborated earlier in Sec.III $`p_r`$ has a statistical origin which drives the system away from its stable state $`x_s`$ to a preassigned remote state $`x_f`$ from which the system relaxes in absence of $`p_r`$. It is thus important to realize that both the dynamical and statistical notions are kept intact in the quantum treatment.
## VI Conclusions
Keeping in view of the quantum nature of the system in terms of the Wigner’s phase space function we examine the semiclassical dynamics of a dissipative system in an Ohmic environment. The weak noise limit of the stochastic process then allows us to capture the essential features of the dynamics within the framework of an auxiliary Hamiltonian description at the semiclassical level. Our results are summarized as follows :
(i) The Wigner’s quantum correction to classical Liouville equation gives rise to an optimal force in addition to usual the classical force term. This quantum optimal force is essentially a result of an interplay of nonlinearity of the system and the thermal fluctuations of its environment and is derivable in terms of an auxiliary Hamiltonian description.
(ii) This term is also responsible for modification of growth and decay of large fluctuations from equilibrium for the appropriately thermalized quantum system (compared to its classical counterpart). The symmetry of the fluctuational and the relaxational paths, signifying the detailed balance, however, as expected is kept preserved.
(iii) The quantum correction term implies the existence of a potential for the dissipative semiclassical system.
Since the fluctuational and the relaxational paths have been experimentally demonstrated as a part of physical reality by analogue experiments in the realm of large fluctuations, we believe that the essential modification of the integrable, nonintegrable and the singular topological features of the dynamics due to semiclassical correction might be relevant in several contexts. We hope to address some of these issues in a future communication.
## ACKNOWLEDGMENTS
B. C. Bag is indebted to the Council of Scientific and Industrial Research for a fellowship.
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# 1 Introduction
## 1 Introduction
The vacuum expectation values (VEV)s of local fields play an important role in quantum field theory (QFT) and in statistical mechanics . In statistical mechanics the VEVs determine the “generalized susceptibilities”, i.e. the linear response of a system to external fields. Furthermore, the VEVs provide all the information about correlation functions in QFT defined as a perturbed conformal field theory (CFT) that is not accessible through a direct calculation in conformal perturbation theory . A few years ago, some important progress was made in the calculation of such quantities in integrable (1+1) QFT. In ref. , an explicit expression for the VEVs of the exponential field in the sinh-Gordon and sine-Gordon models was proposed. In ref. it was shown that this result can be obtained using the “reflection amplitude” of the Liouville field theory. This method was also applied in the so-called Bullough-Dodd model with real and imaginary coupling. It is known that $`c<1`$ minimal CFT perturbed by the operators $`\mathrm{\Phi }_{12}`$ and $`\mathrm{\Phi }_{13}`$ can be obtained by a quantum group (QG) restriction of the sine-Gordon and imaginary Bullough-Dodd model with special values of the coupling. The VEVs of primary fields were then calculated. The same method was applied later to integrable QFTs involving more than one field. For instance, the VEV for a two-parameter family of integrable QFTs gave rise to the VEV of local operators in parafermionic sine-Gordon models and in integrable perturbed $`SU(2)`$ coset CFT . The VEV of simply-laced affine Toda field theories (ATFT)s is known for a long time and the case of non-simply laced dual pairs was recently studied in for which a general expression for the VEVs was derived.
Such perturbed CFTs have recently attracted much attention in condensed matter physics, such as in the context of point contacts in the fractional quantum Hall effect and impurities in quantum wires . In such cases the property of integrability has provide a non-perturbative answer for experimentally important strongly interacting solid state physics problems . Particulary, on-shell results were obtained using exact relativistic scattering and related form factor techniques .
In this paper, we are interested in exact off-shell results for two coupled conformal field theories <sup>1</sup><sup>1</sup>1In literature, the first example of such integrable coupled models was studied in . for which the inter-layer coupling preserves integrability. The on-shell dynamics of these models were studied in .
Let us briefly recall the ideas of . We consider two planar systems corresponding to two coupled minimal models $`_{p/p^{}}`$ which interact through a relevant operator which leads to an integrable theory. The resulting action can be written :
$`𝒜=_{p/p^{}}+_{p/p^{}}+\lambda {\displaystyle d^2x\mathrm{\Phi }_{12}^{(1)}\mathrm{\Phi }_{12}^{(2)}}`$ (1)
or
$`\stackrel{~}{𝒜}=_{p/p^{}}+_{p/p^{}}+\widehat{\lambda }{\displaystyle d^2x\mathrm{\Phi }_{21}^{(1)}\mathrm{\Phi }_{21}^{(2)}},`$ (2)
where we denote respectively $`\mathrm{\Phi }_{12}^{(1)}(\mathrm{\Phi }_{21}^{(1)})`$ and $`\mathrm{\Phi }_{12}^{(2)}(\mathrm{\Phi }_{21}^{(2)})`$ as two specific primary operators of each unperturbed minimal models and where the parameters $`\lambda `$ and $`\widehat{\lambda }`$ characterize the strength of the interaction. Here, we will be interested in exact one-point functions in such system.
In section 2 we introduce the notations and those known results which are useful for our purpose. In section 3 using the exact result for the VEV of the $`C_2^{(1)}`$ ATFT derived in , we deduce the exact VEV $`<\mathrm{\Phi }_{rs}^{(1)}(x)\mathrm{\Phi }_{r^{}s^{}}^{(2)}(x)>`$ for any values of $`(r,s),(r^{},s^{})`$ in the model with action (1). To do it, we relate the parameter $`\lambda `$ in (1) to the masses of the particles and we perform the QG restriction of the $`C_2^{(1)}`$ ATFT with imaginary coupling which leads to the model (1). For $`(r,s)=(r^{},s^{})=(1,2)`$ this VEV can be calculated exactly as well as the bulk free energy. The specific case $`_{3/4}+_{3/4}`$ coupled by $`\mathrm{\Phi }_{12}^{(1)}\mathrm{\Phi }_{12}^{(2)}`$ is considered in details. It corresponds to two layer Ising models coupled by their magnetization operator $`\sigma ^{(1)}\sigma ^{(2)}`$. The previous approach is also extended to the model described by action (2). In section 4, we extract some limited information about the asymptotics of two-point correlation functions between any pairs of primary operators which belong to the same or different unperturbed models. More precisely, we will distinguish four different cases :
$`(a)`$ $`:`$ $`<\mathrm{\Phi }_{rs}^{(1)}(x)\mathrm{\Phi }_{r^{}s^{}}^{(2)}(y)>\text{for}|xy|0,`$
$`(b)`$ $`:`$ $`<\mathrm{\Phi }_{rs}^{(1)}(x)\mathrm{\Phi }_{r^{}s^{}}^{(2)}(y)>\text{for}|xy|\mathrm{},`$
$`(c)`$ $`:`$ $`<\mathrm{\Phi }_{rs}^{(i)}(x)\mathrm{\Phi }_{r^{}s^{}}^{(i)}(y)>\text{for}|xy|0\text{and}i\{1,2\},`$
$`(d)`$ $`:`$ $`<\mathrm{\Phi }_{rs}^{(i)}(x)\mathrm{\Phi }_{r^{}s^{}}^{(i)}(y)>\text{for}|xy|\mathrm{}\text{and}i\{1,2\}`$
which are depicted in figure 1. We finally give some numerical results for various examples of coupled minimal models such as two energy-energy coupled tricritical Ising, two coupled $`A_5`$-RSOS models and two energy-energy coupled 3-state Potts model. Perspective and conclusions follow in this final section.
Figure 1 \- Two coupled two dimensional models. Short distance
results are obtained by taking the limit $`ϵ0`$.
## 2 Coupled minimal models as restricted $`C_2^{(1)}`$ ATFT
The ATFT with real coupling $`b`$ associated with the affine Lie algebra $`C_2^{(1)}`$ is described by the action in the Euclidean space :
$`𝒜={\displaystyle d^2x\left[\frac{1}{8\pi }(_\mu 𝝋)^2+\mu ^{}e^{2b\phi _1}+\mu ^{}e^{2b\phi _2}+\mu e^{b(\phi _1\phi _2)}\right]}`$ (3)
where we choosed the convention that the length squared of the long roots is four. As the different vertex operators do not renormalize in the same way, we introduced two scale parameters $`\mu `$ and $`\mu ^{}`$. The fields in eq. (3) are normalized such that:
$`<\phi _i(x)\phi _j(y)>=\delta _{ij}\mathrm{ln}|xy|^2.`$ (4)
This model possess two fundamental particles with mass $`M_a`$ which depends on one parameter $`\overline{m}`$ :
$`M_a=2\overline{m}\mathrm{sin}({\displaystyle \frac{\pi a}{H}})\text{for}a=1,2`$ (5)
where we introduced the “deformed” Coxeter <sup>2</sup><sup>2</sup>2Differently to the simply laced case for which the mass ratios take the classical values, the mass ratios for non-simply laced case get quantum corrections. number $`H=h(1B)+h^{}B`$ with $`B=\frac{b^2}{1+b^2}`$ and $`h=4`$, $`h^{}=6`$ are respectively the Coxeter and dual Coxeter numbers. The exact relation between $`\overline{m}`$, $`\mu `$ and $`\mu ^{}`$ was found in and is given by :
$`\left(\pi \mu \gamma (1+b^2)\right)\left(\pi \mu ^{}\gamma (1+2b^2)\right)=2^{\frac{4}{1B}}\left(\overline{m}{\displaystyle \frac{\mathrm{\Gamma }((1B)/H)\mathrm{\Gamma }(1+B/H)}{\mathrm{\Gamma }(1/H)}}\right)^{\frac{H}{1B}}.`$ (6)
In ATFT approach to perturbed CFT, one usually identifies the perturbation with the affine extension of the Lie algebra $`𝒢`$. Instead, the perturbation will be associated here with the standard (length 2) root of $`C_2^{(1)}`$. Removing the last term in the action (3) leaves a $`D_2=SO(4)=SU(2)SU(2)`$ model, i.e. two decoupled Liouville models. To associate the two first terms of the $`C_2^{(1)}`$ Toda potential to two decoupled conformal field theories, we first introduce for each one a background charge at infinity. Then, the total stress-energy tensor $`T(z)`$ is written $`T(z)=T^{(1)}(z)+T^{(2)}(z)`$ where :
$`T^{(i)}(z)={\displaystyle \frac{1}{2}}(\phi _i)^2+Q_i^2\phi _i\text{for}i\{1,2\}`$ (7)
ensures the local conformal invariance of the $`D_2`$ model for the choice $`Q_2=Q_1=b+1/2b`$. With our conventions <sup>3</sup><sup>3</sup>3Here, the length of the longest root is chosen to be 4., the exponential fields
$`e^{a_i\phi _i(x)}\text{for}i\{1,2\}`$ (8)
are spinless conformal primary fields of each Liouville model with conformal dimensions
$`\mathrm{\Delta }(e^{a_i\phi _i(x)})={\displaystyle \frac{a_i^2}{2}}+a_iQ_i\text{and}i\{1,2\}.`$ (9)
In particular, the exponential fields $`e^{2b\phi _1}`$ and $`e^{2b\phi _2}`$ have conformal dimensions 1. As is well known, the “minimal model” $`_{p/p^{}}`$ with central charge $`c=16\frac{(pp^{})^2}{pp^{}}`$ can be obtained from the Liouville case. Consequently, the $`D_2`$ CFT can be identified with two decoupled minimal models by the substitution :
$`bi\beta ,\mu \mu ,\mu ^{}\mu ^{},`$ (10)
and the choice :
$`\beta ^2=\beta _+^2=p/2p^{}\text{or}\beta ^2=\beta _{}^2=p^{}/2p\text{with}p<p^{}.`$ (11)
Similarly, the primary operators of each minimal model $`_{p/p^{}}`$ are obtained through the substitution $`a_ji\eta _j`$ for $`j=1,2`$ in (8). With these substitutions, the conformal dimension of the perturbing operator becomes :
$`\mathrm{\Delta }_{pert}=\mathrm{\Delta }(e^{i\beta (\phi _1\phi _2)})=3\beta ^21.`$ (12)
As long as we consider a relevant perturbation, we are restricted to choose $`\beta ^2<2/3`$. In the following we will consider only the cases for which this condition is satisfied, in particular $`\beta ^2=\beta _+^2`$.
We define respectively $`\{\mathrm{\Phi }_{rs}^{(1)}\}`$ and $`\{\mathrm{\Phi }_{r^{}s^{}}^{(2)}\}`$ as the two sets of primary fields with conformal dimensions :
$`\mathrm{\Delta }_{rs}={\displaystyle \frac{(p^{}rps)^2(pp^{})^2}{4pp^{}}}\text{for}1r<p,1s<p^{}\text{and}p<p^{}.`$ (13)
They are simply related to the vertex operators of each minimal model through the relation:
$`\mathrm{\Phi }_{rs}^{(i)}(x)=N_{rs}^{(i)1}\mathrm{exp}(i\eta _i^{rs}\phi _i(x))\text{with}\eta _1^{rs}=\eta _2^{rs}={\displaystyle \frac{(1r)}{2\beta }}(1s)\beta ,`$ (14)
where we have introduced the normalization factors $`N_{rs}^{(i)}`$ for each model. These numerical factors depend on the normalization of the primary fields. Here, they are chosen in such a way that the primary fields satisfy the conformal normalization condition :
$`<\mathrm{\Phi }_{rs}^{(i)}(x)\mathrm{\Phi }_{rs}^{(i)}(y)>_{CFT}={\displaystyle \frac{1}{|xy|^{4\mathrm{\Delta }_{rs}}}}\text{for}i\{1,2\}.`$ (15)
For further convenience, we write these coefficients $`N_{rs}^{(i)}=N^{(i)}(\eta _i^{rs})`$ where :
$`N^{(1)}(\eta )=\left[\pi \mu ^{}\gamma (2\beta ^2)\right]^{\frac{\eta }{2\beta }}\left[{\displaystyle \frac{\mathrm{\Gamma }(2\beta ^2+2\eta \beta )\mathrm{\Gamma }(1/2\beta ^2\eta /\beta )\mathrm{\Gamma }(22\beta ^2)\mathrm{\Gamma }(21/2\beta ^2)}{\mathrm{\Gamma }(22\beta ^22\eta \beta )\mathrm{\Gamma }(21/2\beta ^2+\eta /\beta )\mathrm{\Gamma }(2\beta ^2)\mathrm{\Gamma }(1/2\beta ^2)}}\right]^{\frac{1}{2}}`$
and $`N^{(2)}(\eta )=N^{(1)}(\eta )`$.
For imaginary values of the coupling $`b=i\beta `$, the $`C_2^{(1)}`$ ATFT possesses complex soliton solutions which interpolate between the degenerate vacua. This QFT possesses the QG symmetry associated to $`U_q(D_3^{(2)})`$ \- as we will recall in the next section. In the $`S`$-matrix of the $`B_2^{(1)}`$ ATFT was constructed using the $`U_q(A_3^{(2)})`$ QG symmetry of this QFT. In particular, using the bootstrap procedure, the authors deduced the breather-breather $`S`$-matrix. It was also shown that a breather-particle identification holds by comparing the $`S`$-matrix elements of the lowest-breathers (breathers with lowest mass) with the $`S`$-matrix elements for the quantum particles in real ATFT. Following the conventions of , the particle spectrum in real $`B_2^{(1)}`$ ATFT is given by $`M_1=2\overline{m}\mathrm{sin}(\pi /H)`$ and $`M_2=\overline{m}`$. Using the results of we find the identification :
$`\overline{m}=2M\mathrm{sin}({\displaystyle \frac{\pi \xi }{42\xi }})\text{with}\xi ={\displaystyle \frac{\beta ^2}{1\beta ^2}}`$ (16)
where we denote $`M`$ as the mass of the lowest kink. Eq. (16) holds for our case due to the identification $`B_2^{(1)}C_2^{(1)}`$ and $`A_3^{(2)}D_3^{(2)}`$.
## 3 Expectation values in coupled minimal models
In ref. we derived the exact VEV $`G(𝒂)=<\mathrm{exp}(𝒂.𝝋)(x)>`$ for all non-simply laced ATFTs using the so-called “reflection relations” which relate different fields with the same quantum numbers. We refer the reader to these papers for more details. Although the model (1) is very different from the $`C_2^{(1)}`$ ATFT (3) in its physical content (the model (1) contains solitons and excited solitons), there are good reasons to believe that the expectation values obtained in the real coupling case provide also the expectation values for imaginary coupling. The calculation of the VEVs in both cases ($`b`$ real or imaginary) within the standard perturbation theory agree through the identification $`b=i\beta `$. Following the analysis done for the Bullough-Dodd model , it is then straightforward to obtain the VEV of primary operators which belong to different minimal models. With the substitutions (10) one gets $`𝒢(𝜼)=G(i𝜼)`$ :
$`𝒢(𝜼)`$ $`=`$ $`\left[\pi ^2\mu \mu ^{}\gamma ({\displaystyle \frac{1}{1+\xi }})\gamma ({\displaystyle \frac{1\xi }{1+\xi }})\right]^{\frac{(1+\xi )}{2\xi }(\mathrm{\Delta }_1+\mathrm{\Delta }_2)}\left[\pi \mu ^{}\gamma ({\displaystyle \frac{1\xi }{1+\xi }})\right]^{\frac{(1+\xi )}{2\xi }\beta (\eta _1\eta _2)}\left[{\displaystyle \frac{1}{1+\xi }}\right]^{𝜼^2}`$ (17)
$`\times \mathrm{exp}\left[{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t}}\left({\displaystyle \frac{\chi (𝜼,t)}{\mathrm{sinh}((1+\xi )t)\mathrm{sinh}(2t\xi )\mathrm{sinh}((42\xi )t)}}𝜼^2e^{2t}\right)\right]`$
with
$`\chi (𝜼,t)`$ $`=`$ $`2[\mathrm{sinh}((\eta _1\eta _2)\beta (1+\xi )t)\mathrm{sinh}(((\eta _1\eta _2)\beta (1+\xi )+2\xi 2)t)`$
$`+\mathrm{sinh}^2((\eta _1+\eta _2)\beta (1+\xi )t)]\mathrm{sinh}(t)\mathrm{cosh}(t\xi )`$
$`+[\mathrm{sinh}(2\eta _2\beta (1+\xi )t)\mathrm{sinh}((2\eta _2\beta (1+\xi )2)t)`$
$`+\mathrm{sinh}(2\eta _1\beta (1+\xi )t)\mathrm{sinh}((2\eta _1\beta (1+\xi )+2)t)]\mathrm{sinh}((1\xi )t).`$
As usual we denote $`\gamma (x)=\mathrm{\Gamma }(x)/\mathrm{\Gamma }(1x)`$. The integral in (17) is convergent if :
$`{\displaystyle \frac{2}{(1+\xi )}}<(\eta _1+\eta _2)\beta <{\displaystyle \frac{2}{(1+\xi )}}\text{and}1<(\eta _1\eta _2)\beta <{\displaystyle \frac{3\xi }{(1+\xi )}},`$ (18)
and is defined via analytic continuation outside this domain. In (17) we defined $`\mathrm{\Delta }_1`$ and $`\mathrm{\Delta }_2`$, the conformal dimensions in the imaginary coupling case by $`\mathrm{\Delta }_1=\mathrm{\Delta }(e^{i\eta _1\phi _1})=\frac{\eta _1^2}{2}+\frac{\eta _1\beta }{2}(\frac{\xi 1}{\xi })`$ and similarly for $`\mathrm{\Delta }_2`$ with the change $`\eta _1\eta _2`$.
If we want to express the final result for the VEV in terms of the parameter $`\lambda `$ in the action (1), we need the exact relation between $`\lambda `$ and the parameters $`\mu ,\mu ^{}`$ in the $`C_2^{(1)}`$ ATFT with imaginary coupling. We obtain :
$`\lambda ={\displaystyle \frac{\pi \mu \mu ^{}}{(4\beta ^21)^2}}\gamma (4\beta ^2)\gamma ^2(12\beta ^2),`$ (19)
which corresponds to $`\lambda =\mu N^{(1)}(\beta )N^{(2)}(\beta )`$.
Like any other ATFT, the model (3) for imaginary values of the coupling has non-local conserved charges $`\{Q_k,\overline{Q}_k\}`$ for $`k=0,1,2`$ generated respectively by the purely chiral and anti-chiral components :
$`J_{𝐞_k^{}}(z)=e^{\frac{i}{\beta }𝐞_k^{}.\phi (z)}\text{and}\overline{J}_{𝐞_k^{}}(\overline{z})=e^{\frac{i}{\beta }𝐞_k^{}.\overline{\phi }(\overline{z})},`$ (20)
where the fundamental vector field $`𝝋(z,\overline{z})=\phi (z)+\overline{\phi }(\overline{z})`$ and $`\{𝐞_k^{}\}`$, $`k=0,1,2`$ is the set of dual simple roots of the non-simply laced affine Lie algebra $`C_2^{(1)}`$. We also define the topological charge :
$`H_k={\displaystyle \frac{\beta }{2i\pi }}{\displaystyle d^2x𝐞_k^{}}._x𝝋(x,t).`$ (21)
Using the equal-time braiding relations for all $`x,y`$ :
$`J_{𝐞_k^{}}(x,t)\overline{J}_{𝐞_l^{}}(y,t)=q_l^{a_{kl}}\overline{J}_{𝐞_l^{}}(y,t)J_{𝐞_k^{}}(x,t)\text{with}q_l=e^{\frac{2i\pi }{|𝐞_k|^2\beta ^2}},`$ (22)
where $`a_{kl}=\frac{2𝐞_k.𝐞_l}{|𝐞_k|^2}`$ denotes the extended Cartan matrix of $`C_2^{(1)}`$, one can show that the charges $`\{Q_k,\overline{Q}_k,H_k\}`$ for $`k=0,1,2`$ satisfy the quantum universal enveloping algebra $`U_q(D_3^{(2)})`$ with deformation parameter $`qq_0=e^{\frac{i\pi }{2\beta ^2}}`$. If we express these generators in terms of the standard Chevalley basis $`\{E_k^+,E_k^{},H_k\}`$ by :
$`Q_kE_k^+q^{H_k}\text{and}\overline{Q}_kE_k^{}q^{H_k}\text{for}k=0,1,2,`$ (23)
then we have :
$`[H_k,H_l]=0,`$ (24)
$`[H_k,E_l^\pm ]=\pm a_{l,k}E_l^\pm ,`$
$`[E_k^+,E_l^{}]=\delta _{kl}{\displaystyle \frac{q^{2H_l}q^{2H_l}}{q_lq_l^1}}.`$
The $`U_q(D_3^{(2)})`$ has two subalgebras $`U_q(D_2)`$ as well as $`U_{q^2}(C_2)`$ where the subalgebra $`U_q(D_2)`$ is generated by $`\{Q_0,\overline{Q}_0,Q_2,\overline{Q}_2,H_0,H_2\}`$. As was found in the $`S`$-matrix in the unrestricted form acts as an intertwinner on the modules of the $`U_q(D_2)`$ representations :
$`S_{a,b}(\theta ):𝒱_{\rho _a}𝒱_{\rho _b}𝒱_{\rho _b}𝒱_{\rho _a}`$ (25)
where $`𝒱_\rho `$ is the module with highest weight $`\rho `$. The scattering of two solitons of species $`a`$ and $`b`$ is then described by $`S_{a,b}(\theta _{12})`$ with $`\theta _{12}`$ being their relative rapidity. There are two such fundamental multiplets denoted $`\{4\}`$ and $`\{6\}`$ in ref. . In addition, there are also scalar bound states and excited solitons depending on the values of $`(p,p^{})`$ chosen.
To understand the restricted $`C_2^{(1)}`$ (denoted $`RC_2^{(1)}`$ below), we use the general framework of superselection sectors (for details see ). The model (1) is a perturbation of the two minimal models by the operator $`\mathrm{\Phi }_{12}^{(1)}\mathrm{\Phi }_{12}^{(2)}`$. Each minimal model contains a finite number of primary fields (13,14). Using the superselection criterion for the present case <sup>4</sup><sup>4</sup>4We only consider the holomorphic part of the primary operators but keep the same notation.:
$`\left(\mathrm{\Phi }_{2j+\mathrm{1\; 1}}^{(1)}(z)\mathrm{\Phi }_{2\stackrel{~}{j}+\mathrm{1\; 1}}^{(2)}(z)\right)\left(\mathrm{\Phi }_{12}^{(1)}(w)\mathrm{\Phi }_{12}^{(2)}(w)\right){\displaystyle \frac{1}{(zw)^l}}\mathrm{\Phi }_{2j+\mathrm{1\; 2}}^{(1)}(w)\mathrm{\Phi }_{2\stackrel{~}{j}+\mathrm{1\; 2}}^{(2)}(w),l`$ (26)
we find $`j+\stackrel{~}{j}`$ where $`j\stackrel{~}{j}`$ denote the representations of $`U_q(D_2)`$ with $`j`$ the spin-$`j`$ representation of $`SU(2)`$ with dimension $`2j+1`$ (and similarly for $`\stackrel{~}{j}`$). If $`q`$ is a root of unity i.e if eq. (11) is satisfied then using the Coulomb gas representation condition (13) $`12j+1p1`$ and $`12\stackrel{~}{j}+1p1`$ one has :
$`0jp/21\text{and}0\stackrel{~}{j}p/21.`$ (27)
The superselection sectors $`_{j\stackrel{~}{j}}^{RC_2^{(1)}}`$ of the $`RC_2^{(1)}`$ model (1) are thus :
$`^{RC_2^{(1)}}={\displaystyle \underset{(j,\stackrel{~}{j})\{0,1/2,\mathrm{},p/21\}}{}}_{j\stackrel{~}{j}}^{RC_2^{(1)}}\text{with}j+\stackrel{~}{j}.`$ (28)
As shown in for the unitary series ($`p^{}=p+1`$), after the quantum group retriction (11) for $`p>3`$ the fundamental solitons in the $`\{4\}`$ and $`\{6\}`$ representation of $`U_q(D_2)`$ become the RSOS kink $`K_{j_2\stackrel{~}{j}_2j_1\stackrel{~}{j}_1}^{j\stackrel{~}{j}}`$. These kinks interpolate between different vacua $`|0_{j_1\stackrel{~}{j}_1}>`$ and $`|0_{j_2\stackrel{~}{j}_2}>`$ which are connected using the $`U_q(sl_2)`$ fusion ring at $`q=e^{i\frac{\pi }{p}}`$:
$`j_1\times j={\displaystyle \underset{j_2=|j_1j|}{\overset{min(j_1+j,p2j_1j)}{}}}j_2`$ (29)
and similarly for $`\stackrel{~}{j}`$. However, for $`p=3`$, $`(j,\stackrel{~}{j})\{0,1/2\}`$ and then the $`\{6\}`$ is projected out of the spectrum, leaving only the $`\{4\}`$. We refer the reader to for more details.
From the previous remarks and the identification $`D_2=SO(4)=SU(2)SU(2)`$, by analogy with the primary fields $`\mathrm{\Phi }_{1s}^{(1)}`$ and $`\mathrm{\Phi }_{1s^{}}^{(2)}`$ commute with the generators in (3) (for $`k=0,2`$) of the subalgebra $`U_q(D_2)U_q(D_3^{(2)})`$. If one interpret the fields $`\mathrm{\Phi }_{2j+\mathrm{1\; 1}}^{(1)}(z)\mathrm{\Phi }_{2\stackrel{~}{j}+\mathrm{1\; 1}}^{(2)}(z)`$ as the highest component fields in the multiplet, it can be shown that the primary operators $`\mathrm{\Phi }_{rs}^{(1)}`$ and $`\mathrm{\Phi }_{r^{}s^{}}^{(2)}`$ are not invariant with respect to $`U_q(D_2)`$. Together with some non-local fields they form finite-dimensional representation of this algebra. Consequently, the VEV in (1) should take into account the factor $`d_{rs,r^{}s^{}}^{j\stackrel{~}{j}}`$ coming from the QG restriction of the QFT (3). Following the conjecture of it takes the form :
$`d_{rs,r^{}s^{}}^{j\stackrel{~}{j}}={\displaystyle \frac{\mathrm{sin}(\frac{\pi (2j+1)}{p}|p^{}rps|)}{\mathrm{sin}(\frac{\pi (2j+1)}{p}(p^{}p))}}{\displaystyle \frac{\mathrm{sin}(\frac{\pi (2\stackrel{~}{j}+1)}{p}|p^{}r^{}ps^{}|)}{\mathrm{sin}(\frac{\pi (2\stackrel{~}{j}+1)}{p}(p^{}p))}}.`$ (30)
Using the notations introduced in the previous section, and eqs. (14), (17), the outcome for the VEV between different primary operators is :
$`<0_{j\stackrel{~}{j}}|\mathrm{\Phi }_{rs}^{(1)}(x)\mathrm{\Phi }_{r^{}s^{}}^{(2)}(x)|0_{j\stackrel{~}{j}}>`$ $`=`$ $`d_{rs,r^{}s^{}}^{j\stackrel{~}{j}}\left[{\displaystyle \frac{\pi \lambda \gamma (\frac{1}{1+\xi })(1+\xi )^{\frac{42\xi }{1+\xi }}}{\gamma (\frac{3\xi 1}{1+\xi })\gamma (\frac{1\xi }{1+\xi })}}\right]^{\frac{(1+\xi )}{2\xi }(\mathrm{\Delta }_{rs}+\mathrm{\Delta }_{r^{}s^{}})}`$ (31)
$`\times \mathrm{exp}𝒬_{12}((1+\xi )r2\xi s,(1+\xi )r^{}2\xi s^{}).`$
The function $`𝒬_{12}(\theta ,\theta ^{})`$ for $`|\theta \pm \theta ^{}|<4\xi `$ and $`\xi >\frac{1}{3}`$ is given by the integral :
$`𝒬_{12}(\theta ,\theta ^{})={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t}}\left({\displaystyle \frac{\mathrm{\Psi }_{12}(\theta ,\theta ^{},t)}{\mathrm{sinh}((1+\xi )t)\mathrm{sinh}(2t\xi )\mathrm{sinh}((42\xi )t)}}{\displaystyle \frac{\theta ^2+\theta _{}^{}{}_{}{}^{2}2(1\xi )^2}{4\xi (\xi +1)}}e^{2t}\right)`$
with
$`\mathrm{\Psi }_{12}(\theta ,\theta ^{},t)`$ $`=`$ $`[\mathrm{cosh}((\theta +\theta ^{})t)\mathrm{cosh}((\theta \theta ^{})t)`$
$`\mathrm{cosh}((22\xi )t)]\mathrm{sinh}((1\xi )t)\mathrm{cosh}((42\xi )t)`$
$`+\left[\mathrm{cosh}((\theta +\theta ^{})t)+\mathrm{cosh}((\theta \theta ^{})t)\mathrm{cosh}((22\xi )t)1\right]\mathrm{sinh}(t)\mathrm{cosh}(t\xi ).`$
and defined by analytic continuation outside this domain. Notice that eq. (31) satisfies:
$`<0_{j\stackrel{~}{j}}|\mathrm{\Phi }_{rs}^{(1)}(x)\mathrm{\Phi }_{r^{}s^{}}^{(2)}(x)|0_{j\stackrel{~}{j}}>=<0_{j\stackrel{~}{j}}|\mathrm{\Phi }_{prp^{}s}^{(1)}(x)\mathrm{\Phi }_{pr^{}p^{}s^{}}^{(2)}(x)|0_{j\stackrel{~}{j}}>.`$ (32)
A particular case of eq. (31) is the expectation value of the perturbing operator which can be calculated explicitly to give the result :
$`<0_{j\stackrel{~}{j}}|\mathrm{\Phi }_{12}^{(1)}(x)\mathrm{\Phi }_{12}^{(2)}(x)|0_{j\stackrel{~}{j}}>`$ $`=`$ $`{\displaystyle \frac{1}{\lambda }}\left[{\displaystyle \frac{\pi \lambda \gamma (\frac{1}{1+\xi })}{\gamma (\frac{3\xi 1}{1+\xi })\gamma (\frac{1\xi }{1+\xi })}}\right]^{\frac{1+\xi }{2\xi }}{\displaystyle \frac{2^{\frac{3\xi 2}{2\xi }}(1+\xi )}{\pi (\xi 2)}}{\displaystyle \frac{\gamma (\frac{\xi }{42\xi })\gamma (\frac{1}{42\xi })}{\gamma (\frac{1+\xi }{42\xi })}}.`$ (33)
By using the exact relation between the mass parameter $`\overline{m}`$ and the mass of the kink $`M`$ (16) and eqs. (6), (19) we immediately obtain the relation between $`M`$ and $`\lambda `$ :
$`M={\displaystyle \frac{2^{\frac{\xi }{2\xi }}\mathrm{\Gamma }(\frac{\xi }{42\xi })\mathrm{\Gamma }(\frac{1}{42\xi })}{\pi \mathrm{\Gamma }(\frac{1+\xi }{42\xi })}}\left[{\displaystyle \frac{\pi \lambda \gamma (\frac{1}{1+\xi })}{\gamma (\frac{3\xi 1}{1+\xi })\gamma (\frac{1\xi }{1+\xi })}}\right]^{\frac{1+\xi }{42\xi }}.`$ (34)
Consequently, according to eqs. (19), (33), (34) and $`\beta ^2<2/3`$, the perturbed CFTs (1) develop a massive spectrum for ($`m(\lambda )=0`$) :
$`(i)0<\xi <1/3,\lambda >0\text{i.e.}0<{\displaystyle \frac{p}{p^{}}}<{\displaystyle \frac{1}{2}},`$ (35)
$`(ii)1/3<\xi <1,\lambda <0\text{i.e.}{\displaystyle \frac{1}{2}}<{\displaystyle \frac{p}{p^{}}}<1.`$
where $`\xi =\frac{p}{2p^{}p}`$. In particular, the condition $`(ii)`$ is always satisfied for the coupled unitary minimal models defined by (1).
Finally, the expectation value (33) can be used to derive the bulk free energy :
$`f_{12}=\text{ }\text{limV→∞}{\displaystyle \frac{1}{V}}\mathrm{ln}Z,`$ (36)
where $`V`$ is the volume of the 2D space and $`Z`$ is the singular part of the partition function associated with action (1). Using
$`_\lambda f_{12}=<0_{j\stackrel{~}{j}}|\mathrm{\Phi }_{12}^{(1)}\mathrm{\Phi }_{12}^{(2)}|0_{j\stackrel{~}{j}}>,`$ (37)
and eqs. (33), (34), the result for the bulk free energy follows :
$`f_{12}={\displaystyle \frac{M^2\mathrm{sin}(\frac{\pi }{42\xi })}{2}}{\displaystyle \frac{\mathrm{sin}(\frac{\pi \xi }{42\xi })}{\mathrm{sin}(\frac{\pi (1+\xi )}{42\xi })}}.`$ (38)
As was suggested in , a particular case of the model (1) can be related to the SG theory at the reflectionless point as well as the $`D_8^{(1)}`$ ATFT with imaginary coupling. Let us first consider the SG theory with action :
$`𝒜_{SG}={\displaystyle d^2x\left[\frac{1}{16\pi }(_\nu \varphi )^2+\mathrm{\Lambda }\mathrm{cos}(\widehat{\beta }\varphi )\right]}.`$ (39)
This model is integrable and its on-shell solutions, i.e. spectrum of particles and $`S`$-matrix are well-kown . This theory contains soliton, antisoliton and soliton-antisoliton bound-state (breathers) $`B_n`$, $`n=1,\mathrm{},1/\widehat{\xi }`$ where $`\widehat{\xi }=\frac{\widehat{\beta }^2}{1\widehat{\beta }^2}`$. The lightest bound-state $`B_1`$ coincides with the field $`\varphi `$ in the perturbative treatment of the QFT (39). Its mass is given by $`m_1=2M_{SG}\mathrm{sin}(\frac{\pi \widehat{\xi }}{2})`$ where $`M_{SG}`$ denotes the soliton mass of the SG model. For $`p=3`$, $`p^{}=4`$ the model (1) corresponds to two magnetically coupled Ising models. It is a $`c=1`$ CFT perturbed by the operator $`\mathrm{\Phi }_{12}^{(1)}\mathrm{\Phi }_{12}^{(2)}`$ with conformal dimension $`1/8`$. Indeed, the VEV of this operator must be simply related to the VEV of $`<\mathrm{cos}(\widehat{\beta }\varphi )>`$ in the SG model for $`\widehat{\beta }^2=1/8`$, which is also a $`c=1`$ CFT with perturbing operator of conformal dimension $`1/8`$. For this latter value, $`\widehat{\xi }=1/7`$ and the SG model possesses 6 neutral excitations :
$`m_a=2M_{SG}\mathrm{sin}({\displaystyle \frac{\pi a}{14}}),a=1,\mathrm{},6;`$ (40)
$`m_1<m_2<M_{SG}<m_3<\mathrm{}<m_6.`$
Similarly, for $`\beta ^2=3/8`$, $`\xi =3/5`$ the spectrum of the model (1) was described in and summarize here. Here, we denote the lowest mass by $`M`$. If we compare the spectrum of each model, using the notations of we have :
$`m_1M,m_2m_{B_1^{(1)}},M_{SG}m_{B_1^{(2)}},\text{etc}\mathrm{}`$ (41)
Furthermore, the bulk free energy of the SG model (39) is well known and given by :
$`f_{SG}={\displaystyle \frac{M_{SG}^2}{4}}\mathrm{tan}({\displaystyle \frac{\pi \widehat{\xi }}{2}}).`$ (42)
By expressing it in terms of the lightest particle $`m_1`$ using (40) and the relation
$`8\mathrm{sin}(\pi /14)\mathrm{sin}(3\pi /14)\mathrm{cos}(\pi /7)=1`$ we find the same result as in eq. (38) for $`\beta ^2=3/8`$, i.e. $`\xi =3/5`$. Also, using the results of for $`\widehat{\beta }^2=1/8`$, one has :
$`\mathrm{\Lambda }={\displaystyle \frac{2\gamma (1/8)}{\pi }}\left[{\displaystyle \frac{M_{SG}\sqrt{\pi }}{2}}{\displaystyle \frac{\mathrm{\Gamma }(4/7)}{\mathrm{\Gamma }(1/14)}}\right]^{\frac{7}{4}}.`$ (43)
Similarly, the VEV of the perturbing operator is given by :
$`<\mathrm{cos}(\widehat{\beta }\varphi )>_{\widehat{\beta }^2=1/8}={\displaystyle \frac{8}{7}}{\displaystyle \frac{\mathrm{\Gamma }(1/7)\mathrm{\Gamma }(6/7)}{\mathrm{\Gamma }^2(4/7)\mathrm{\Gamma }^2(13/14)}}\left[{\displaystyle \frac{\pi }{\gamma (1/8)}}\right]^{8/7}({\displaystyle \frac{\mathrm{\Lambda }}{2}})^{1/7}.`$ (44)
Comparing eqs. (43), (44) with eqs. (33) and (34) respectively for $`\xi =3/5`$ and using eq. (41), one find the relations :
$`\mathrm{\Lambda }2^{\frac{1}{2}}\lambda \text{and}\mathrm{cos}(\widehat{\beta }\varphi )|_{\widehat{\beta }^2=1/8}2^{\frac{1}{2}}\mathrm{\Phi }_{12}^{(1)}\mathrm{\Phi }_{12}^{(2)}`$ (45)
in perfect agreement <sup>5</sup><sup>5</sup>5At $`\widehat{\beta }^2=1/8`$, the operator can be expressed in terms of two independent Ising spin fields, $`\sigma ^{(i)}=\mathrm{\Phi }_{12}^{(i)}`$ for $`i=1,2`$. with eq. (B.21) in .
Let us now turn to the relation with the $`D_8^{(1)}`$ ATFT. The calculation of the bulk free energy for the simply-laced $`D_8^{(1)}`$ ATFT gives :
$`f_{D_8^{(1)}|_b}={\displaystyle \frac{\overline{m}^2}{8}}{\displaystyle \frac{\mathrm{sin}(\frac{\pi }{14})}{\mathrm{sin}(\frac{\pi B}{14})\mathrm{sin}(\frac{\pi (1B)}{14})}}`$ (46)
where $`B`$ is defined in section 2 and the mass of the particles are related with the mass parameter $`\overline{m}`$ by : $`m_8^2=m_7^2=2\overline{m}^2`$, $`m_a^2=8\overline{m}^2\mathrm{sin}^2(\frac{\pi a}{14})`$ for $`a=1,\mathrm{},6`$. For imaginary coupling, one expects the particle-breather identification :
$`8\overline{m}^2\mathrm{sin}^2({\displaystyle \frac{\pi }{14}})=\left(2M\mathrm{sin}({\displaystyle \frac{\pi \xi }{14}})\right)^2`$ (47)
where $`\xi `$ is defined as in eq. (16). This yields :
$`f_{D_8^{(1)}|_{b=i\beta }}={\displaystyle \frac{M^2}{16}}{\displaystyle \frac{\mathrm{sin}(\frac{\pi \xi }{14})}{\mathrm{sin}(\frac{\pi }{14})\mathrm{sin}(\frac{\pi (1+\xi )}{14})}},`$ (48)
which for $`\xi =7`$ is in perfect agreement with eq. (38) evaluated at $`\xi =3/5`$.
It is also interesting to study the behaviour of the model (1) for $`p^{}=p+1`$ in the limit $`p\mathrm{}`$. If we consider the sine-Gordon model with action (39) there is an equivalent description in terms of the massive Thirring model :
$`𝒜_{MT}={\displaystyle d^2x\left[i\overline{\psi }\gamma ^\nu _\nu \psi M_{SG}\overline{\psi }\psi \frac{g}{2}(\overline{\psi }\gamma ^\nu \psi )^2\right]}`$ (49)
where $`\psi ,\overline{\psi }`$ is a Dirac field. The four fermion coupling constant $`g`$ relates to $`\widehat{\beta }`$ in (39) by $`\frac{g}{\pi }=\frac{1}{2\widehat{\beta }^2}1`$. Also the free fermion point is reached for $`\widehat{\beta }^2\frac{1}{2}`$. Using the results of ref. concerning the one-point function $`<e^{ia\varphi }>`$ in the SG model one gets :
$`<\mathrm{cos}(\widehat{\beta }\varphi )>\text{ }\text{limϵ→0}M_{SG}\mathrm{\Gamma }(ϵ)\text{and}\mathrm{\Lambda }{\displaystyle \frac{M_{SG}}{\pi }}\text{for}\widehat{\beta }^2{\displaystyle \frac{1}{2}}.`$ (50)
Using the boson-fermion correspondence, one has the identification $`\overline{\psi }\psi \mathrm{cos}(\widehat{\beta }\varphi )/\pi `$ and $`<\overline{\psi }\psi >\text{ }\text{limϵ→0}\frac{M_{SG}}{\pi }\mathrm{\Gamma }(ϵ)`$. Let us now return to the case $`p^{}=p+1`$ and the limit $`p\mathrm{}`$ in (1). It corresponds to $`\beta ^2\frac{1}{2}`$, i.e. $`\xi 1`$ and similarly to eq. (50) the VEV of the perturbing operator $`\mathrm{\Phi }_{12}^{(1)}\mathrm{\Phi }_{12}^{(2)}`$ also becomes singular. Consequently, the perturbing operator in (1) can be conveniently rewritten in terms of two Dirac fields (with the same mass $`M_f`$) of two independent massive-Thirring models at the free fermion point :
$`\lambda \mathrm{\Phi }_{12}^{(1)}\mathrm{\Phi }_{12}^{(2)}M_f\overline{\psi }\psi M_f\overline{\psi }^{}\psi ^{}.`$ (51)
In fact, this situation is not really surprising : when $`p\mathrm{}`$ the coset algebra of each $`_{p/p+1}`$ reduces to a level-1 $`SU(2)`$ current algebra. The two models are coupled via their primary fields in the spin $`\frac{1}{2}`$ representation of dimension $`\frac{1}{4}`$ . As expected, in this limit the model becomes the free field theory of two complex massive fermions with undeformed $`SO(4)`$ symmetry.
Finally, let us consider the case $`p^{}=p+1`$ and $`p=2`$ in the action (1). Each minimal model is then identified to $`_{2/3}`$ with central charge $`c=0`$. For each model, the unitary representation is the vacuum one with conformal weight $`\mathrm{\Delta }_{11}=\mathrm{\Delta }_{12}=0`$, for which cases we have the identifications $`\mathrm{\Phi }_{12}^{(i)}=𝕀^{(i)}`$ for $`i\{1,2\}`$. This corresponds to the choice $`\beta ^2=\frac{1}{3}`$, i.e. $`\xi =\frac{1}{2}`$ and, using eqs. (33), (34), (38), one can check that
$`<0|\mathrm{\Phi }_{12}^{(1)}\mathrm{\Phi }_{12}^{(2)}|0>=1\text{and}f_{12}=\lambda ,`$ (52)
as expected.
Let us now turn to the model associated with action (2). Due to eq. (11) it corresponds to the choice $`\beta ^2=\beta _{}^2`$. The condition $`4p>3p^{}`$ guarantees that the perturbing operator is relevant. Then, the vacuum structure is expected to be similar to that of (1). From eqs. (30) and (31) and the substitutions :
$`pp^{},(r,r^{})(s,s^{}),\xi {\displaystyle \frac{1+\xi }{3\xi 1}}`$ (53)
the result for the VEV immediatly follows :
$`<0_{j\stackrel{~}{j}}|\mathrm{\Phi }_{rs}^{(1)}(x)\mathrm{\Phi }_{r^{}s^{}}^{(2)}(x)|0_{j\stackrel{~}{j}}>`$ $`=`$ $`\stackrel{~}{d}_{rs,r^{}s^{}}^{j\stackrel{~}{j}}\left[{\displaystyle \frac{\pi \widehat{\lambda }\gamma (\frac{3\xi 1}{4\xi })(2\xi )^{\frac{5\xi 3}{2\xi }}}{\gamma (\frac{1}{\xi })\gamma (\frac{\xi 1}{2\xi })}}\right]^{\frac{4\xi }{5\xi 3}(\mathrm{\Delta }_{rs}+\mathrm{\Delta }_{r^{}s^{}})}`$ (54)
$`\times \mathrm{exp}𝒬_{21}((1+\xi )r2\xi s,(1+\xi )r^{}2\xi s^{}).`$
The function $`𝒬_{21}(\theta ,\theta ^{})`$ is given by the integral :
$`𝒬_{21}(\theta ,\theta ^{})={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t}}\left({\displaystyle \frac{\mathrm{\Psi }_{21}(\theta ,\theta ^{},t)}{\mathrm{sinh}((1\xi )t)\mathrm{sinh}(2t\xi )\mathrm{sinh}((35\xi )t)}}{\displaystyle \frac{\theta ^2+\theta _{}^{}{}_{}{}^{2}2(1\xi )^2}{4\xi (\xi +1)}}e^{2t}\right)`$
with
$`\mathrm{\Psi }_{21}(\theta ,\theta ^{},t)`$ $`=`$ $`\left[\mathrm{cosh}((\theta +\theta ^{})t)\mathrm{cosh}((\theta \theta ^{})t)\mathrm{cosh}((22\xi )t)\right]`$
$`\times \mathrm{sinh}((1\xi )t)\mathrm{cosh}((35\xi )t)`$
$`\left[\mathrm{cosh}((\theta +\theta ^{})t)+\mathrm{cosh}((\theta \theta ^{})t)\mathrm{cosh}((22\xi )t)1\right]`$
$`\times \mathrm{sinh}((3\xi 1)t/2)\mathrm{cosh}((1+\xi )t/2)`$
and defined by analytic continuation outside this domain. The prefactor associated with the QG restriction $`\stackrel{~}{d}_{rs,r^{}s^{}}^{j\stackrel{~}{j}}=d_{sr,s^{}r^{}}^{j\stackrel{~}{j}}|_{pp^{}}`$. For $`(r,s)=(r^{},s^{})=(2,1)`$ (54) becomes :
$`<0_{j\stackrel{~}{j}}|\mathrm{\Phi }_{21}^{(1)}(x)\mathrm{\Phi }_{21}^{(2)}(x)|0_{j\stackrel{~}{j}}>`$ $`=`$ $`{\displaystyle \frac{1}{\widehat{\lambda }}}\left[{\displaystyle \frac{\pi \widehat{\lambda }\gamma (\frac{3\xi 1}{4\xi })}{\gamma (\frac{1}{\xi })\gamma (\frac{\xi 1}{2\xi })}}\right]^{\frac{4\xi }{5\xi 3}}{\displaystyle \frac{2^{\frac{53\xi }{5\xi 3}}(4\xi )}{\pi (35\xi )}}{\displaystyle \frac{\gamma (\frac{1+\xi }{10\xi 6})\gamma (\frac{3\xi 1}{10\xi 6})}{\gamma (\frac{2\xi }{5\xi 3})}}`$ (55)
with the relation between the mass of the lightest kink $`M`$ and $`\widehat{\lambda }`$ :
$`M={\displaystyle \frac{2^{\frac{1+\xi }{5\xi 3}}\mathrm{\Gamma }(\frac{1+\xi }{10\xi 6})\mathrm{\Gamma }(\frac{3\xi 1}{10\xi 6})}{\pi \mathrm{\Gamma }(\frac{2\xi }{5\xi 3})}}\left[{\displaystyle \frac{\pi \widehat{\lambda }\gamma (\frac{3\xi 1}{4\xi })}{\gamma (\frac{1}{\xi })\gamma (\frac{\xi 1}{2\xi })}}\right]^{\frac{2\xi }{5\xi 3}}.`$ (56)
For the coupled minimal models defined by (2), the massive phase corresponds to the domain :
$`(iii){\displaystyle \frac{3}{5}}<\xi <1,\lambda <0\text{i.e.}{\displaystyle \frac{3}{4}}<{\displaystyle \frac{p}{p^{}}}<1.`$ (57)
One also obtains the bulk free energy associated with action (2) :
$`f_{21}={\displaystyle \frac{M^2\mathrm{sin}(\frac{\pi (3\xi 1)}{10\xi 6})}{2}}{\displaystyle \frac{\mathrm{sin}(\frac{\pi (1+\xi )}{10\xi 6})}{\mathrm{sin}(\frac{\pi (2\xi )}{5\xi 3})}}.`$ (58)
## 4 Application and concluding remarks
Accepting the conjectures (31) and (54), one can easily deduce interesting predictions <sup>6</sup><sup>6</sup>6The same method was applied in to obtain a prediction about the long distance asymtotic of two-point correlation function $`<\sigma _0\sigma _n>_n\mathrm{}`$ in the XXZ spin chain. for the short and long distance asymptotic of two-point correlation functions in the model (1) or (2). For each case depicted in figure 1, we can express the result in terms of :
$`<0_{j\stackrel{~}{j}}|\mathrm{\Phi }_{rs}^{(1)}(x)\mathrm{\Phi }_{r^{}s^{}}^{(2)}(x)|0_{j\stackrel{~}{j}}>=𝒢_{rs,r^{}s^{}}^{j\stackrel{~}{j}}`$ (59)
given by eq. (31) or (54), respectively. First, using the short distance approximation (operator product expansion) and eq.(4) we get $`e^{ia\phi _1(x)}e^{ib\phi _2(y)}e^{i(a\phi _1(x)+b\phi _2(y))}`$. Then, using the Coulomb gas representation of each primary operator (14) one has
<0jj~|Φrs(1)(x)Φrs(2)(y)|0jj~>
|x-y|→0
<0jj~|Φrs(1)(x)Φrs(2)(x)|0jj~>.quantum-operator-productsubscript0𝑗~𝑗subscriptsuperscriptΦ1𝑟𝑠𝑥subscriptsuperscriptΦ2superscript𝑟superscript𝑠𝑦subscript0𝑗~𝑗
|x-y|→0
quantum-operator-productsubscript0𝑗~𝑗subscriptsuperscriptΦ1𝑟𝑠𝑥subscriptsuperscriptΦ2superscript𝑟superscript𝑠𝑥subscript0𝑗~𝑗\displaystyle<0_{j{\tilde{j}}}|\Phi^{(1)}_{rs}(x)\Phi^{(2)}_{r^{\prime}s^{\prime}}(y)|0_{j{\tilde{j}}}>\ \ \raisebox{-6.544pt}{~{}\shortstack{ $\longrightarrow$ \\
${\vspace{-0.2cm}{}_{|x-y|\rightarrow 0}}$}}\ \ <0_{j{\tilde{j}}}|\Phi^{(1)}_{rs}(x)\Phi^{(2)}_{r^{\prime}s^{\prime}}(x)|0_{j{\tilde{j}}}>. (60)
In this limit - case $`(a)`$ \- , the two-point functions then become :
<0jj~|Φrs(1)(x)Φrs(2)(y)|0jj~>
|x-y|→0
𝒢rs,rsjj~.quantum-operator-productsubscript0𝑗~𝑗subscriptsuperscriptΦ1𝑟𝑠𝑥subscriptsuperscriptΦ2superscript𝑟superscript𝑠𝑦subscript0𝑗~𝑗
|x-y|→0
superscriptsubscript𝒢𝑟𝑠superscript𝑟superscript𝑠𝑗~𝑗\displaystyle<0_{j{\tilde{j}}}|\Phi^{(1)}_{rs}(x)\Phi^{(2)}_{r^{\prime}s^{\prime}}(y)|0_{j{\tilde{j}}}>\ \raisebox{-6.544pt}{~{}\shortstack{ $\longrightarrow$ \\
${\vspace{-0.2cm}{}_{|x-y|\rightarrow 0}}$}}\ {\cal G}_{rs,r^{\prime}s^{\prime}}^{j{\tilde{j}}}.\ \ (61)
Secondly, in the long distance approximation the asymptotic two-point function simply reduces to the product of two one-point functions as $`<e^{i𝒂.𝝋(x)}e^{i𝒃.𝝋(y)}><e^{i𝒂.𝝋(x)}><e^{i𝒃.𝝋(y)}>`$ when $`|xy|\mathrm{}`$. Then we obtain - case $`(b)`$ and $`(d)`$ -
<0jj~|Φrs(1)(x)Φrs(2)(y)|0jj~>
|x-y|→∞
<0jj~|Φrs(1)(x)𝕀(2)|0jj~><0jj~|𝕀(1)Φrs(j)(x)|0jj~>;quantum-operator-productsubscript0𝑗~𝑗subscriptsuperscriptΦ1𝑟𝑠𝑥subscriptsuperscriptΦ2superscript𝑟superscript𝑠𝑦subscript0𝑗~𝑗
|x-y|→∞
quantum-operator-productsubscript0𝑗~𝑗subscriptsuperscriptΦ1𝑟𝑠𝑥superscript𝕀2subscript0𝑗~𝑗quantum-operator-productsubscript0𝑗~𝑗superscript𝕀1subscriptsuperscriptΦ𝑗superscript𝑟superscript𝑠𝑥subscript0𝑗~𝑗\displaystyle<0_{j{\tilde{j}}}|\Phi^{(1)}_{rs}(x)\Phi^{(2)}_{r^{\prime}s^{\prime}}(y)|0_{j{\tilde{j}}}>\ \ \raisebox{-6.544pt}{~{}\shortstack{ $\longrightarrow$ \\
${\vspace{-0.2cm}{}_{|x-y|\rightarrow\infty}}$}}\ \ <0_{j{\tilde{j}}}|\Phi^{(1)}_{rs}(x){\mathbb{I}}^{(2)}|0_{j{\tilde{j}}}><0_{j{\tilde{j}}}|{\mathbb{I}}^{(1)}\Phi^{(j)}_{r^{\prime}s^{\prime}}(x)|0_{j{\tilde{j}}}>;
<0jj~|Φrs(1)(x)Φrs(1)(y)|0jj~>
|x-y|→∞
<0jj~|Φrs(1)(x)𝕀(2)|0jj~><0jj~|Φrs(1)(x)𝕀(2)|0jj~>.quantum-operator-productsubscript0𝑗~𝑗subscriptsuperscriptΦ1𝑟𝑠𝑥subscriptsuperscriptΦ1superscript𝑟superscript𝑠𝑦subscript0𝑗~𝑗
|x-y|→∞
quantum-operator-productsubscript0𝑗~𝑗subscriptsuperscriptΦ1𝑟𝑠𝑥superscript𝕀2subscript0𝑗~𝑗quantum-operator-productsubscript0𝑗~𝑗subscriptsuperscriptΦ1superscript𝑟superscript𝑠𝑥superscript𝕀2subscript0𝑗~𝑗\displaystyle<0_{j{\tilde{j}}}|\Phi^{(1)}_{rs}(x)\Phi^{(1)}_{r^{\prime}s^{\prime}}(y)|0_{j{\tilde{j}}}>\ \ \raisebox{-6.544pt}{~{}\shortstack{ $\longrightarrow$ \\
${\vspace{-0.2cm}{}_{|x-y|\rightarrow\infty}}$}}\ \ <0_{j{\tilde{j}}}|\Phi^{(1)}_{rs}(x){\mathbb{I}}^{(2)}|0_{j{\tilde{j}}}><0_{j{\tilde{j}}}|\Phi^{(1)}_{r^{\prime}s^{\prime}}(x){\mathbb{I}}^{(2)}|0_{j{\tilde{j}}}>.
Indeed, it gives :
<0jj~|Φrs(1)(x)Φrs(i)(y)|0jj~>
|x-y|→∞
𝒢rs,11jj~[δi1𝒢rs,11jj~+δi2𝒢11,rsjj~]fori{1,2}quantum-operator-productsubscript0𝑗~𝑗subscriptsuperscriptΦ1𝑟𝑠𝑥subscriptsuperscriptΦ𝑖superscript𝑟superscript𝑠𝑦subscript0𝑗~𝑗
|x-y|→∞
superscriptsubscript𝒢𝑟𝑠11𝑗~𝑗delimited-[]subscript𝛿𝑖1superscriptsubscript𝒢superscript𝑟superscript𝑠11𝑗~𝑗subscript𝛿𝑖2superscriptsubscript𝒢11superscript𝑟superscript𝑠𝑗~𝑗for𝑖12\displaystyle<0_{j{\tilde{j}}}|\Phi^{(1)}_{rs}(x)\Phi^{(i)}_{r^{\prime}s^{\prime}}(y)|0_{j{\tilde{j}}}>\ \raisebox{-6.544pt}{~{}\shortstack{ $\longrightarrow$ \\
${\vspace{-0.2cm}{}_{|x-y|\rightarrow\infty}}$}}\ {\cal G}_{rs,11}^{j{\tilde{j}}}\big{[}\delta_{i1}{\cal G}_{r^{\prime}s^{\prime},11}^{j{\tilde{j}}}+\delta_{i2}{\cal G}_{11,r^{\prime}s^{\prime}}^{j{\tilde{j}}}\big{]}\ \ \ \mbox{for}\ \ i\in\{1,2\} (62)
with $`\delta _{ii^{}}`$ the Krönecker symbol. Obsviously similar results are obtained for the two-point function $`<0_{j\stackrel{~}{j}}|\mathrm{\Phi }_{rs}^{(i)}(x)\mathrm{\Phi }_{r^{}s^{}}^{(2)}(y)|0_{j\stackrel{~}{j}}>`$ using the $`_2`$ symmetry ($`12`$).
Finally, using the fusion rules in the short distance approximation (the two primary fields belong to the same space of states) $`\mathrm{\Phi }_{rs}^{(i)}\times \mathrm{\Phi }_{r^{}s^{}}^{(i)}`$ $``$ $`\mathrm{\Phi }_{r^{\prime \prime }s^{\prime \prime }}^{(i)}`$, the two-point function is expanded in terms of the one-point functions $`<0_{j\stackrel{~}{j}}|\mathrm{\Phi }_{r^{\prime \prime }s^{\prime \prime }}^{(1)}(x)𝕀^{(2)}|0_{j\stackrel{~}{j}}>`$ or $`<0_{j\stackrel{~}{j}}|𝕀^{(1)}\mathrm{\Phi }_{r^{\prime \prime }s^{\prime \prime }}^{(2)}(x)|0_{j\stackrel{~}{j}}>`$ for $`i\{1,2\}`$ respectively. Then, we have for $`i=1`$ \- case $`(c)`$
<0jj~|Φrs(1)(x)Φrs(1)(y)|0jj~>
|x-y|→0
r′′s′′Crs,rsr′′s′′|xy|2(Δr′′s′′ΔrsΔrs)𝒢r′′s′′,11jj~quantum-operator-productsubscript0𝑗~𝑗subscriptsuperscriptΦ1𝑟𝑠𝑥subscriptsuperscriptΦ1superscript𝑟superscript𝑠𝑦subscript0𝑗~𝑗
|x-y|→0
subscriptsuperscript𝑟′′superscript𝑠′′superscriptsubscript𝐶𝑟𝑠superscript𝑟superscript𝑠superscript𝑟′′superscript𝑠′′superscript𝑥𝑦2subscriptΔsuperscript𝑟′′superscript𝑠′′subscriptΔ𝑟𝑠subscriptΔsuperscript𝑟superscript𝑠superscriptsubscript𝒢superscript𝑟′′superscript𝑠′′11𝑗~𝑗\displaystyle<0_{j{\tilde{j}}}|\Phi^{(1)}_{rs}(x)\Phi^{(1)}_{r^{\prime}s^{\prime}}(y)|0_{j{\tilde{j}}}>\ \raisebox{-6.544pt}{~{}\shortstack{ $\longrightarrow$ \\
${\vspace{-0.2cm}{}_{|x-y|\rightarrow 0}}$}}\ \sum_{r^{\prime\prime}s^{\prime\prime}}C_{rs,r^{\prime}s^{\prime}}^{r^{\prime\prime}s^{\prime\prime}}|x-y|^{2(\Delta_{r^{\prime\prime}s^{\prime\prime}}-\Delta_{rs}-\Delta_{r^{\prime}s^{\prime}})}{\cal G}_{r^{\prime\prime}s^{\prime\prime},11}^{{j{\tilde{j}}}} (63)
and similarly for $`i=2`$ where $`C_{rs,r^{}s^{}}^{r^{\prime \prime }s^{\prime \prime }}`$ are the structure constants of the minimal model operator algebra .
Example 1 : Two magnetically coupled Ising models
It is now straightforward to compute different two-point correlation functions in one of the simplest (non-trivial) cases : two-magnetically coupled Ising models in the massive phase. It corresponds to $`\beta ^2=3/8`$ i.e. $`\xi =3/5`$ in (31). In this case we have the identification $`\mathrm{\Phi }_{12}^{(i)}=\sigma ^{(i)}`$ with conformal dimension $`\mathrm{\Delta }_\sigma =1/16`$ \- the spin operator - and $`\mathrm{\Phi }_{13}^{(i)}=ϵ^{(i)}`$ with conformal dimension $`\mathrm{\Delta }_ϵ=1/2`$ \- the energy operator. There are two degenerate ground states denoted $`|0_{00}>|>`$ and $`|0_{1/\mathrm{2\; 1}/2}>|+>`$. For simplicity, we write $`<\pm |\mathrm{}|\pm ><\mathrm{}>_\pm `$. Sometimes, the reflection relation :
$`𝒢(\eta _1,\eta _2)=S_L(\eta _2)𝒢(\eta _1,\eta _2+2\beta 1/\beta )`$ (64)
with
$`S_L(\eta _2)=\left[\pi \mu ^{}\gamma (2\beta ^2)\right]^{\frac{\eta _2\beta +\beta ^21/2}{\beta ^2}}\times {\displaystyle \frac{\mathrm{\Gamma }(2\eta _2\beta +2\beta ^2)\mathrm{\Gamma }(\eta _2/\beta +1/2\beta ^2)}{\mathrm{\Gamma }(2+2\eta _2\beta 2\beta ^2)\mathrm{\Gamma }(2\eta _2/\beta 1/2\beta ^2)}}`$
is useful for analytic continuation. Using eq. (31) and $`d_{11,1s}^\pm =d_{1s,11}^\pm `$ for all $`s`$ we obtain for instance :
$`<\sigma ^{(i)}>_\pm `$ $`=`$ $`𝒢_{11,12}^\pm =\pm 1.297197220\mathrm{}(\lambda )^{1/14};`$
$`<ϵ^{(i)}>_\pm `$ $`=`$ $`𝒢_{11,13}^\pm =2.278284275\mathrm{}(\lambda )^{4/7};`$
$`<\sigma ^{(1)}(0)\sigma ^{(2)}(0)>_\pm `$ $`=`$ $`𝒢_{12,12}^\pm =1.698928047\mathrm{}(\lambda )^{1/7};`$
$`<\sigma ^{(1)}(0)\sigma ^{(2)}(\mathrm{})>_\pm `$ $`=`$ $`(𝒢_{11,12}^\pm )^2=1.682720628\mathrm{}(\lambda )^{1/7};`$
$`<\sigma ^{(1)}(0)ϵ^{(2)}(0)>_\pm `$ $`=`$ $`𝒢_{12,13}^\pm =3.311880669\mathrm{}(\lambda )^{9/14};`$
$`<\sigma ^{(1)}(0)ϵ^{(2)}(\mathrm{})>_\pm `$ $`=`$ $`𝒢_{11,12}^\pm 𝒢_{11,13}^\pm =2.955384028\mathrm{}(\lambda )^{9/14};`$
$`<ϵ^{(1)}(0)ϵ^{(2)}(\mathrm{})>_\pm `$ $`=`$ $`(𝒢_{11,13}^\pm )^2=5.160349412\mathrm{}(\lambda )^{8/7};`$
where the parameter $`\lambda `$ is related to the mass of the lowest kink by :
$`\lambda =0.2379062104\mathrm{}M^{7/4}.`$ (65)
Notice that $`<\sigma ^{(1)}(0)\sigma ^{(2)}(0)>_\pm `$ and $`<\sigma ^{(1)}(0)\sigma ^{(2)}(\mathrm{})>_\pm `$ differ by less than $`0.7\%`$ as expected <sup>7</sup><sup>7</sup>7For instance, it is known that form factors are able to reproduce with high accuracy the UV behaviour of the correlation functions ..
Example 2 : Two energy-energy coupled tricritical Ising models
The case $`p=4`$, $`p^{}=5`$ in (1) describes two tricritical Ising models which interact through their leading energy density operators $`\mathrm{\Phi }_{12}^{(i)}=ϵ^{(i)}`$ of conformal dimension $`\mathrm{\Delta }_ϵ=1/10`$. It corresponds to $`\beta ^2=2/5`$ i.e. $`\xi =2/3`$ in (31). Beside $`ϵ^{(i)}`$ and the identity operator, each minimal model contains the sub-leading energy density operator $`\mathrm{\Phi }_{13}^{(i)}=ϵ_{}^{}{}_{}{}^{(i)}`$ with $`\mathrm{\Delta }_ϵ^{}=3/5`$ (“vacancy operator”) , two magnetic operators $`\mathrm{\Phi }_{22}^{(i)}=\sigma ^{(i)}`$ with $`\mathrm{\Delta }_\sigma =3/80`$, $`\mathrm{\Phi }_{21}^{(i)}=\sigma _{}^{}{}_{}{}^{(i)}`$ with $`\mathrm{\Delta }_\sigma ^{}=7/16`$ and $`\mathrm{\Phi }_{14}^{(i)}`$. Due to the obvious property $`d_{rs,r^{}s^{}}^{j\stackrel{~}{j}}=d_{11,r^{}s^{}}^{j\stackrel{~}{j}}d_{rs,11}^{j\stackrel{~}{j}}`$, similarly to the previous case we find for instance for any vacuum $`|j\stackrel{~}{j}>`$ :
$`<\sigma ^{(1)}>_{j\stackrel{~}{j}}`$ $`=`$ $`𝒢_{22,11}^{j\stackrel{~}{j}}=d_{22,11}^{j\stackrel{~}{j}}\times 1.144656674\mathrm{}(\lambda )^{3/64};`$
$`<ϵ^{(1)}>_{j\stackrel{~}{j}}`$ $`=`$ $`𝒢_{12,11}^{j\stackrel{~}{j}}=d_{12,11}^{j\stackrel{~}{j}}\times 1.529866659\mathrm{}(\lambda )^{1/8};`$
$`<\sigma ^{(1)}(0)\sigma ^{(2)}(0)>_{j\stackrel{~}{j}}`$ $`=`$ $`𝒢_{22,22}^{j\stackrel{~}{j}}=d_{22,22}^{j\stackrel{~}{j}}\times 1.315726811\mathrm{}(\lambda )^{3/32};`$
$`<\sigma ^{(1)}(0)\sigma ^{(2)}(\mathrm{})>_{j\stackrel{~}{j}}`$ $`=`$ $`𝒢_{22,11}^{j\stackrel{~}{j}}𝒢_{11,22}^{j\stackrel{~}{j}}=d_{22,22}^{j\stackrel{~}{j}}\times 1.310238901\mathrm{}(\lambda )^{3/32};`$
$`<ϵ^{(1)}(0)ϵ^{(2)}(0)>_{j\stackrel{~}{j}}`$ $`=`$ $`𝒢_{12,12}^{j\stackrel{~}{j}}=d_{12,12}^{j\stackrel{~}{j}}\times 2.419476973\mathrm{}(\lambda )^{1/4};`$
$`<ϵ^{(1)}(0)ϵ^{(2)}(\mathrm{})>_{j\stackrel{~}{j}}`$ $`=`$ $`𝒢_{12,11}^{j\stackrel{~}{j}}𝒢_{11,12}^{j\stackrel{~}{j}}=d_{12,12}^{j\stackrel{~}{j}}\times 2.340491994\mathrm{}(\lambda )^{1/4};`$
where the parameter $`\lambda `$ is related to the mass of the lowest kink by :
$`\lambda =0.2566343706\mathrm{}M^{8/5}.`$ (66)
Notice that $`<\sigma ^{(1)}(0)\sigma ^{(2)}(0)>_{j\stackrel{~}{j}}`$ and $`<\sigma ^{(1)}(0)\sigma ^{(2)}(\mathrm{})>_{j\stackrel{~}{j}}`$$`<ϵ^{(1)}(0)ϵ^{(2)}(0)>_{j\stackrel{~}{j}}`$ and $`<ϵ^{(1)}(0)ϵ^{(2)}(\mathrm{})>_{j\stackrel{~}{j}}`$ differ by less than $`0.5\%`$ and $`3\%`$ respectively.
Example 3 : Two coupled $`A_5`$ RSOS models
The case $`p=5`$, $`p^{}=6`$ in (1) describes two $`A_5`$ RSOS models coupled by their primary operators $`\mathrm{\Phi }_{12}^{(i)}`$ with conformal dimension $`\mathrm{\Delta }_{12}=1/8`$. It corresponds to $`\beta ^2=5/12`$ i.e. $`\xi =5/7`$ in (31). Each minimal model also contains the primary operator $`\mathrm{\Phi }_{22}^{(i)}`$ with $`\mathrm{\Delta }_{22}=1/40`$. As before we obtain :
$`<\mathrm{\Phi }_{22}^{(1)}>_{j\stackrel{~}{j}}`$ $`=`$ $`𝒢_{22,11}^{j\stackrel{~}{j}}=d_{22,11}^{j\stackrel{~}{j}}\times 1.090446894\mathrm{}(\lambda )^{1/30};`$
$`<\mathrm{\Phi }_{12}^{(1)}>_{j\stackrel{~}{j}}`$ $`=`$ $`𝒢_{12,11}^{j\stackrel{~}{j}}=d_{12,11}^{j\stackrel{~}{j}}\times 1.726352342\mathrm{}(\lambda )^{1/6};`$
$`<\mathrm{\Phi }_{22}^{(1)}(0)\mathrm{\Phi }_{22}^{(2)}(0)>_{j\stackrel{~}{j}}`$ $`=`$ $`𝒢_{22,22}^{j\stackrel{~}{j}}=d_{22,22}^{j\stackrel{~}{j}}\times 1.191588988\mathrm{}(\lambda )^{1/15};`$
$`<\mathrm{\Phi }_{22}^{(1)}(0)\mathrm{\Phi }_{22}^{(2)}(\mathrm{})>_{j\stackrel{~}{j}}`$ $`=`$ $`𝒢_{22,11}^{j\stackrel{~}{j}}𝒢_{11,22}^{j\stackrel{~}{j}}=d_{22,22}^{j\stackrel{~}{j}}\times 1.189074429\mathrm{}(\lambda )^{1/15};`$
where the parameter $`\lambda `$ is related to the mass of the lowest kink by :
$`\lambda =0.2697511940\mathrm{}M^{3/2}.`$ (67)
Notice that $`<\mathrm{\Phi }_{22}^{(1)}(0)\mathrm{\Phi }_{22}^{(2)}(0)>_{j\stackrel{~}{j}}`$ and $`<\mathrm{\Phi }_{22}^{(1)}(0)\mathrm{\Phi }_{22}^{(2)}(\mathrm{})>_{j\stackrel{~}{j}}`$ differ by $`2\%`$.
Example 4 : Two energy-energy coupled 3-state Potts models
The case $`p=5`$, $`p^{}=6`$ in (2) describes two 3-state Potts models coupled by their energy density operator $`\mathrm{\Phi }_{21}^{(i)}=ϵ^{(i)}`$ with conformal dimension $`\mathrm{\Delta }_{21}=2/5`$. It corresponds to $`\xi =5/7`$ in (54). Each minimal model also contains the primary operator $`\mathrm{\Phi }_{23}^{(i)}=\sigma ^{(i)}`$ \- the spin operator - with $`\mathrm{\Delta }_{23}=1/15`$. We obtain for instance :
$`<\sigma ^{(1)}>_{j\stackrel{~}{j}}`$ $`=`$ $`𝒢_{23,11}^{j\stackrel{~}{j}}=\stackrel{~}{d}_{23,11}^{j\stackrel{~}{j}}\times 1.9079\mathrm{}(\widehat{\lambda })^{1/3};`$
$`<\sigma ^{(1)}(0)\sigma ^{(2)}(0)>_{j\stackrel{~}{j}}`$ $`=`$ $`𝒢_{23,23}^{j\stackrel{~}{j}}=\stackrel{~}{d}_{23,23}^{j\stackrel{~}{j}}\times 4.50\mathrm{}(\widehat{\lambda })^{2/3};`$
$`<\sigma ^{(1)}(0)\sigma ^{(2)}(\mathrm{})>_{j\stackrel{~}{j}}`$ $`=`$ $`𝒢_{23,11}^{j\stackrel{~}{j}}𝒢_{11,23}^{j\stackrel{~}{j}}=\stackrel{~}{d}_{23,23}^{j\stackrel{~}{j}}\times 3.64\mathrm{}(\widehat{\lambda })^{2/3};`$
where the parameter $`\widehat{\lambda }`$ is related to the mass of the lowest kink by :
$`\widehat{\lambda }=0.2612863655\mathrm{}M^{2/5}.`$ (68)
To conclude, we would like to mention that here we studied only a special case of a much more general class of integrable coupled models . There exist many other examples which can be worked out along the same lines. Let us also note that the exact form factor techniques as well as the truncated conformal space approach may be used and similar numerical analyses can be performed for the correlation functions.
### Aknowledgements
I am grateful to C. Ahn, G. Delfino, G. Delius, D. Reynaud, P. Simon, C. Kim, C. Rim, L. Zhao and especially to V.A. Fateev and E. Corrigan for useful discussions. Work supported by the EU under contract ERBFMRX CT960012 and Marie Curie fellowship HPMF-CT-1999-00094.
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# ON THE THERMAL INSTABILITY IN A CONTRACTING GAS CLOUD AND FORMATION OF A BOUND CLUSTER
## 1 Introduction
The Galactic globular clusters are suggested to be fossil record of the Galaxy formation era (e.g., Fall $`\&`$ Rees 1985). So any knowledge about their formation mechanism will be helpful to understand the protogalactic environment. Young counterparts of globular clusters are absent in the Milky Way, whereas in galaxies such as the LMC, SMC and M33, young globulars exist (e.g., Kumai, Basu, $`\&`$ Fujimoto 1993; Efremov 1995). Such galaxies are in many case ongoing mergers and merger remnants or starburst galaxies, but young globulars also exist in normal spiral, dwarf and irregular galaxies (Schweizer 1999). The conditions that enable globular cluster formation in these galaxies may be also realized at the epoch of the Galactic globular cluster formation. Since globular clusters are gravitationally bound, we attempt to clarify the condition of bound cluster formation.
At the formation epoch of the globular clusters, high star formation efficiency (SFE) must be achieved somehow, since high SFE is required in order to form a bound cluster. The threshold SFE value for bound cluster formation is 0.5 in the case of rapid gas removal (e.g., Hills 1980), which is significantly greater than average SFE ($``$ 0.01) in the giant molecular clouds (GMCs). Thus, to understand the conditions of bound cluster formation, we should know in what environment the SFE can be high. In regard of this problem, Elmegreen $`\&`$ Efremov (1997) claimed that high-pressure environment is required for bound cluster formation because disruption of the cluster-forming cloud by massive stars is difficult in such an environment. They also argued that bound cluster formation is efficient in a low-metallicity environment, as it is difficult to disrupt the cloud by radiation pressure of massive stars. Here, considering the formation process of dense clouds, we point out that high SFE and efficient bound cluster formation can be achieved in low-metallicity and/or strong-radiation environments because of cloud fragmentation via thermal instability.
Molecular clouds, where stellar clusters form, usually belong to some larger structure such as superclouds (Elmegreen $`\&`$ Elmegreen 1983), the largest cloud complexes in galaxies. Due to the similarity between the masses of superclouds and the critically unstable Jeans mass, large-scale gravitational instability is considered to produce these superclouds (Elmegreen 1987). Warm H I gas occupies large fraction of these superclouds, since roughly half of the H I in the Milky Way disk is in the form of warm gas (Kulkarni $`\&`$ Heiles 1988). When the gas disk becomes gravitationally unstable and fragments into superclouds, this warm gas should be converted into cold H I gas since the gas becomes thermally unstable when the density is somewhat higher. From this cold H I gas, GMCs and, subsequently, stellar clusters are expected to form. Similarly, GMC formation in a supercloud via thermal instability was also considered by Kolesnik (1991). However, his method was based on density structure of static isothermal cloud and our treatment is different.
For modeling the above-mentioned top-down scenario of cluster formation from a supercloud (see also Efremov 1995), we consider a large, warm gas cloud ($`T10^4\mathrm{K}`$) contracting because of its self-gravity. As the cloud density increases, it becomes thermally unstable and breaks up into cool, dense condensations. We investigate this process based on linear analysis of thermal instability in the cloud. For the evolution of condensations, we adopt the scenario that these condensations collide each other and self-gravitating clumps are formed in the cloud, and then stellar clusters form. To judge whether the cluster is bound or not, we should know the mean density of the cloud at the point of breakup, since we can relate the mean density of the cloud to its velocity dispersion, which determines SFE and hence whether the cluster is bound or unbound. So we calculate the thermal evolution of a contracting gas cloud and estimate the density when the cloud becomes thermally unstable. The effects of metallicity and radiation flux are also examined.
In the following sections, we present the model of the cloud and the instability analysis method. Based on the results of the calculations, we discuss the conditions necessary for high SFE and bound cluster formation.
## 2 Thermal Evolution of a Contracting Gas Cloud
We calculate the evolution of density, temperature, and ionization degree of a contracting gas cloud. For simplicity, we consider a spherical and homogeneous gas cloud. As an initial condition, we take the static cloud filled with warm H I gas ($`T10^4\mathrm{K}`$) in thermal and ionization equilibrium. The contraction is simplified as a free-fall collapse. Then the thermal evolution of the cloud is given by the following basic equations:
$$\frac{d}{dt}\mathrm{ln}n=\left\{24\pi G\rho \left[1\left(\frac{n_i}{n}\right)^{1/3}\right]\right\}^{1/2},$$
(1)
$$\frac{d\epsilon }{dt}=\frac{P}{\rho ^2}\frac{d\rho }{dt},$$
(2)
$$\frac{dx_e}{dt}=X(T,\rho ,x_e).$$
(3)
Here $`\rho `$ and $`n`$ are the mass and number density of hydrogen of the cloud, respectively, and $`n_i`$ is the initial value of $`n`$. $`P`$, $`T`$, $`\epsilon `$, and $`x_e`$ are the pressure, temperature, internal energy, and ionization degree of the cloud, respectively. Function $``$ represents the net cooling rate. As heating sources, we consider interstellar UV radiation, soft X-ray, and cosmic ray. As coolants, we consider hydrogen $`\mathrm{Ly}\alpha `$ cooling, line cooling by O I, C II, Si II, Fe II, dust cooling, and so on. Ionization states of these metal species are fixed according to their ionization potentials. In this work we approximate the cloud as optically thin. Function $`X`$ represents the net ionization rate of hydrogen atoms. We consider soft X-ray and cosmic ray as ionizing sources. To determine the X-ray heating and ionization rate, the hydrogen column density of the cloud, $`N_\mathrm{H}`$, must be specified. we take $`N_\mathrm{H}=10^{19}\mathrm{cm}^3`$ throughout the cloud evolution, though $`N_\mathrm{H}`$ increases as the cloud contracts. This simplification is not serious, since we performed the same calculation including the change of $`N_\mathrm{H}`$ and got almost the same result. For these thermal processes, we refer to Wolfire et al. (1995) and references there in.
To verify the dependence on metallicity and radiation flux, we introduce parameters $`Z`$ and $`X_0`$. In our assumptions, the number density of dust grains and the gas phase abundances of O I, C II, etc. change in proportion to $`Z`$. The UV radiation flux, X-ray flux and cosmic ray flux change proportional to $`X_0`$. Here $`Z=1`$ and $`X_0=1`$ correspond to solar values.
## 3 Formation of Condensations Due to Thermal Instability
We estimate the mass and the formation timescale of condensations using linear approximation. The growth of density perturbation in the parent cloud is governed by the equation
$`\sigma ^3+\sigma ^2c_s\left(k_T+{\displaystyle \frac{k^2}{k_K}}\right)+\sigma c_{s}^{}{}_{}{}^{2}k^2+{\displaystyle \frac{c_{s}^{}{}_{}{}^{3}k^2}{\gamma }}`$
$`\times [k_Tk_\rho +{\displaystyle \frac{k^2}{k_K}}{\displaystyle \frac{\gamma }{c_{s}^{}{}_{}{}^{3}}}(\gamma 1)_0]=0,`$ (4)
where $`k`$ and $`\sigma `$ are the wave number and the growth rate of the perturbation, respectively. This equation is the same as in Field (1965), except in the last term, where it represents nonequilibrium effect. Thus, the net cooling rate $`_0=(T_0,\rho _0,x_{e0})`$ is not necessary zero, where $`T_0`$, $`\rho _0`$, and $`x_{e0}`$ are the background temperature, density, and ionization degree, respectively. The definitions of $`k_\rho `$, $`k_T`$, and $`k_K`$ are the same as in Field (1965):
$`k_\rho ={\displaystyle \frac{\mu \left(\gamma 1\right)\rho _0_\rho }{Rc_sT_0}},k_T={\displaystyle \frac{\mu \left(\gamma 1\right)_T}{Rc_s}},`$
$`k_K={\displaystyle \frac{Rc_s\rho _0}{\mu \left(\gamma 1\right)K_0}},`$ (5)
where $`\mu `$, $`\gamma `$, $`R`$, and $`c_s`$ are the mean molecular weight, ratio of specific heats, gas constant, and sound speed, respectively. We take the value of $`\mu `$ as $`1.4`$ and $`\gamma `$ as 5/3. The sound speed is expressed as $`c_s=\left(\gamma P_0/\rho _0\right)^{1/2}`$, where $`P_0`$ is the background pressure. The symbols $`_\rho `$, $`_T`$, and $`K_0`$ are represented as follows:
$$_\rho =\left(\frac{}{\rho }\right)_T,_T=\left(\frac{}{T}\right)_\rho ,K_0=K\left(T_0\right),$$
(6)
where $`K(T)`$ is thermal conductivity, and we adopt $`K(T)=2.0\times 10^3T^{1/2}`$, the thermal conductivity for neutral hydrogen.
Among three roots of equation (4), two roots correspond to sound wave modes and one root corresponds to isobaric condensation mode (Field 1965). When this equation has unstable isobaric condensation mode, condensations form because of the growth of unstable perturbations (e.g., Goldsmith 1970). This mode has a maximum growth rate of $`\sigma _m`$ at wave number $`k_m`$ (Fig. 1). Among various density perturbations with different intensities and wave numbers in the cloud, the perturbations with wave number $`k_m`$ grow fastest. So we estimate the mass of the condensation, $`M`$, as
$$M=\frac{4\pi }{3}\rho _0\left(\frac{\lambda _m}{2}\right)^3,\lambda _m=\frac{2\pi }{k_m}.$$
(7)
We also estimate the formation timescale as
$$t_g=\frac{1}{\sigma _m}.$$
(8)
Because equation (4) is derived under the assumption that contraction of the cloud is negligible, it must keep in mind that equation (4) is applicable only when the growth timescale of perturbation is shorter than the contraction timescale of the cloud. Also note that our use of equation (4) is justified when $`\lambda _m`$ is much smaller than the parent cloud, because equation (4) is for an infinite homogeneous medium. We can neglect this point, since the parent cloud ($``$ a few kpc ) is much larger than $`\lambda _m`$ ( $``$ a few pc).
## 4 Results
We perform the calculation in the parameter ranges $`Z=1`$-$`10^2`$, $`X_0=10`$-$`0.1`$. For the purpose of illustration, we show the results of the cases ($`Z=1,X_0=1`$) and ($`Z=10^2,X_0=1`$) in Figures 2 and 3. In the figures, evolutionary paths in $`nT`$ plane of both the free-fall collapsing gas cloud and the quasi-statically contracting gas cloud keeping thermal equilibrium are presented. Equations (1)-(3) are used for the free-fall collapsing paths and the equations, $`(T,\rho ,x_e)=0`$ and $`X(T,\rho ,x_e)=0`$, are used for thermal equilibrium paths. The condensation mass $`M`$ and $`t_g/t_{dyn}`$, the ratio of formation timescale of condensation to dynamical timescale of the parent cloud, are also shown in the same figures. Here the dynamical timescale of the parent cloud, $`t_{\mathrm{dyn}}`$, is defined as $`\rho /\left(d\rho /dt\right)`$, where $`\rho `$ is the density of the parent cloud. The contours of $`M`$ and $`t_\mathrm{g}/t_{\mathrm{dyn}}`$ in Figures 2 and 3 are calculated using equation (4), with the assumption of ionization equilibrium (i.e., only $`n`$ and $`T`$ are used for calculation, and $`x_e`$ is a function of $`n`$ and $`T`$). From the figures we can see the tendency that the smaller metallicity, the more the evolutionary path of the free-fall collapsing cloud deviates from the path of the quasi-statically contracting cloud. For the case $`Z=1`$, the two evolutionary paths deviate in temperature $`1000\mathrm{K}`$ at density $`1\mathrm{cm}^3`$. In contrast, for the case $`Z=10^2`$, the two paths deviate $`8000\mathrm{K}`$ at density $`100\mathrm{cm}^3`$. For the paths of the free-fall collapsing gas cloud, dependence on metallicity is stronger than equilibrium paths. Such behavior is explained by the fact that a nonequilibrium state is realized for free-fall cases because the cooling timescale can be long relative to dynamical time and that degree of nonequilibrium is stronger for lower metallicity since cooling time is in inverse proportion to $`Z`$.
As the cloud density exceeds some critical value, the gas cloud becomes thermally unstable. At this point, the growth timescale of perturbation, $`t_g`$, is longer than the dynamical timescale of the cloud, $`t_{\mathrm{dyn}}`$. So the perturbation cannot grow to the condensation, because the parent cloud evolves faster than the perturbation. As the cloud evolves, $`t_g`$ becomes smaller than $`t_{\mathrm{dyn}}`$. Thus, condensation will form at the moment $`t_g=t_{\mathrm{dyn}}`$. For the case $`Z=1`$, condensations form at the density $`1\mathrm{cm}^3`$ (Fig. 2a), whereas for the case $`Z=10^2`$, condensations form at the density $`300\mathrm{cm}^3`$ (Fig. 3a). Condensation mass is a few $`10M_{}`$ at the density $`1\mathrm{cm}^3`$ for $`Z=1`$ (Fig. 2b) and a few $`M_{}`$ at the density $`300\mathrm{cm}^3`$ for $`Z=10^2`$ (Fig 3b).
For the case when $`Z`$ and $`X_0`$ are simultaneously changed, the results are shown in Figure 4. The calculational procedure for Figure 4 differs from Figures 2 and 3, in a point that $`t_g`$ and $`M`$ are calculated using equations (1)-(8) to draw Figure 4, whereas the assumption of ionization equilibrium is adopted to draw the contours of Figures 2 and 3. In Figure 4a, the cloud density $`n_f`$ at the time of condensation formation is shown. For $`X_0=1`$, the density $`n_f`$ increases from $`1\mathrm{c}\mathrm{m}^3`$ at $`Z=1`$ to $`300\mathrm{c}\mathrm{m}^3`$ at $`Z=10^2`$. For $`Z=1`$, $`n_f`$ increases from $`0.1\mathrm{cm}^3`$ at $`X_0=0.1`$ to $`10\mathrm{c}\mathrm{m}^3`$ at $`X_0=10`$. In Figure 4b, the condensation mass is shown. For $`X_0=1`$, the mass $`M`$ decreases from $`30M_{}`$ at $`Z=1`$ to $`1M_{}`$ at $`Z=10^2`$. For $`Z=1`$, $`M`$ decreases from $`300M_{}`$ at $`X_0=0.1`$ to $`1M_{}`$ at $`X_0=10`$. Thus, from Figure 4, we can see that for the smaller metallicity, the cloud density $`n_f`$ is the larger, and the condensation mass $`M`$ is the smaller. We can also see that for the stronger radiation, the density $`n_f`$ is the larger and the mass $`M`$ is the smaller.
Previously, Elmegreen & Parravano (1994) calculated the critical pressure for coexistence of cool and warm gas, changing radiation flux and metallicity. This critical pressure roughly corresponds to $`n_f`$ in this work. Contrary to our results, metallicity dependence of the critical pressure is small in their calculation. They assumed static cloud for the background, whereas we assume free-falling cloud. Then, as noted above, for the smaller metallicity, the evolutionary path of free-fall contracting cloud deviates more from the path of the quasi-statically contracting cloud, because of the reduced cooling rate. On the other hand, for the static cloud, thermal equilibrium is achieved so that metallicity dependence is weak. This is the reason of the difference between our work and their work.
We estimate the density of condensation as follows. After the criterion $`t_g=t_{\mathrm{dyn}}`$ is achieved, the overdensity region (size $``$ $`\lambda _m`$) in the parent cloud grows in pressure equilibrium with surrounding gas, since crossing time is short compared with cooling time for the fastest growing perturbation. Then the overdensity evolves along an isobaric path starting from the point $`t_g=t_{\mathrm{dyn}}`$ (this path is not drawn in Figs.2 and 3). The runaway growth of the overdensity continues until it reaches thermochemical equilibrium, where it is thermally stable (e.g., Parravano 1987). We search the intersection of the isobaric path (the path with $`nT=\mathrm{constant}`$, starting from the point $`t_g=t_{\mathrm{dyn}}`$) with the equilibrium path (Figs.2 and 3, dashed line) and identify it as “condensation”. From the intersection, we can know the density of condensation, $`n_c`$, and the temperature of the condensation, $`T_c`$. We perform the calculation and find that $`n_c`$ is roughly 100 times larger than $`n_f`$ in most parameter ranges. In the above-described procedure, we assume that the intersection is thermally stable and self-gravity of the overdensity can be neglected. We confirm that all intersections in this calculation are thermally stable, so that the first assumption is justified. The second assumption is also justified by the later discussion. Using the values of $`n_c`$ and $`T_c`$, the Jeans mass can be derived, and it is much larger than the condensation mass for most parameter values. Thus, self-gravity of the condensation is negligible and the density inside the condensation is expected to be almost uniform. After formation, condensations accrete surrounding gas and can potentially grow to be self-gravitating. To estimate the mass growth rate of condensation, we should calculate the rate of mass flow across the boundary between the condensation and surrounding gas. According to previous works (e.g., Penston & Brown 1970; Parravano 1987), the growth timescale of the condensation mass is much larger than the free-fall timescale of the parent cloud. So we can neglect the mass growth of condensation through accretion.
## 5 Application to bound cluster formation
After formation of condensations, the parent cloud continues free-fall contraction until collisions between condensations begin. Then the formation of massive self-gravitating clumps through collisional buildup of small condensations is expected. Star formation inside such massive clumps will follow. However, evaluation of such evolution is difficult task. The onset of collisions depends on velocity distribution and space density of condensations, which we cannot know precisely. Subsequent evolution of condensations also depends on velocity and mass distribution of condensations through the collisional process. This collisional process is not so simple (e.g., Hausman 1981) and makes the calculation of evolution complicated. In regard of these difficulties, we adopt a simplified picture of the evolution of condensations. We neglect the random velocity of condensations and assume homologous contraction of the parent cloud composed of condensations. Then the collisions begin when the mean parent cloud density increases to the density of condensation, $`n_c`$, because the parent cloud is almost filled with condensations in such situation. As a result of collisions (we assume that the output of collision is mostly adhesion), massive self-gravitating clumps will form. At this point, the density of massive clumps is larger than the mean parent cloud density $`n_c`$, because of self-gravity. So the parent cloud contract further. We assume collisionless collapse of the system composed of self-gravitating clumps, for simplicity. Then the virialization of the system and the beginning of star formation in the clumps is expected.
It is reasonable to consider that the cluster-forming region is only a part of the parent cloud, not the entire cloud. We define the mass of the region as $`M_{\mathrm{cl}}`$ and take the reference value as $`10^5M_{}`$ according to typical mass of globular clusters. Then virial velocity of the region is
$$c2\times \left(\frac{M_{\mathrm{cl}}}{10^5M_{}}\right)^{1/3}\left(\frac{n_{\mathrm{cl}}}{1\mathrm{c}\mathrm{m}^3}\right)^{1/6}\mathrm{km}\mathrm{s}^1,$$
(9)
where $`n_{\mathrm{cl}}`$ is the mean hydrogen number density of the cluster-forming region. Because (1) we assume collisionless collapse of the parent cloud composed of self-gravitating clumps, where the mean cloud density at the onset of collapse is the same as condensation density $`n_c`$, and (2) condensation density $`n_c`$ is typically 100 times higher than the parent cloud density at the time of condensation formation, $`n_f`$ (see the last paragraph in § 4), $`n_{\mathrm{cl}}`$ and $`n_f`$ are related as
$$n_{\mathrm{cl}}800n_f,$$
(10)
where the relations, $`n_{\mathrm{cl}}8n_c`$ from the point (1) and $`n_c100n_f`$ from the point (2), are used.
Whether the cluster is bound or not depends on SFE, the ratio of stellar mass to the mass of cluster-forming region. And SFE is determined by the interaction between stellar activity and the cluster-forming region. For example, stellar wind, radiation pressure, ionization pressure of the H II region or supernova explosion may blow off remaining gas and lead to lower SFE. In this paper we consider two cases, gas removal by the expanding H II region and supernova explosion.
First we discuss gas removal by expanding H II region. The pressure of H II region corresponds to velocity dispersion $`c_{\mathrm{HII}}10\mathrm{k}\mathrm{m}\mathrm{s}^1`$. If virial velocity of the region, $`c`$, is smaller than $`c_{\mathrm{HII}}`$, the H II region pushes out the surrounding gas and blows it off, and then the cluster will be unbound. In opposite case, the H II region cannot blow off the surrounding gas and high SFE can be achieved; then the cluster will be bound. So the criterion of bound cluster formation is considered as $`c>c_{\mathrm{HII}}`$. This corresponds to the condition
$$n_f>10\times \left(\frac{C}{800}\right)^1\left(\frac{c_{\mathrm{HII}}}{10\mathrm{k}\mathrm{m}\mathrm{s}^1}\right)^6\left(\frac{M_{\mathrm{cl}}}{10^5M_{}}\right)^2\mathrm{cm}^3,$$
(11)
where $`C`$ represents density enhancement, i.e., the ratio of $`n_{\mathrm{cl}}`$ to $`n_f`$. Note that derivation of equation (11) relies on some rough assumptions and depends on uncertain factors such as $`C`$, $`c_{\mathrm{HII}}`$, and $`M_{\mathrm{cl}}`$. Moreover, dependence on $`c_{\mathrm{HII}}`$ is very strong. So we should regard equation (11) as qualitative criterion.
Next we discuss gas removal by supernova (SN) explosion. The proto-globular cluster cloud is disrupted by several to tens of concurrent SN explosions (Dopita & Smith 1986; Morgan & Lake 1989). A high star formation rate at the formation epoch of globular cluster will cause such multiple SN explosions so that relic gas will be blown off. When the star formation time scale $`t_{\mathrm{SF}}`$ is shorter than the lifetime of massive star, $`\tau `$, a large fraction of gas in the cluster-forming region can be converted into stars before SN explosions; then high SFE is achieved and a bound cluster can be formed. In contrast, when $`t_{\mathrm{SF}}`$ is longer than $`\tau `$, SN explosions will remove surrounding gas and stop later star formation. In this case, low SFE and unbound cluster formation is expected (e.g., Yoshii and Arimoto 1987). The star formation timescale $`t_{\mathrm{SF}}`$ is expected as
$$t_{\mathrm{SF}}=At_{ff},$$
(12)
where $`A`$ represents a proportional constant of order 1 and $`t_{ff}=\left(3\pi /32G\rho \right)^{1/2}`$ is free-fall timescale of the cluster-forming region. We adopt the lifetime of massive star, $`\tau `$, as
$$\tau 5\times 10^6\mathrm{yr}.$$
(13)
The star formation timescale $`t_{\mathrm{SF}}`$, evaluated at density $`n_{\mathrm{cl}}`$, is
$$t_{\mathrm{SF}}=At_{ff}1.3\times 10^7A\left(\frac{n_{\mathrm{cl}}}{10^3\mathrm{cm}^3}\right)^{1/2}\mathrm{yr}.$$
(14)
Thus, the condition for high SFE ($`\tau >t_{\mathrm{SF}}`$) in term of $`n_f`$ is
$$n_f>10\times \left(\frac{\tau }{5\times 10^6\mathrm{yr}}\right)^2\left(\frac{A}{1}\right)^2\left(\frac{C}{800}\right)^1\mathrm{cm}^3.$$
(15)
This criterion is qualitative, since it depends some uncertain factors, the same as equation (11).
For both case, $`n_f>10\mathrm{cm}^3`$ is required for high SFE and bound cluster formation. From Figure 4a we can see that the density $`n_f`$ increases as metallicity decreases and also as radiation increases. The condition $`n_f>10\mathrm{cm}^3`$ corresponds to the parameter range $`Z0.05`$ for $`X_0=1`$ and $`X_08`$ for $`Z=1`$. This means high SFE and efficient bound cluster formation are realized in a low-metallicity and/or strong-radiation environment in our model. From Figure 4a, we can also see that when $`Z1`$, $`n_f`$ is more sensitive to $`X_0`$ than $`Z`$ and when $`Z0.1`$, $`n_f`$ is more sensitive to $`Z`$ than $`X_0`$. Thus SFE mainly depends on $`X_0`$ around $`Z1`$ and on $`Z`$ below $`Z0.1`$.
The qualitative behavior of the above result will be independent of the details about evolution of condensations and feedback from star formation, since the essential ingredients are only two points, (1) low-metallicity and/or strong-radiation fields lead to high breakup density of the parent cloud, and (2) high-density environments result in high SFE. However, our conclusion may be changed if we consider nonspherical evolution of the cloud. The dynamical timescale of an initially nonspherical cloud or a cloud with rotation is similar to that of a spherical cloud at first, e.g., $`t_{\mathrm{dyn}}n^{1/2}`$, where $`n`$ is the density of the cloud. But later disklike collapse will follow and the dynamical timescale becomes $`t_{\mathrm{dyn}}n^1`$ (Susa, Uehara, $`\&`$ Nishi 1996). Then the ratio of the cooling time $`t_{\mathrm{cool}}n^1Z^1`$ to the dynamical time is
$$\frac{t_{\mathrm{cool}}}{t_{\mathrm{dyn}}}n^{1/2}Z^1,$$
(16)
for initial spherical-like phase and
$$\frac{t_{\mathrm{cool}}}{t_{\mathrm{dyn}}}Z^1,$$
(17)
for later disklike phase. When the cloud becomes thermally unstable, we can estimate whether or not condensations form, using the value of $`t_{\mathrm{cool}}/t_{\mathrm{dyn}}`$. If this ratio is less than unity, condensation will form, and vice versa. In case the cloud becomes unstable when it is still in a spherical-like phase, condensations will form regardless of the metallicity $`Z`$, as we calculated in this work. This can be explained since the density $`n`$, when the cloud becomes unstable, increases as $`Z`$ decreases (see Figs. 2 and 3), so that $`t_{\mathrm{cool}}/t_{\mathrm{dyn}}`$ does not change so much. But, in case the cloud becomes unstable when it is in a disklike phase, condensations may not form for the low-metallicity case since $`t_{\mathrm{cool}}/t_{\mathrm{dyn}}`$ increases as $`Z`$ decreases. In this case our discussion in this section is not applicable. To summarize, in the presence of rotation or deviation from spherical symmetry, if the cloud becomes thermally unstable in the early phase of collapse, condensations will form and our conclusion in this section is adequate. In contrast, if the cloud becomes thermally unstable in the late phase of collapse, condensations may not form throughout the collapse and another discussion may be required.
Finally, we present some applications. A low-metallicity environment is realized in dwarf galaxies and in the early stage of our Galaxy, for example. And strong-radiation fields are realized in starburst galaxies and probably in the early stage of our Galaxy, where young bulge stars are expected to provide strong radiation. Thus, according to the result, young globular clusters in dwarf galaxies and starburst galaxies may be produced through the mechanism we show in this work. Also, the disk population of globular clusters in our Galaxy (Zinn 1985) may be produced through the same mechanism.
## 6 Summary
In this paper we perform linear analysis of thermal instability in a collapsing warm gas cloud and study the effect of metallicity and radiation flux. For the calculation, we adopt a one-zone approximation and the cloud contraction is simplified as free-fall collapse. When the cloud density reaches critical value $`n_f`$, the cloud fragments into cool dense condensations via thermal instability. From our calculation, the critical density $`n_f`$ increases as metallicity decreases, and also as radiation increases. Condensations collide with each other and self-gravitating clumps will be produced when the mean cloud density becomes sufficiently high; then stars will form. Expansion of the H II region around the massive star and supernova explosions will blow off surrounding gas and stop star formation process. When the mean density at the time of star formation is high, high virial velocity prevents expansion of the H II region. Also, in such high-density environments, the star formation timescale is shorter than the lifetime of a massive star. Then the gas in cluster-forming region will be converted into stars efficiently, before the gas is dispersed by expanding H II region or supernova explosions. In our calculation, such high density is realized in the contracting low-metallicity gas. And if the formation of a contracting gas cloud is possible, a strong-radiation environment is another candidate. Thus, it is suggested that high star formation efficiency and bound cluster formation are likely achieved in low-metallicity and/or strong-radiation environments. Such environments exist in dwarf galaxies, the early stage of our Galaxy and starburst galaxies. According to the result, globular clusters in these galaxies may be produced through the mechanism we show in this paper.
We woud like to thank H. Sato for valuable comments and H. Kamaya for useful discssion. This work is supported in part by the Japanese Grant-in-Aid for Scientific Research on Priority Areas (No. 10147105) of the Ministry of Education, Science, Sports and Culture of Japan.
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# CuSiO3 : a quasi - one - dimensional 𝑆=1/2 antiferromagnetic chain system
## I Introduction
The one dimensional spin system CuGeO<sub>3</sub> has attracted considerable attention in the past years, since it was the first (and up to now the only one) inorganic compound which exhibits a spin - Peierls transition (SP). The simplicity of crystal growth and the large variety of possible substitutions on the Cu and Ge sites promoted a huge amount of experimental and theoretical studies. Completely new phenomena like the coexistence of the SP state with a long - range antiferromagnetically ordered state in slightly doped CuGeO<sub>3</sub> (Si on the Ge site or Zn on the Cu site) and the the strong influence of the frustration due to the next - nearest - neighbor exchange on the magnetic and thermodynamic properties were reported. Partial substitution of Ge by Si in CuGe<sub>1-x</sub>Si<sub>x</sub>O<sub>3</sub> has been an important subject in this field, but despite considerable efforts, it was not possible to substitute more than 50 % Ge by Si without changing the structure (x $``$ 0.1 for single crystals, x $``$ 0.5 for polycrystals ). Pure CuSiO<sub>3</sub> was considered to be non existent. However, recently Otto et al. succeeded in the synthesis of reasonable amounts of pure isostructural CuSiO<sub>3</sub> by using the silicate mineral dioptase Cu<sub>6</sub>Si<sub>6</sub>O$`{}_{18}{}^{}`$6H<sub>2</sub>O as a starting material. Compared to the Ge homologue the unit cell volume of CuSiO<sub>3</sub> is reduced ($``$ \- 3.8 % ) due to the smaller size of Si. This naturally leads a modification of bond angles and lengths, which should have a strong influence on the strength of the magnetic interaction governed by the super - exchange between neighboring Cu<sup>2+</sup> ions via the O<sup>2-</sup> ions. The obvious question is how these structural changes do affect the ground state properties of this system. Here we present the first investigation of the physical properties of CuSiO<sub>3</sub>, based on susceptibility, specific heat and Cu- nuclear quadrupole resonance experiments.
## II Experimental details
High quality crystals of dioptase from the type locality Altyn Tyube (Kazakhstan) were heated up to 873 K and held at this temperature for six hours in order to obtain dehydrated dioptase (black dioptase). The black dioptase was decomposed by a subsequent heat treatment at higher temperatures (1050 K) and for 20 hours under nitrogen atmosphere. The finely ground samples of darkish brown color were characterized by X - ray powder diffraction using the Guinier method. The diffraction pattern reveals a mixture of three different phases. A quantitative phase analysis reveals that about 76 wt.-% of the mixture consist of the new phase CuSiO<sub>3</sub>. The other phases are CuO (tenorite, 13.7 wt.-%) and SiO<sub>2</sub> (amorphous, 10.3 wt.-%). The orthorhombic unit - cell of the new phase CuSiO<sub>3</sub> was refined from Guinier data giving a = 4.6357(6) Å, b = 8.7735(11) Å, and c = 2.8334(4) Å. The lattice constants of CuSiO<sub>3</sub> are in good agreement with an extrapolation of the results from the diluted system CuGe<sub>1-x</sub>Si<sub>x</sub>O<sub>3</sub> to x = 1. The crystal structure was determined to be isostructural to CuGeO<sub>3</sub>. Details of the synthesis and structure characterization are given in Ref. . Magnetization, specific heat and NQR measurements were performed on samples taken from the same batch. The dc - magnetization measurements at low fields $`\mu _0`$ H $``$ 1 T were carried out using a commercial SQUID magnetometer (MPMS, Quantum Design). The ac - and dc - magnetization measurements in higher fields (1 T $`\mu _0`$ H $``$ 14 T) and the specific heat measurements were performed in a commercial multi - purpose device (PPMS, Quantum Design). Ac- and dc- magnetization are extracted from the induction signal of a mutual inductance coil arrangement. The specific heat was determined by standard relaxation technique and an advanced two - tau model was applied to analyze the thermal response. The data obtained were corrected by subtracting the specific heat contributions of the sample holder, thermometers, heaters and glue. The NQR measurements were performed on a conventional pulsed spectrometer using the point - by point method. The masses of the samples are m = 91.40 mg, 8.72 mg and 180.20 mg for susceptibility, specific heat and NQR - measurements, respectively.
## III Results
The temperature dependence of the magnetic susceptibility $`\chi `$(T) of the CuSiO<sub>3</sub> powder sample in a magnetic field of 1 T is shown in Fig. 1. Below room temperature, $`\chi `$(T) increases with decreasing temperature, confirming localized Cu<sup>2+</sup> moments. Between 200 K and 300 K the susceptibility data follow nicely a Curie - Weiss law with a Weiss \- temperature of $`\theta `$ = -7.2 K, indicating rather weak antiferromagnetic coupling. An effective moment of $`\mu _{eff}`$ = 1.56 $`\mu _B`$ per $`Cu^{+2}`$ ion is determined which corresponds to a g - factor of 1.80. Both impurity phases CuO and $`SiO_2`$ exhibit a small susceptibility which could be neglected in a first analysis. Scaling the measured susceptibility with the estimated amount of pure CuSiO<sub>3</sub> in the sample results in an $`\mu _{eff}`$ \- value of $`1.79\mu _B`$ and a $`g`$-factor of 2.06. These values are very close to those found for CuGeO<sub>3</sub> and expected theoretically for free S = 1/2 spins with $`g=2`$ and $`\mu =g\mu _B\sqrt{S(S+1)}=1.73\mu _B`$.
At lower temperatures the susceptibility exhibits a broad maximum at $`T_{m,\chi }`$ = 13.5 K, which is a hallmark for low - dimensional spin systems. Above 8 K, $`\chi `$(T) could be fitted very well with the numerical results of Klümper for S=1/2 Heisenberg chains. The best fit is obtained with a nearest - neighbor coupling of J/k<sub>B</sub> = 2$`T_{m,\chi }`$/1.282 = 21 K and a slightly enhanced $`\mu _{eff}`$ \- value of 1.60 $`\mu _B`$ (note that we used the unscaled susceptibility for the fit and that we define $`J`$ by the exchange Hamiltonian: $`H=J𝐬_𝐢𝐬_{𝐢+\mathrm{𝟏}}+`$ $`\alpha 𝐬_𝐢𝐬_{𝐢+\mathrm{𝟐}}`$ ). The quality of the fit suggests that frustration effects, i.e. an antiferromagnetic interaction $`J^{}`$ between next - nearest - neighbors, are negligible in CuSiO<sub>3</sub> $`(\alpha =J^{}/J=0)`$. This is in contrast to CuGeO<sub>3</sub> where $`\alpha `$ values between 0.24 and 0.35 are proposed. Furthermore it is evident that the Cu - O(2) - Cu exchange in CuSiO<sub>3</sub> is much weaker (J/k<sub>B</sub> = 21 K) than in CuGeO<sub>3</sub> (J/k<sub>B</sub> $``$ 160 K). Below 8 K, $`\chi `$(T) decreases very rapidly and saturates at the lowest temperatures. The derivation of the susceptibility d$`\chi `$/dT shows a peak at $`T_0`$ = 7.9 K which gives clear evidence for a cooperative phase transition at this temperature (see inset Fig.1). The signature of the transition in the susceptibility looks more like a long - range antiferromagnetic order in a polycrystaline sample than a spin - Peierls transition. For the latter scenario, one expects a vanishing susceptibility at lowest temperatures $`\chi `$(T $``$ 0) = 0, which is not observed. The absence of a Curie - like tail in the susceptibility at the lowest temperatures indicates the absence of defects and thus points to a high crystalline perfection of the CuSiO<sub>3</sub> phase.
The presence of a phase transition at $`T_0`$ is clearly confirmed by the specific heat results. The temperature dependence of the total specific heat of three CuSiO<sub>3</sub> \- crystals (total mass of 8.72 mg, 1 mol = 139.6 g) below 20 K is shown in Fig. 2. For a comparison, the data of Liu et al. for CuGeO<sub>3</sub> are plotted in the same figure. The specific heat of CuSiO<sub>3</sub> shows a very clear $`\lambda `$ \- shaped, asymmetric anomaly at $`T_0`$ = 7.9 K with a specific heat jump of $`\mathrm{}C(1.50\pm \mathrm{\hspace{0.17em}0.05})J/molK`$. This is comparable to the $`\mathrm{}C`$ \- value of approximately $`1.8J/molK`$ found for CuGeO<sub>3</sub>. The total specific heat obtained is the sum of the lattice term $`C_{ph}(T)`$ from the phonons and the magnetic term $`C_m(T)`$ from the spin system. The separation of the two contributions is not trivial. At low temperatures, the phonon contribution should follow a $`T^3`$ law, i.e., $`C_{ph}=\beta T^3`$. Because of the smaller mass of the Si ions compared to the Ge ions ($`m_{Si}/m_{Ge}\mathrm{\hspace{0.17em}0.39}`$) one can expect a harder phonon spectra in CuSiO<sub>3</sub> and thus a reduced phonon contribution ( $`\beta _{CuSiO_3}<\beta _{CuGeO_3}\mathrm{\hspace{0.17em}0.32}mJ/molK^4`$). Thus, the estimate of the phonon contribution for CuGeO<sub>3</sub> (see Fig. 2, dashed line) yields therefore an upper limit for the phonon contribution in CuSiO<sub>3</sub>. This clearly demonstrates, that below 20 K, the specific heat of CuSiO<sub>3</sub> is dominated by the magnetic contribution. Well below $`T_0`$, C(T) follows a $`T^3`$ power law with a coefficient $`\beta _m`$ ($`4.5mJ/molK^4`$), more than one order of magnitude larger than that expected for the phonon contribution , indicating that it has to be related to the magnetic interactions. Such a power law is expected for long - range ordered 3D antiferromagnets with weak or absent anisotropy. The $`T^3`$ power law at low temperatures is in contrast to the experimental findings on the spin - Peierls compound CuGeO<sub>3</sub> where the opening of a gap in the magnetic excitation spectra leads to an exponential decrease of the specific heat. Therefore our specific heat measurements, too, support long - range antiferromagnetic order rather than a spin - Peierls dimerization at $`T_0`$.
At higher temperatures $`T>T_0`$, the specific heat $`C_m(T)`$ of a Heisenberg chain without frustration exhibits a maximum at $`T_{m,c}=0.75`$ $`T_{m,\chi }`$ $`=10.1K`$ with a value of $`C_m(T_{m,c})`$ $`0.35R2.9J/molK`$, which is independent of J. Scaling this value with the amount of pure CuSiO<sub>3</sub> in the sample results in an expected magnetic contribution of 2.2 J/molK. The experimental value of the specific heat at 10 K is only slightly lower. Further measurements and analysis are currently under progress to improve the estimation of the magnetic specific heat contribution. At even higher temperatures $`(T>J/k_B)`$ $`C_m(T)`$ becomes very small and C(T) originates mainly from the phonon contribution. Therefore the germanate exhibits a larger total specific heat than the silicate as evidenced from Fig.2.
The influence of magnetic fields up to 14 T on the magnetic susceptibility and the specific heat is shown in Fig. 3. In the susceptibility $`\chi (T,H)`$ the signature of the transition at $`T_0`$ = 7.9 K is smeared out, the drop in the susceptibility is reduced and the susceptibility maximum at $`T_{m,\chi }`$ is shifted significantly to lower temperatures with field. This shift of $`T_{m,\chi }`$, which indicates the suppression of the antiferromagnetic in - chain correlations with increasing field, is in good agreement with the theoretical calculations of Klümper (see inset in Fig. 3) for the $`S=1/2`$ Heisenberg chain. In the specific heat C(T) the transition at $`T_0`$ = 7.9 K is clearly observable in fields up to 14 Tesla and the $`T^3`$ power law at low temperatures is preserved. The antiferromagnetic order temperature $`T_0`$ shows only a weak field dependence as indicated in the magnetic phase diagram plotted in Fig. 4a. No other transitions are visible in our $`C(T,H)`$ measurements. The temperature dependent $`\chi (T)`$ measurements at fixed fields (Fig. 3) and additionally field dependent ac - susceptibilty measurements at fixed temperatures (Fig. 4b) evidences a broadened (due to the random orientation of the powder particles) transition at $`\mu _0H_{SF}3T`$, which looks very similar to a spin \- flop transition. The phase diagram corresponds to that observed for the AF - phase in doped CuGeO<sub>3</sub>, but it is quite different from that expected and confirmed for a spin- Peierls transition. The almost field independent transition temperature and the presence of a spin - flop - like transition are strong evidences for a antiferromagnetically ordered ground state in CuSiO<sub>3</sub>.
We have investigated the nuclear quadrupole resonance (NQR) of Cu in CuSiO<sub>3</sub> to deduce microscopic informations at nuclear sites. A spectrum at 4.2 K is shown in Fig. 5. The lines are fitted well by a Gaussian function and the central frequencies $`{}_{}{}^{63,65}\nu _{NQR}^{}`$ and the line widths $`{}_{}{}^{63,65}\mathrm{}\nu _{NQR}`$ were determined at different temperatures (see inset of Fig.5). The NQR signals at 4.2 K have been found at 26.88 $`\pm `$ 0.02 MHz for <sup>63</sup>Cu and 24.88 $`\pm `$ 0.02 MHz for <sup>65</sup>Cu. The ratio of the frequencies of the NQR signals are in good agreement with that expected from the nuclear quadrupole moments of Cu ($`{}_{}{}^{63}Q`$ / $`{}_{}{}^{65}Q`$ = 1.081 ). Also the ratio of the signal intensities $`{}_{}{}^{63}I/^{65}I=2.8`$ corresponds to that of the natural abundance of the isotopes <sup>63</sup>Cu and <sup>65</sup>Cu with $`{}_{}{}^{63}I/^{65}I=2.20`$. Furthermore the presence of the impurity phase CuO is confirmed nicely by NQR measurements. At room temperature the <sup>63</sup>Cu NQR frequency of CuO is found at $`{}_{}{}^{63}\nu _{NQR}^{}`$ = 20.6 MHz, which is in good agreement with the literature. Due to the AF order in CuO at $`T_N`$ = 230 K the $`{}_{}{}^{63}\nu _{NQR}^{}`$ is shifted to much higher frequencies below $`T_N`$ ($`{}_{}{}^{63}\nu _{NQR}^{}137MHz`$ for the central line at 4.2 K). Therefore the observed resonance signals presented in Fig. 5 could be clearly assigned to the Cu NQR lines of the CuSiO<sub>3</sub> phase. No other resonance lines are observed in the frequency range of 20 - 90 MHz indicating the crystallographically equivalence of the Cu site. Compared to CuGeO<sub>3</sub> with 34.23 $`\pm `$0.02 MHz for $`{}_{}{}^{63}Cu`$ and 31.66 $`\pm `$0.02 MHz for $`{}_{}{}^{65}Cu`$ at 4.2 K the NQR lines in CuSiO<sub>3</sub> are shifted to lower frequencies. One possible explanation is the effect of the modified bond lengths and angles on the electric field gradient EFG which affects strongly $`{}_{}{}^{63,65}\nu _{NQR}^{}`$. The EFG at the Cu site originates mainly from the ionic charge distribution of the surrounding ions (so called lattice contribution) and the 3d Cu charge distribution (valence) itself. Compared to CuGeO<sub>3</sub> the Cu-O(2) bond length is nearly the same whereas the Cu-Cu distance is reduced ( $`3.5\%`$). In a simple point charge model for the lattice contribution this could reduce the EFG and therefore lower the NQR frequencies. Details of the NQR results and a complex analysis of the data will be presented elsewhere. Surprisingly $`{}_{}{}^{63,65}\nu _{NQR}^{}`$ and the line widths $`{}_{}{}^{63,65}\mathrm{}\nu _{NQR}`$ exhibit only a weak temperature dependence between 4.2 and 40 K and especially around the transition at $`T_0`$ = 7.9 K, no anomaly is observed (see inset Fig. 5). Usually an antiferromagnetic phase transition is associated with the appearance of strong internal magnetic fields on the Cu site. This should result in a remarkable transformation of the pure NQR spectrum at T $`>T_N`$ to a high frequency AFMR spectrum for T $`<T_N`$ which is perturbed by quadrupole interaction. The origin of this absence of any signature of the transition is presently not clear.
## IV Discussion
CuSiO<sub>3</sub>, isotypic to the spin - Peierls compound CuGeO<sub>3</sub> was synthesized from the mineral dioptase and the physical properties of this new compound were determined by means of susceptibility, specific heat and Cu nuclear quadrupole resonance measurements. The susceptibility of CuSiO<sub>3</sub> is in good agreement with the theoretical results for a quasi one dimensional $`S=1/2`$ Heisenberg antiferromagnet with a nearest - neighbor exchange constant of $`J/k_B=\mathrm{\hspace{0.17em}21}K`$ and without a significant next - nearest neighbor exchange interaction ($`J^{}=\mathrm{\hspace{0.17em}0}`$). This is in contrast to the findings for CuGeO<sub>3</sub> where frustration effects due to next - nearest neighbor interactions play a crucial role ($`\alpha =J^{}/J\mathrm{\hspace{0.17em}0.35}`$) and where a much higher exchange coupling constant of $`J/k_B\mathrm{\hspace{0.17em}160}K`$ was found. According to the so - called Goodenough - Kanamori - Anderson (GKA) rules a change from an antiferromagnetic exchange to a ferromagnetic exchange is expected when the bond angle is near $`90^{}`$. Therefore the much smaller J value in CuSiO<sub>3</sub> could easily be attributed to the reduction of the Cu - O(2) - Cu bond angle from $`99^{}`$ in CuGeO<sub>3</sub> to $`94^{}`$ in CuSiO<sub>3</sub>. In contrast, the disappearance of the frustration is surprising, since the corresponding bond angles are fare away from $`90^{}`$. One possible explanation is, that the smaller Cu - O(2) - Cu bond angle of $`94^{}`$ leads to a cancellation of the antiferromagnetic and ferromagnetic exchange contributions and a vanishing nearest - neighbor exchange $`J`$. The remaining next - nearest - neighbor exchange $`J^{}`$ would then transform one atomic chain into two independent spin chains. Since the susceptibility (and the specific heat) is the same as that of a single spin chain, the present experimental results do not allow to distinguish between both cases. However, theoretical calculations do not support this scenario.
The large residual susceptibility below $`T_N`$, the $`T^3`$ power law in the magnetic specific heat at low temperatures, the spin - flop like transition in a magnetic field as well as the H-T magnetic phase diagram are very strong evidences for the antiferromagnetic nature of the transition. The comparatively large ratio between the ordering temperature $`T_N`$ and the temperature of the maximum in the susceptibility $`T_{m,\chi }`$ indicate that the ratio between inter - and intra - chain exchange ( $`J_{}/J`$) is significantly larger in CuSiO<sub>3</sub> than in CuGeO<sub>3</sub>. This is a natural consequence of the weakness of the intra - chain exchange. An estimate of the inter - chain coupling constant for quasi - one - dimensional chains is given by the following expression: $`J_{}=\frac{T_N}{1.28\sqrt{ln(5.8J/T_N)}}=3.7K`$. From this we obtain a ratio of $`J_{}/J0.18`$. LDA band structure calculations also reveal a much weaker intra - chain coupling J and a larger $`J_{}/J`$ ratio ($`0.14`$) compared to CuGeO<sub>3</sub>, which is in agreement with our experimental results. However, they also suggests that the frustration should still be significant in CuSiO<sub>3</sub>. The weaker intra - chain coupling J and the much larger $`J_{}/J`$ ratio are obvious responsible for the occurrence of an AF - transition instead of a SP - transition.
In conclusion it is shown, that CuSiO<sub>3</sub> is a quasi - one - dimensional S=1/2 Heisenberg chain system which undergoes a transition to long range antiferromagnetic order at $`T_N=7.9K`$. There is no direct evidence for spin - Peierls transition in this new compound. Field dependent susceptibility measurements reveal a spin - flop phase at $`\mu _0H_{SF}3T`$. Among all inorganic low dimensional Cu based spin - systems the edge - sharing compound CuSiO<sub>3</sub> exhibits the smallest intra - chain exchange AF - coupling constant of $`J/k_B=21K`$. The existence of the two homologeus compounds CuGeO<sub>3</sub> and CuSiO<sub>3</sub> with different ground state properties provides an excellent basis for the application of theoretical models and methods of low dimensional physics.
## Acknowledgments
We acknowledge fruitful discussions with H. Rosner and S.-L. Drechsler.
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# Infrared observations of hot gas and cold ice toward the low mass protostar Elias~29Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) and with the participation of ISAS and NASA.
## 1 Introduction
The general picture of low mass star formation has been formed since the 1980’s with the availability of infrared and millimeter wavelength broad band photometry from the ground, and with the IRAS satellite and KAO observatory (e.g., Lada & Wilking lad84 (1984); Adams et al. ada87 (1987); Hillenbrand et al. hil92 (1992); André et al. and93 (1993)). A classification scheme was made, where the continuum emission of Class 0 and I objects peaks in the submillimeter and far-infrared. These objects are still deeply embedded in their accreting envelopes. In the Class II phase, the wind of the protostar has cleared its surrounding environment, such that it becomes optically visible, and shows H i emission lines. The continuum emission of these objects peaks in the near-infrared, but there is still significant excess emission above the stellar continuum. They are believed to be surrounded by optically thick dusty disks. Finally, little dust emission remains for Class III objects, when the disk is optically thin, and planetary companions may have been formed.
Our knowledge of the physical and chemical state and evolution of the material surrounding protostars, has progressed with the availability of medium and high resolution spectroscopic instrumentation at near and mid-infrared wavelengths ($`220`$ $`\mu \mathrm{m}`$). The progress made, is best illustrated by the observations of high mass protostars, which are bright and easy to observe. The (ro-)vibrational bands of various molecules (CO, $`\mathrm{H}_2\mathrm{O}`$, $`\mathrm{CH}_3\mathrm{OH}`$, silicates) were observed from the ground, revealing profile variations of the 3.07 $`\mu \mathrm{m}`$, and 4.67 $`\mu \mathrm{m}`$ $`\mathrm{H}_2\mathrm{O}`$ and CO ice bands (e.g. Smith et al. smi89 (1989); Tielens et al. tie91 (1991); Chiar et al. chi98 (1998)). This was interpreted as evaporation of the volatile CO ice, and crystallization of $`\mathrm{H}_2\mathrm{O}`$ ice in the molecular envelopes. This heating effect is strengthened by the detection of hot gas in 4.6 $`\mu \mathrm{m}`$ CO observations (Mitchell et al. mit91 (1991)). With the launch of the Infrared Space Observatory in 1995 (ISO; Kessler et al. kes96 (1996)), it became possible to observe all other molecular bands in the infrared (Whittet et al. whi96 (1996)). It was shown that high mass protostellar evolution can be traced in the gas-to-solid abundance ratios (van Dishoeck et al. dis96 (1996); van Dishoeck & Blake disb98 (1998)), and the profiles of the ice bands, in particular solid $`\mathrm{CO}_2`$ (Gerakines et al. ger99 (1999); Boogert et al. boo00 (2000)). Thus, there is overwhelming evidence that thermal processing, i.e. evaporation and crystallization of ices in and around hot molecular cores, plays an important role in the evolution of high mass molecular envelopes.
The composition and evolution of the molecular material around low mass protostars are not as well studied. It seems unlikely that the molecular material evolves similar to that around high mass protostars. Low mass protostars evolve much slower, release less radiative energy, drive less energetic winds, and form disks. It is not established whether low mass objects possess hot cores as well, and whether the ices survive the process of star formation. If (some of) the ices survive, are they included into comets, and if so, are the ice structure and composition still the same compared to interstellar ices? How important are energetic processes, such as cosmic ray bombardment, in altering the ice composition on the long time scale of the formation of low mass stars?
To investigate the influence of low mass protostars on their molecular envelope, we make an infrared spectroscopic study Elias~29, also called WL 15 and YLW 7 (Elias eli78 (1978), Wilking et al. wil83 (1983), Young et al. you86 (1986)). On a large scale, Elias~29 lies in core E, which is in the south-east corner of the 1$`\times `$2 pc extended compact CO ridge L 1688 (Loren et al. lor90 (1990)) in the densest part of the $`\rho `$ Ophiuchi cloud, at a distance of $``$160 pc from the earth (Wilking & Lada wil83 (1983); Whittet whi74 (1974)). It is the reddest object found in the near-infrared survey of this cloud by Elias (eli78 (1978)), without a counterpart at optical wavelengths. For our observations, we used Elias’ coordinates (J2000):
$`\alpha =16^\mathrm{h}27^\mathrm{m}09^\mathrm{s}.3`$ $`\delta =24^\mathrm{o}37^{}21^{\prime \prime }`$.
The overall spectrum of Elias~29 is typical for a heavily embedded Class I source, probably in a late accretion phase (Wilking et al. wil89 (1989); André & Montmerle and94 (1994); Greene & Lada gre96 (1996); Saraceno et al. sar96 (1996)). The embedded nature is also revealed by its high extinction, and by the cold compact envelope observed at millimeter wavelengths (André & Montmerle and94 (1994); Motte et al. mot98 (1998)). Elias~29 is associated with a molecular outflow (Bontemps et al. bon96 (1996); Sekimoto et al. sek97 (1997)). With a bolometric luminosity of $`36L_{}`$ (Chen et al. che95 (1995)), Elias~29 is the most luminous protostar in the $`\rho `$ Oph cloud, which makes this source very suitable for spectroscopic studies. The relatively high luminosity, and high bolometric temperature ($`T_{\mathrm{bol}}410`$ K) imply an age in the range 0.5–4 10<sup>5</sup> yr (Chen et al. che95 (1995)). In the pre-main-sequence evolutionary tracks of Palla & Stahler (pal93 (1993)), this corresponds to a star with end mass 3.0–3.5 $`M_{}`$. Elias~29 might thus be a precursor Herbig AeBe star. This classification is however uncertain. For example, it has been argued from the SED, and the absence of mid-infrared emission features, that Elias~29 is a 1 $`M_{}`$ protostar with a large accretion luminosity, and a spectral type of K3-4 at the birth line (Greene & Lada gre00 (2000)).
This Paper is structured as follows. Technical details on the ISO infrared observations are given in Sect. 2. All the observed emission and absorption features are discussed in detail in Sect. 3. Section 3.1 gives a description of the continuum shape, and a comparison to other lines of sight. The ice composition and thermal history, and the silicate band depth with inferred extinction and column densities toward Elias~29 are discussed in Sect. 3.2. Then, numerous lines of gaseous CO and $`\mathrm{H}_2\mathrm{O}`$ are detected, and modeled to derive gas temperatures and column densities (Sect. 3.3). The molecular abundances and gas-to-solid ratios of Elias~29 are compared to a sample of sight-lines, ranging from dark cloud cores to evolved protostars. A comparison with high mass protostars is made (Sect. 4.1). Section 4.2 discusses the origin of the wealth of observed emission and absorption features and puts them in a geometrical picture, where we review the evidence for an extended envelope and an accretion disk. We conclude in Sect. 5 with a summary and suggestions for future observations.
## 2 Observations
### 2.1 The 2.3–45 $`\mu `$m spectrum
A low resolution ($`R=\lambda /\mathrm{\Delta }\lambda =400`$), full 2.3–45 $`\mu \mathrm{m}`$ spectrum of Elias~29 was obtained with the ISO Short Wavelength Spectrometer (ISO–SWS; de Graauw et al. gra96 (1996)) during revolution 267 (August 10 1996). The ISO–SWS pipeline and calibration files, available in July 1998 at SRON Groningen, were applied. The spectrum is generally of good quality, with well-matching up and down scans, and no serious dark current problems, except for band 2C (7–12 $`\mu \mathrm{m}`$). Here, we found that the up and down scans deviate over the silicate band. One scan showed good agreement with a ground-based spectrum of Hanner et al. (han95 (1995)), and we used this to correct the deviating scan. Standard after-pipeline steps were applied, such as low order flat-fielding, sigma clipping and re-binning (see also Boogert et al. boo98 (1998)). The twelve sub-spectra in the 2-45 $`\mu \mathrm{m}`$ range match fairly well at the overlap regions. Small correction factors ($`<`$15%) were applied to correct for the band jumps.
At selected wavelength ranges (3–3.6, 4–9, and 19.5–28 $`\mu \mathrm{m}`$), we also obtained high resolution ($`R=1500`$) ISO–SWS grating spectra, in revolution 292 (September 04 1996). These were reduced similarly to the low resolution spectrum. We found that the overall shape of the spectrum near 4–5 $`\mu \mathrm{m}`$ is quite badly affected by detector memory effects, presumably due to the occurrence of scan breaks (de Graauw et al. gra96 (1996)). We corrected for this, by applying a wavelength-dependent shift to match the low resolution spectrum. This does not affect our conclusions, since the high resolution spectrum was only used to study narrow features. Also, near 6.9 $`\mu \mathrm{m}`$ the scans deviate significantly because of memory effects. This problem is reflected in the large error bars given in this paper, as they were derived from the difference between the average up and down scans.
### 2.2 The 45–190 $`\mu `$m spectrum
Elias~29 was observed during Revolution 484 (March 14 1997) with the ISO Long Wavelength Spectrometer (ISO–LWS; Clegg et al. cle96 (1996)). We obtained 15 scans covering the range from 43 $`\mu \mathrm{m}`$ to 197 $`\mu \mathrm{m}`$ in the low resolution mode ($`R`$200) for a total of 2611 sec of integration time. The data were reduced using the Off-Line-Processing package (OLP) version 7 and the ISO-Spectral-Analysis-Package (ISAP) version 1.3. The spectra were flux calibrated using Uranus (Swinyard et al. swi96 (1996)). We find that at the ISO–SWS/LWS overlap region near 45 $`\mu \mathrm{m}`$, LWS has 35% higher flux than SWS. This difference is only slightly larger than the absolute calibration uncertainties of the two instruments, and thus it is doubtful that this can be ascribed to the presence of extended emission in the larger aperture of ISO–LWS ($`80^{\prime \prime }`$ versus $`25^{\prime \prime }`$). We therefore decided to multiply the LWS spectrum down with this factor.
## 3 Results
### 3.1 The spectral energy distribution (SED)
Elias~29 is only visible at wavelengths larger than $``$1.5 $`\mu \mathrm{m}`$ (Greene & Lada gre96 (1996); Elias eli78 (1978)). Our ISO observations show that the continuum emission rises steeply between 2–3 $`\mu \mathrm{m}`$, reaches a maximum of $`\lambda F_\lambda `$=15$`\times 10^{16}`$ W $`\mathrm{cm}^2`$ at $`\lambda `$$`\mu \mathrm{m}`$, and is remarkably flat with $`\lambda F_\lambda `$$`8\times 10^{16}`$ W $`\mathrm{cm}^2`$ between 20 and 100 $`\mu \mathrm{m}`$ (Fig. 1). The emission has dropped to $`\lambda F_\lambda `$$`4\times 10^{16}`$ W $`\mathrm{cm}^2`$ at 200 $`\mu \mathrm{m}`$, and by four orders of magnitude at 1300 $`\mu \mathrm{m}`$. Our near-infrared spectral continuum fluxes are in excellent agreement with broad band fluxes from ground-based observations (Elias eli78 (1978)). Also the ground-based small beam 10 and 20 $`\mu \mathrm{m}`$ observations, as well as the large beam 12 and 25 $`\mu \mathrm{m}`$ IRAS fluxes, match the ISO–SWS observation well, thus indicating that at these wavelengths the emission is well confined within a region of 8<sup>′′</sup> in diameter (Fig. 1; Lada & Wilking lad84 (1984); Young et al. you86 (1986)). The reasonable match of the ISO–SWS and LWS spectra (Sect. 2.2) indicates that also at 45 $`\mu \mathrm{m}`$ the emission is not very extended ($`<25^{\prime \prime }`$). At 100 $`\mu \mathrm{m}`$, however, some large scale emission may be present, since the IRAS flux, observed in a 5.5 times larger aperture, is a factor of 2 larger compared to the ISO measurement.
The observed SED of Elias~29 is different from that of massive protostars such as GL 2591, and GL 7009S, which peak in the far-infrared (Fig. 2). It has been proposed that the shape of SEDs is independent of the luminosity of the central object, and is rather determined by the total dust column density along the line of sight (Ivezic & Elitzur ive97 (1997)). In the “standard model” of Ivezic & Elitzur, GL 2591 and GL 7009S would have a column density corresponding to an $`A_\mathrm{V}`$ of several hundred. Elias 29 must have an $`A_\mathrm{V}<100`$, because it does not peak in the far infrared. However, the flatness of the SED up to 100 $`\mu \mathrm{m}`$ is not reproduced in these models. A lower column density alone thus cannot explain the differences between the SED of Elias~29 and massive protostars. Other factors, such as a different density gradient, and the presence of a circumstellar disk are probably important. We will discuss the structure of Elias~29, in relation to the detected gas and ice absorption features, in Sect. 4.2.
### 3.2 Ice and dust absorption bands
Numerous absorption bands of ices and silicates are present in the infrared spectrum of Elias~29 (Fig. 3), such that hardly any ’clean’ continuum emission is left. We identify each band, derive column densities, and, when possible, determine the ice mantle composition and thermal history. The full spectrum also allows a determination of upper limits of abundances for undetected, though astrophysically relevant molecules. The column densities for the species discussed below are summarized in Table 1.
#### 3.2.1 H<sub>2</sub>O ice
The infrared spectrum of Elias~29 shows all the vibration modes of $`\mathrm{H}_2\mathrm{O}`$ ice in absorption (Fig. 4). We see the O–H stretching mode at 3.0 $`\mu \mathrm{m}`$ (“$`\nu _1`$, $`\nu _3`$” in spectroscopic notation), the O–H bending mode at 6.0 $`\mu \mathrm{m}`$ (“$`\nu _2`$”), the libration or hindered rotation mode at $``$12 $`\mu \mathrm{m}`$ (“$`\nu _\mathrm{L}`$”), the combination mode at 4.5 $`\mu \mathrm{m}`$ (“$`3\nu _\mathrm{L}`$” or “$`\nu _2+\nu _\mathrm{L}`$”), and perhaps the lattice mode at $``$45 $`\mu \mathrm{m}`$.
The continuum determination is complicated by the large width of all these bands. In accordance with other studies (Smith et al. smi89 (1989); Schutte et al. sch96 (1996); Keane et al. kea00 (2000)), we used single blackbodies to fit the continuum locally, directly adjacent to the absorption bands. For the 6 $`\mu \mathrm{m}`$ band we took into account that laboratory spectra of the bending mode of $`\mathrm{H}_2\mathrm{O}`$ ice show a prominent wing on the long wavelength side, extended up to 8 $`\mu \mathrm{m}`$ (e.g. Hudgins et al. hud93 (1993), Maldoni et al. mal98 (1998)). We simultaneously fitted a blackbody continuum, normalized at 5.1 $`\mu \mathrm{m}`$, and a laboratory ice spectrum to the observed flux at 8 $`\mu \mathrm{m}`$ and the shape of the 6.0 $`\mu \mathrm{m}`$ feature (Fig. 5).
The shape of the 6.0 $`\mu \mathrm{m}`$ $`\mathrm{H}_2\mathrm{O}`$ bending mode is particularly sensitive to the ice temperature (e.g. Maldoni et al. mal98 (1998)). At higher $`T`$ in the laboratory, the strength of the main 6.0 $`\mu \mathrm{m}`$ component decreases at the expense of more absorption in the long wavelength wing. In the spectrum of Elias~29, the wing cannot be seen as a separate feature, since at 8 $`\mu \mathrm{m}`$ it blends with the very deep silicate band. However, the observed 6.0 $`\mu \mathrm{m}`$ band is relatively sharp, and it can only be fitted with $`\mathrm{H}_2\mathrm{O}`$ ice at $`T<80`$ K, with a best fit at $`T=40`$ K (Fig. 5). The excellent fit to the 6.0 $`\mu \mathrm{m}`$ band in Elias~29 indicates that the 5.83 and 6.2 $`\mu \mathrm{m}`$ excess absorptions detected toward several massive protostars (Schutte et al. sch96 (1996), sch98 (1998); Keane et al. kea00 (2000)), are not seen in this source (Sect. 3.2.7).
The observed peak position of the stretching mode of $`\mathrm{H}_2\mathrm{O}`$ ice toward Elias~29 is 3.07$`\pm `$0.01 $`\mu \mathrm{m}`$. The short wavelength wing is well matched with a laboratory ice at $`T=`$40 K, as for the bending mode (Fig. 5). The long wavelength wing, however, is poorly fitted. It has been realized since long that light scattering by large ice grains leads to extra extinction on the long wavelength wing (e.g. Léger et al. leg83 (1983)). To illustrate this, we calculate the extinction cross section for spherical silicate grains coated with ice mantles, applying the code given in Bohren & Huffman (boh83 (1983)) and the optical constants of Draine & Lee (dra84 (1984)) and Hudgins et al. (hud93 (1993)). Indeed, grains with a core+mantle radius of $``$0.6 $`\mu \mathrm{m}`$ provide a much better fit to the long wavelength wing than small grains do (Fig. 5). This effect is unimportant for the 6.0 $`\mu \mathrm{m}`$ band since it is intrinsically weaker, and the grains are smaller compared to the wavelength. In a more realistic approach, a distribution of grain sizes, as well as constraints to other observables such as continuum extinction, the total grain and ice column densities, and polarization need to be taken into account. Although there is a general consensus that large grains need to be invoked (e.g. Leǵer et al. leg83 (1983); Pendleton et al. pen90 (1990); Smith et al. smi89 (1989); Martin & Whittet mar90 (1990)), there is no unified grain model yet that obeys all the observational constraints (e.g. Smith et al. smi93 (1993); Tielens tie82 (1982)). Alternative absorbers at the long wavelength wing have been proposed, such as $`\mathrm{H}_2\mathrm{O}`$.$`\mathrm{NH}_3`$ bondings. As illustrated in Fig. 5, this effect is however small at the low column density ratio of $`\mathrm{NH}_3`$/$`\mathrm{H}_2\mathrm{O}`$ $`<0.13`$ toward Elias~29 (Sect. 3.2.7). A small contribution is also made by absorption by hydrocarbons (Sect. 3.2.4).
The peak optical depth of the 3.0 $`\mu \mathrm{m}`$ band is 1.85$`\pm `$0.08, which is in excellent agreement with the study of Tanaka et al. (tan90 (1990)). Using an integrated band strength $`A=2.0\times 10^{16}`$ cm molecule<sup>-1</sup>, we derive a column density of $`N`$($`\mathrm{H}_2\mathrm{O}`$)= $`3.0\times 10^{18}`$ $`\mathrm{cm}^2`$ for the small grain model, and $`3.7\times 10^{18}`$ $`\mathrm{cm}^2`$ for the large grain model. Since at present we can not favor one of these two cases, we will assume an average value of $`N`$($`\mathrm{H}_2\mathrm{O}`$)= (3.4$`\pm 0.6`$)$`\times `$10<sup>18</sup> $`\mathrm{cm}^2`$ in this paper. The error bar also includes the uncertainty in band strength, which increases with 10% when the ice is heated from 10 to 100 K (Gerakines et al. ger95 (1995)). Note that a column density determination from the 6.0 $`\mu \mathrm{m}`$ bending mode is more uncertain due to the unreliable continuum on the long wavelength side (Fig. 5). At this column density of $`\mathrm{H}_2\mathrm{O}`$ ice, the depth of the other vibrational modes is in good agreement with the observed spectrum of Elias~29 (Fig. 4).
#### 3.2.2 CO ice
The CO ice band at 4.67 $`\mu \mathrm{m}`$ in Elias~29 is contaminated by gas phase CO lines from low $`J`$ levels (Figs. 6 and 10). In particular, the P(1) line lies in the center of the ice band at 4.674 $`\mu \mathrm{m}`$. To study the band profile, we subtracted a model for the gaseous lines at $`T_{\mathrm{ex}}=750`$ K, $`N=5\times 10^{18}`$ $`\mathrm{cm}^2`$, and $`b_\mathrm{D}`$=5 $`\mathrm{km}\mathrm{s}^1`$ (Sect. 3.3). This increases the ice band width by $`0.9`$ $`\mathrm{cm}^1`$, to FWHM=4.40 $`\mathrm{cm}^1`$ (0.010 $`\mu \mathrm{m}`$). With a peak position of 4.673 $`\mu \mathrm{m}`$ (2140.1 $`\mathrm{cm}^1`$), the CO ice band observed toward Elias~29 is similar to that of the luminous protostar NGC~7538~:~IRS9 (Fig. 6; Tielens et al. tie91 (1991); Chiar et al. chi98 (1998)). The main, narrow component at 4.673 $`\mu \mathrm{m}`$ is attributed to pure solid CO, or CO embedded in an environment of apolar molecules. In particular, mixtures with O<sub>2</sub>, at an O<sub>2</sub>/CO ratio as much as 5 (Elsila et al. els97 (1997); Chiar et al. chi98 (1998)) provide good fits. Mixtures of CO with $`\mathrm{CO}_2`$ are generally too broad (Ehrenfreund et al. ehr97 (1997)). While the apolar, volatile component dominates the spectrum, both Elias~29 and NGC~7538~:~IRS9 show evidence for a wing on the long wavelength side. This is attributed to CO diluted in a mixture of polar molecules such as $`\mathrm{H}_2\mathrm{O}`$ and $`\mathrm{CH}_3\mathrm{OH}`$ (Chiar et al. chi98 (1998); Tielens et al. tie91 (1991)). Assuming a band strength $`A=1.1\times 10^{17}`$ cm molecule <sup>-1</sup> for both the polar and apolar components (Gerakines et al. ger95 (1995)), we derive $`N`$(CO ice)=1.7$`\times 10^{17}`$ $`\mathrm{cm}^2`$ with an apolar/polar ratio of $``$8, comparable to NGC~7538~:~IRS9. These results are in good agreement with the ground-based study of Kerr et al. (ker93 (1993)). Although NGC~7538~:~IRS9 seems to have a larger polar CO component in Fig. 6, this difference may merely reflect uncertainties in the continuum subtraction, and the fact that the NGC~7538~:~IRS9 spectrum is not corrected for gas phase CO lines.
#### 3.2.3 $`\mathrm{CO}_2`$ ice
The absorption bands of $`\mathrm{CO}_2`$ ice are prominently present in the infrared spectrum of Elias~29 (Fig. 3). We see the stretching and bending modes at 4.27 and 15.2 $`\mu \mathrm{m}`$ respectively. Not visible in this spectrum is the stretching mode of solid $`{}_{}{}^{13}\mathrm{CO}_{2}^{}`$ at 4.38 $`\mu \mathrm{m}`$, although the high resolution spectrum (Fig. 10) shows a hint of its presence. A very sensitive observation is presented elsewhere (Boogert et al. boo00 (2000)). The $`{}_{}{}^{12}\mathrm{CO}_{2}^{}`$ bending mode and the $`{}_{}{}^{13}\mathrm{CO}_{2}^{}`$ stretching mode have proven to be very sensitive to ice mantle composition and thermal history. In Elias~29, these bands do not show the narrow substructures seen in many other protostars, and attributed to heated polar $`\mathrm{CO}_2`$ ices (Boogert et al. boo00 (2000); Gerakines et al. ger99 (1999)). As for the CO ice band (Fig. 6), the width and peak position of the $`{}_{}{}^{13}\mathrm{CO}_{2}^{}`$ band very much resemble that of the luminous protostar NGC~7538~:~IRS9. Thus, the $`\mathrm{CO}_2`$ ice toward Elias~29 is mixed in with polar molecules, and is not much affected by heating. The $`{}_{}{}^{12}\mathrm{CO}_{2}^{}`$ column density is 22$`\pm `$4% relative to $`\mathrm{H}_2\mathrm{O}`$ ice, which is comparable to the values reported for high mass protostars (Gerakines et al. ger99 (1999)). Finally, we derive an isotope ratio of $`{}_{}{}^{12}\mathrm{CO}_{2}^{}`$/$`{}_{}{}^{13}\mathrm{CO}_{2}^{}`$=81$`\pm `$11 in the ice toward Elias~29, which is well within the range found for the local ISM (Boogert et al. boo00 (2000)).
#### 3.2.4 The 3.47 $`\mu \mathrm{m}`$ band
The long wavelength wing of the deep 3.0 $`\mu \mathrm{m}`$ absorption band shows a change of slope at 3.38 $`\mu \mathrm{m}`$, indicative of a shallow absorption feature (Fig. 7). This feature is also detected in an independent ground based study of Elias~29 (Brooke et al. bro99 (1999)). For consistency with ground based studies, the continuum on each side of the feature was assumed to start at 3.37 and 3.61 $`\mu \mathrm{m}`$. It must be emphasized however, that in particular on the long wavelength side, the continuum is poorly defined. Fitting a smooth 6-th order polynomial results in an absorption band centered on 3.49$`\pm `$0.03 $`\mu \mathrm{m}`$ with a peak optical depth of $`\tau `$=0.06 (Fig. 7). The width is FWHM=120$`\pm `$40 $`\mathrm{cm}^1`$, where the uncertainty includes the poorly constrained continuum on the long wavelength side. Features of similar width and peak position have been detected in several massive protostellar objects (Allamandola et al. all92 (1992)) and in low mass objects and quiescent molecular cloud material (Chiar et al. chi96 (1996)). A likely candidate for this 3.47 $`\mu \mathrm{m}`$ band is the C–H stretching mode of hydrocarbons. From the correlation of peak optical depths of this feature and the 3.0 $`\mu \mathrm{m}`$ ice band, it is concluded that the carrier for the 3.47 $`\mu \mathrm{m}`$ band resides in ices rather than in refractory dust (Brooke et al. bro96 (1996)). We find that with $`\tau `$(3.47 $`\mu \mathrm{m}`$)=0.06 and $`\tau `$(3.0 $`\mu \mathrm{m}`$)=1.85, Elias~29 follows this correlation very well.
#### 3.2.5 CH<sub>3</sub>OH ice
In several high mass protostars, the 3.47 $`\mu \mathrm{m}`$ band is blended with a distinct narrow feature centered on 3.54 $`\mu \mathrm{m}`$ (Allamandola et al. all92 (1992)). This feature is ascribed to the C–H stretching mode of solid $`\mathrm{CH}_3\mathrm{OH}`$. A direct comparison with the high mass protostar NGC~7538~:~IRS9 shows that, although the 3.47 $`\mu \mathrm{m}`$ bands have similar shapes, the 3.54 $`\mu \mathrm{m}`$ feature is absent in Elias~29 (Fig. 7). We determine a 3$`\sigma `$ upper limit to the peak optical depth of $`\tau `$(3.54 $`\mu \mathrm{m}`$)$`<`$0.036. Scaling with the observed depth and column density in NGC~7538~:~IRS9 (Brooke et al. bro99 (1999)), then results in an upper limit to the $`\mathrm{CH}_3\mathrm{OH}`$ ice column density $`N`$($`\mathrm{CH}_3\mathrm{OH}`$ ice)$`<1.5\times 10^{17}`$ $`\mathrm{cm}^2`$, or less than 5% of $`\mathrm{H}_2\mathrm{O}`$ ice toward Elias~29 (Table 3). The other modes of $`\mathrm{CH}_3\mathrm{OH}`$ ice are either much weaker, or are severely blended with the strong $`\mathrm{H}_2\mathrm{O}`$ and silicate bands (e.g. the C–O stretching mode at 9.7 $`\mu \mathrm{m}`$; Schutte et al. sch91 (1991); Skinner et al. ski92 (1992)) and thus do not provide better constraints on the $`\mathrm{CH}_3\mathrm{OH}`$ ice column density. Toward other low mass objects, and quiescent dark clouds, low upper limits have been set to the $`\mathrm{CH}_3\mathrm{OH}`$ ice abundance as well. The $`\mathrm{CH}_3\mathrm{OH}`$ ice abundance found in massive protostars is generally of the same magnitude (Chiar et al. chi96 (1996)), but in a few objects significantly larger (Dartois et al. dar99 (1999)), than these upper limits.
#### 3.2.6 The 6.85 $`\mu \mathrm{m}`$ band
Elias~29 is the first low mass protostar in which the 6.85 $`\mu \mathrm{m}`$ absorption band is detected (Fig. 5). After subtraction of the $`\mathrm{H}_2\mathrm{O}`$ ice band and the gas phase $`\mathrm{H}_2\mathrm{O}`$ lines (Fig. 8), we find that it has a peak optical depth of $`\tau `$0.07 and an integrated optical depth $`\tau _{\mathrm{int}}=7.8\pm 1.6`$ $`\mathrm{cm}^1`$. When scaled to the $`\mathrm{H}_2\mathrm{O}`$ ice column density, the strength of the 6.85 $`\mu \mathrm{m}`$ band toward Elias~29 is similar to high mass protostars (Keane et al. kea00 (2000)). The band profile, e.g. the sharp edge at 6.60 $`\mu \mathrm{m}`$, agrees very well with several high mass objects, in particular those tracing ’cold’ gas and dust (NGC~7538~:~IRS9, W 33A, GL 989). It clearly deviates from warmer lines of sight (e.g. S 140 : IRS1; Fig. 8). Thus, in this picture, we find that the material responsible for the 6.85 $`\mu \mathrm{m}`$ band toward Elias~29 is not significantly thermally processed. Given the low upper limits to the $`\mathrm{CH}_3\mathrm{OH}`$ ice column density toward Elias~29, only a fraction of the band, as for high mass objects, can be explained by the C–H bending mode of $`\mathrm{CH}_3\mathrm{OH}`$ ices (Schutte et al. sch96 (1996)). For a detailed band profile analysis and a discussion on the origin of the 6.85 $`\mu \mathrm{m}`$ band, we refer to Keane et al. (kea00 (2000)).
#### 3.2.7 Upper limits to solid $`\mathrm{CH}_4`$, $`\mathrm{NH}_3`$, $`\mathrm{H}_2\mathrm{CO}`$, HCOOH, OCS, and ‘XCN’
Several solid state species have been detected toward luminous protostars, but are absent toward Elias~29. The deformation mode of solid $`\mathrm{CH}_4`$ was detected toward protostars, with a peak position at 1303 $`\mathrm{cm}^1`$ (7.67 $`\mu \mathrm{m}`$), and a width FWHM=11 $`\mathrm{cm}^1`$ (Boogert et al. boo96 (1996); Dartois et al. 1998b ). For Elias~29 we can exclude this band to a peak optical depth of $`\tau <`$0.03, corresponding to $`N`$($`\mathrm{CH}_4`$)/$`N`$($`\mathrm{H}_2\mathrm{O}`$)$`<`$1.5%. This 3$`\sigma `$ upper limit is comparable to the detection in NGC~7538~:~IRS9 (Boogert et al. boo96 (1996)).
Solid $`\mathrm{NH}_3`$ was detected by its 9.10 $`\mu \mathrm{m}`$ inversion mode toward NGC~7538~:~IRS9 (Lacy et al. lac98 (1998)). Using the band strength determined in Kerkhof et al. (ker99 (1999)), the $`\mathrm{NH}_3`$ column density is 13% of $`\mathrm{H}_2\mathrm{O}`$ ice. Recent detections in other highly obscured lines of sight give similar (W 33A; Gibb et al. gib00 (2000)), or a factor 2 larger $`\mathrm{NH}_3`$ abundances (Galactic Center; Chiar et al. chi00 (2000)). To find this band in the deep silicate feature of Elias~29 we take the same approach as Lacy et al., by fitting a local straight line continuum to the wavelength regions 8.52–8.69 and 9.20–9.55 $`\mu \mathrm{m}`$. As a check, we perform the same procedure to the ISO–SWS spectrum of NGC~7538~:~IRS9 (Fig. 9). We confirm the detection of Lacy et al., although the peak optical depth $`\tau 0.16`$ is a factor 2 lower in our case. We ascribe this difference to the calibration uncertainties of ISO–SWS at this wavelength (Leech lee00 (2000)). For Elias~29, a feature with $`\tau 0.06`$ might be present. However, due to the poorly defined long wavelength side of the continuum (Fig. 9) and the ISO–SWS calibration uncertainties, we will assume a conservative upper limit to this band of $`\tau <`$0.1. This corresponds to a column density of $`N`$($`\mathrm{NH}_3`$)$`<3.5\times 10^{17}`$$`\mathrm{cm}^2`$, i.e. $`N`$($`\mathrm{NH}_3`$)/$`N`$($`\mathrm{H}_2\mathrm{O}`$)$`<`$13%. Other vibrational bands of $`\mathrm{NH}_3`$ do not provide better constraints. The equally strong N–H stretching mode at 2.90 $`\mu \mathrm{m}`$ (d’Hendecourt & Allamandola hen86 (1986)) is hidden in the steep wing of the 3.0 $`\mu \mathrm{m}`$ ice band, and there is no significant difference in this region between the laboratory spectra of pure $`\mathrm{H}_2\mathrm{O}`$ ice and the mixture $`\mathrm{H}_2\mathrm{O}`$:$`\mathrm{NH}_3`$=10:1 (Fig. 5). A similar problem exists for the N–H deformation mode at 6.16 $`\mu \mathrm{m}`$, which is hidden in the long wavelength wing of the $`\mathrm{H}_2\mathrm{O}`$ bending mode (Keane et al. kea00 (2000)). A feature with a peak optical depth of $`\tau <`$0.025 would be expected here (Sandford & Allamandola san93 (1993)). When subtracting water ice and vapor absorption, a weak band with an optical depth $`\tau =`$ 0.03 perhaps remains present at the expected wavelength (Fig. 8). Given the other positive and negative structure in the spectrum, we regard this also as an upper limit, however.
The $`\mathrm{H}_2\mathrm{O}`$-subtracted 5.0-6.5 $`\mu \mathrm{m}`$ wavelength region (Fig. 8) does not show the features detected toward high mass protostars (Schutte et al. sch96 (1996); sch98 (1998); Keane et al. kea00 (2000)). At 6.25 $`\mu \mathrm{m}`$ (not to confuse with the feature of $`\mathrm{NH}_3`$ ice at slightly shorter wavelength; see above), a feature has been associated with absorption by carbonaceous dust (PAH). At 5.83 $`\mu \mathrm{m}`$ a broad feature has been assigned to the C=O stretching mode of solid HCOOH, and a narrow feature of solid H<sub>2</sub>CO (Keane et al. kea00 (2000)). Scaling the features observed toward NGC~7538~:~IRS9 to the lower $`\mathrm{H}_2\mathrm{O}`$ ice band column density toward Elias~29, one would expect peak optical depths $`\tau _{5.83}=0.06`$ and $`\tau _{6.25}=0.03`$. Our spectra indicate upper limits to these features of $`\tau <`$0.03 (Fig. 8). Thus, in particular the 5.83 $`\mu \mathrm{m}`$ feature toward Elias~29 is significantly less pronounced compared to high mass protostars. Using the band strengths and typical widths given in Keane et al. (kea00 (2000)), we derive 3$`\sigma `$ column density upper limits of $`N(\mathrm{H}_2\mathrm{CO})<6\times 10^{16}`$ $`\mathrm{cm}^2`$, and $`N(\mathrm{HCOOH})<3\times 10^{16}`$ $`\mathrm{cm}^2`$. With abundance upper limits of 1–2% with respect to $`\mathrm{H}_2\mathrm{O}`$, these aldehydes are thus minor ice components. For comparison, toward high mass objects it is typically 3% or higher.
An absorption feature has been detected at 2042$`\pm `$$`\mathrm{cm}^1`$ (4.90 $`\mu \mathrm{m}`$) in lines of sight toward several massive protostars (Palumbo et al. pal97 (1997)). With a width FWHM=23$`\pm `$$`\mathrm{cm}^1`$, it has been ascribed to absorption by solid OCS. For Elias~29 this feature is not detected with a peak optical depth $`\tau <`$0.01 (3$`\sigma `$), corresponding to $`N`$(OCS)$`<1.5\times 10^{15}`$ $`\mathrm{cm}^2`$ or $`<0.05`$% of $`\mathrm{H}_2\mathrm{O}`$ ice. This upper limit is of the same order of magnitude as the detections in W 33A and Mon R2 : IRS2 (Palumbo et al. pal97 (1997)).
Finally, toward several high and low mass protostars a feature has been detected at $``$2166 $`\mathrm{cm}^1`$ (4.62 $`\mu \mathrm{m}`$) with a width FWHM$``$20 $`\mathrm{cm}^1`$ (Lacy et al. lac84 (1984); Tegler et al. teg95 (1995)). This feature is absent in Elias~29, with a peak optical depth $`\tau <`$0.01 (3$`\sigma `$). If this feature is caused by the C$``$N stretching mode in ‘XCN’, this corresponds to a column density $`N`$(XCN)$`<6.7\times 10^{15}`$ $`\mathrm{cm}^2`$, or less than 0.2% of $`\mathrm{H}_2\mathrm{O}`$ ice (applying $`A=3\times 10^{17}`$ cm molecule<sup>-1</sup>; Tegler et al. teg95 (1995)). This is considerably less than the detections made toward high mass objects (e.g. W 33A) and several low mass objects (Elias 18; L 1551 : IRS5; Tegler et al. teg95 (1995)). This feature has not been detected in the quiescent regions of the Taurus molecular cloud (Elias 16; Table 3). For a more elaborate discussion on this feature, and the proposed carriers, we refer to Pendleton et al. (pen99 (1999)).
#### 3.2.8 Silicates
The absorption bands of the Si–O stretching and bending modes of silicate dust are prominently present at 9.7 $`\mu \mathrm{m}`$ and 18 $`\mu \mathrm{m}`$ (Fig. 3). We derive a peak absorption optical depth of the 9.7 $`\mu \mathrm{m}`$ band $`\tau _{9.7}=1.38`$ (Fig. 4), which is in excellent agreement with the ground-based study of Hanner et al. (han95 (1995)). It is likely that this is a lower limit, since the absorption bands have been partly filled in with silicate emission from hot dust near the protostar. Modeling of the 9.7 $`\mu \mathrm{m}`$ silicate band toward Elias~29, including emission and absorption, shows that $`\tau _{9.7}`$ ranges between 1.51 and 3.38 for optically thick and thin emission respectively (Hanner et al. han95 (1995)). A better fit is obtained for optically thick emission. In contrast, for luminous protostars optically thin emission has been generally assumed. Using the relation $`\tau _{9.7}=1.4\tau _{9.7}(\mathrm{obs})+1.6`$ (Gillet et al. gil75 (1975); Willner et al. wil82 (1982)), yields $`\tau _{9.7}=3.53`$ for Elias~29.
For these values of $`\tau _{9.7}`$, the visual extinction $`A_\mathrm{V}`$ ranges between 28 and 65, assuming the standard relation $`A_\mathrm{V}/\tau _{9.7}`$=18.5 (Roche & Aitken roc84 (1984)). However, these limits are likely overestimated (30-50%), because of the anomalous extinction curve due to larger grains in the $`\rho `$ Oph molecular cloud (Bohlin et al. boh78 (1978); Martin & Whittet mar90 (1990)). Independent extinction determinations, such as $`A_\mathrm{V}<`$48 from the H–K broad band color and $`A_\mathrm{V}<`$80 from C<sup>18</sup>O observations (Wilking & Lada wil83 (1983)), do not help to solve this issue. Millimeter continuum observations (André & Montmerle and94 (1994)), and the near-infrared J–H color (Greene, priv. comm.), suggest a relatively low $`A_\mathrm{V}<`$30.
The total hydrogen column density $`N_\mathrm{H}=N`$(H i)+$`2N(\mathrm{H}_2)`$ is closely related to $`\tau _{9.7}`$, and, in contrast to the derivation of $`A_\mathrm{V}`$, the derived $`N_\mathrm{H}`$ is not strongly affected by the large grain size in $`\rho `$ Oph. Applying standard conversion factors for the diffuse ISM (Bohlin et al. boh78 (1978); Roche & Aitken roc84 (1984)), we find $`N_\mathrm{H}=0.51.2\times 10^{23}`$ $`\mathrm{cm}^2`$, depending on the applied $`\tau _{9.7}`$. To be consistent with studies of high mass protostars, we will assume in the abundance calculations, the value corresponding to optically thin silicate emission, i.e. the high limit $`N_\mathrm{H}=1.2\times 10^{23}`$ $`\mathrm{cm}^2`$ (Table 3).
### 3.3 Gas phase absorption lines
The high resolution 4.00–8.50 $`\mu \mathrm{m}`$ spectrum of Elias~29 shows an impressive number of narrow absorption lines of gaseous CO and $`\mathrm{H}_2\mathrm{O}`$ (Figs. 5 and 10). We determined local continuum points by hand and connected these, using a smooth cubic spline interpolation. Then the data were converted to optical depth scale, and the absorption lines were modeled, using the ro-vibrational spectra of gaseous CO and $`\mathrm{H}_2\mathrm{O}`$ described in Helmich (hel96 (1996)). These models assume the gas is in Local Thermodynamic Equilibrium (LTE), and has a single excitation temperature $`T_{\mathrm{ex}}`$. The absorption lines have a Voigt profile, and are Doppler broadened to a width $`b_\mathrm{D}`$ (=FWHM/$`2\sqrt{\mathrm{ln2}}`$). The line oscillator strengths are calculated from the HITRAN database (Rothman et al. rot92 (1992)). Finally, the spectrum is convolved with a Gaussian to the resolution of our observations ($`R=15002000`$). Thus, three parameters are varied to fit the observed absorption lines: the column density $`N`$, the Doppler parameter $`b_\mathrm{D}`$, and the excitation temperature $`T_{\mathrm{ex}}`$. Reliable column densities can only be derived if $`b_\mathrm{D}`$ is a priori known, which in many studies (like ours) is not the case, since the lines are unresolved. At low values of $`b_\mathrm{D}`$, the lines become easily optically thick, and much larger column densities are needed to fit the observed lines, compared to models with high $`b_\mathrm{D}`$ values, and optically thin lines.
We emphasize that our assumptions of collisional excitation, and LTE at a single $`T_{ex}`$ need not be valid. There is likely a temperature gradient along the line of sight, as expected for a protostellar envelope. The LTE assumption may not apply for the high rotational levels, which have high critical densities. Also, the energy levels may be pumped by infrared photons, rather than being collisionally excited. Bearing these caveats in mind, we will here focus on deriving CO and $`\mathrm{H}_2\mathrm{O}`$ gas column densities and temperatures using the LTE models.
#### 3.3.1 CO gas
The 4.4–5.0 $`\mu \mathrm{m}`$ region shows absorption lines of gas phase $`{}_{}{}^{12}\mathrm{CO}`$, up to rotational quantum number $`J_{\mathrm{low}}`$=33 in the R-branch, and $`J_{\mathrm{low}}`$=36 in the P-branch (Fig. 10). The P(1), P(2) and R(0) lines are blended with the CO ice band at 4.67 $`\mu \mathrm{m}`$ and the H i Pf $`\beta `$ emission line at 4.653 $`\mu \mathrm{m}`$. For all other absorption lines we determined equivalent widths to construct a rotation diagram. A rotation diagram gives a first impression of the temperature components present along the line of sight, as well as their column densities (or lower limits for optically thick lines). For technical details on constructing such a diagram we refer to Mitchell et al. (mit90 (1990)), and Boogert et al. (boo98 (1998)). The equivalent widths were converted to column densities, using the oscillator strengths of Goorvitch (goo94 (1994)). For $`{}_{}{}^{12}\mathrm{CO}`$ (Fig. 11), we find two regimes with very different slopes, corresponding to temperatures $`T_{\mathrm{rot}}=90\pm 45`$ K and $`T_{\mathrm{rot}}=1100\pm 300`$ K respectively (with 3$`\sigma `$ errors). However, the slopes of the R- and P-branch lines of the hot component are different (Fig. 11), resulting in $`T_{\mathrm{rot}}=1700\pm 420`$ K when fitting to the P-branch lines only. A possible explanation is that CO is excited by continuum photons rather than collisions. The rising continuum may lead to a higher $`T_{\mathrm{rot}}`$ for the P-branch with respect to the R-branch. This effect becomes stronger when the photons released after de-excitation of R-branch levels are re-absorbed in P-branch levels. Radiative excitation has also been used to explain the $`\mathrm{H}_2\mathrm{O}`$ ro-vibrational spectrum toward Orion BN/KL (Gonzalez-Alfonso et al. gon98 (1998)). The fact that the $`\mathrm{H}_2\mathrm{O}`$ P-branch lines are seen in emission for Orion BN/KL, rather than in absorption as for CO (and $`\mathrm{H}_2\mathrm{O}`$; Sect. 3.3.2) toward Elias~29, may reflect a different density gradient toward Elias~29, such that the photons are not able to escape the envelope. Additionally, collisional excitation in shocks may be of less importance in Elias~29 compared to Orion BN/KL. A more careful analysis is needed to discriminate between the radiative and collisional excitation models, and alternative explanations, such as non-LTE effects.
The column densities that we derive from the abscissa in the rotation diagram are $`N`$(CO)$`=1.7\times 10^{17}`$ and $`N`$(CO)$`=3.5\times 10^{17}`$ $`\mathrm{cm}^2`$ for the cold and hot CO components respectively. To better constrain the column densities and derive more reliable temperatures, one has to take into account optical depth effects, using the LTE model spectra discussed above. We chose to fit to the frequency range 2170–2290 $`\mathrm{cm}^1`$ ($`J_{\mathrm{low}}>7`$ in R-branch), thus minimizing the contribution from the cold CO component and contamination by $`{}_{}{}^{13}\mathrm{CO}`$ lines (see below). We find that good fits to these high R-branch lines are obtained only for line widths $`b_\mathrm{D}`$$`>3`$ $`\mathrm{km}\mathrm{s}^1`$. Sub-millimeter emission line studies indicate $`b_\mathrm{D}`$=3.6 $`\mathrm{km}\mathrm{s}^1`$ for CO J$`=65`$, but much lower values of $`b_\mathrm{D}`$=1.2 $`\mathrm{km}\mathrm{s}^1`$ for C<sup>18</sup>O J$`=10`$ and CS J$`=54`$ (Boogert, Hogerheijde, et al., in prep.). Indeed, studies of other sources have shown that, as a rule, infrared absorption lines are broader than sub-millimeter emission lines (van der Tak et al. tak99 (1999)). Figure 11 shows the $`\chi _\nu ^2`$ contour diagram of temperature versus column density for two values of the line width $`b_\mathrm{D}`$=5, and $`b_\mathrm{D}`$=10 $`\mathrm{km}\mathrm{s}^1`$. The best fitting models have temperatures $`T_{\mathrm{ex}}=1100\pm `$400 K, in good agreement with the rotation diagram. At $`b_\mathrm{D}`$=10 $`\mathrm{km}\mathrm{s}^1`$ the column density is well constrained to $`N`$(CO)=(1.3$`\pm `$ 0.5)$`\times 10^{18}`$ $`\mathrm{cm}^2`$, which is a factor of 3 larger compared to that derived from the rotation diagram. Thus at $`b_\mathrm{D}`$=10 $`\mathrm{km}\mathrm{s}^1`$ the lines are still somewhat optically thick. At lower $`b_\mathrm{D}`$=5 $`\mathrm{km}\mathrm{s}^1`$, the lines become very optically thick, and the column density is poorly constrained. Although the best fits with $`\chi _\nu ^2`$$`<3`$ have $`N`$(CO)=(8$`\pm `$4)$`\times 10^{18}`$ $`\mathrm{cm}^2`$ at $`T_{\mathrm{ex}}=650\pm 150`$ K, reasonable fits are obtained at any $`N`$(CO)$`>2\times 10^{18}`$ $`\mathrm{cm}^2`$ for this hot CO gas.
Several $`{}_{}{}^{13}\mathrm{CO}`$ lines can be seen in between the $`{}_{}{}^{12}\mathrm{CO}`$ P-branch lines (Fig. 13). At the resolution of our observations, the blending with the $`{}_{}{}^{12}\mathrm{CO}`$ lines hinders analyzing the much weaker $`{}_{}{}^{13}\mathrm{CO}`$ lines. But using several well separated lines, we were able to construct a rotation diagram (Fig. 12). We find that they result from cold gas at $`T_{\mathrm{rot}}=85\pm 57`$ K (3$`\sigma `$ error), in good agreement with the cold $`{}_{}{}^{12}\mathrm{CO}`$ gas temperature. In the optical thin case, the column density of this cold component is $`N`$($`{}_{}{}^{13}\mathrm{CO}`$)=$`(1.1\pm 0.2)\times 10^{17}`$ $`\mathrm{cm}^2`$. However, the detected $`{}_{}{}^{13}\mathrm{CO}`$ lines could still be optically thick. Therefore, we also modeled the $`{}_{}{}^{13}\mathrm{CO}`$ spectrum, and determine the $`\chi _\nu ^2`$ after subtraction of a good fitting hot $`{}_{}{}^{12}\mathrm{CO}`$ gas model (Figs. 13 and 12). In the optically thick case, such as for $`b_\mathrm{D}`$=2.5 $`\mathrm{km}\mathrm{s}^1`$, the column density can have a wide range $`N`$($`{}_{}{}^{13}\mathrm{CO}`$)=$`(2\pm 1.3)\times 10^{17}`$ $`\mathrm{cm}^2`$. Using the isotope abundance ratio $`{}_{}{}^{12}\mathrm{CO}`$/$`{}_{}{}^{13}\mathrm{CO}`$=80 (Boogert et al. boo00 (2000)), the inferred cold $`{}_{}{}^{12}\mathrm{CO}`$ column density is thus $`N`$($`{}_{}{}^{12}\mathrm{CO}`$)=$`(16\pm 10)\times 10^{18}`$ $`\mathrm{cm}^2`$. There is also evidence for $`{}_{}{}^{13}\mathrm{CO}`$ lines of warm gas ($`J_{\mathrm{low}}>9`$), but at low significance ($`2\sigma `$) and no reliable temperature or column density could be derived.
We conclude that the CO gas along the line of sight consists of two temperature components, $`T_{\mathrm{rot}}=90\pm 45`$ K and $`T_{\mathrm{rot}}=1100\pm 300`$ K. The column density of both components depends highly on the assumed line optical thickness (Table 2). Until the intrinsic line width is directly observed by very high spectral resolution observations, we can only give a lower limit of $`N`$(CO–hot)$`>2\times 10^{18}`$ $`\mathrm{cm}^2`$, while $`N`$(CO–cold) is not well constrained, i.e. $`(16\pm 10)\times 10^{18}`$ $`\mathrm{cm}^2`$. Given that $`N_\mathrm{H}=1.2\times 10^{23}`$ $`\mathrm{cm}^2`$ toward Elias~29, a total gas phase CO column density $`N`$(CO)=$`12\times 10^{18}`$ $`\mathrm{cm}^2`$ is expected, assuming that most of the gas along the line of sight is molecular and the conversion factor $`N(\mathrm{H}_2)/\mathrm{N}`$(CO)=5000 applies (Lacy et al. lac94 (1994)). Then, the ratio of hot to cold CO gas along the line of sight must be at least 0.2.
#### 3.3.2 H<sub>2</sub>O gas
We compare the numerous narrow absorption lines detected in the 5-7.3 $`\mu \mathrm{m}`$ spectral region of Elias~29 with model spectra of $`\mathrm{H}_2\mathrm{O}`$ vapor at various physical conditions (Fig. 14). Clearly, the many lines observed at wavelengths longer than $``$6.55 $`\mu \mathrm{m}`$ are explained by $`\mathrm{H}_2\mathrm{O}`$ vapor at a high temperature ($`T_{\mathrm{ex}}>`$100 K). On the other hand, the relative weakness of the lines observed in the range 6.55–6.65 $`\mu \mathrm{m}`$ imposes a strict upper limit to the temperature of this hot gas ($`T_{\mathrm{ex}}<`$1000 K). To further constrain the gas temperature, and the $`\mathrm{H}_2\mathrm{O}`$ vapor column density, we determined the $`\chi _\nu ^2`$ for a large number of models. Reasonable fits to the full 5–7.3 $`\mu \mathrm{m}`$ range are obtained for temperatures of $`T_{\mathrm{ex}}=350\pm 200`$ K. The column density is constrained to $`N`$=$`(7\pm 4)\times 10^{17}`$ $`\mathrm{cm}^2`$ for low line optical depths ($`b_\mathrm{D}`$$``$5). For narrower lines the column density can be an order of magnitude larger.
In a second approach, we test whether both hot and cold $`\mathrm{H}_2\mathrm{O}`$ vapor components could be present along the line of sight, much like the hot and cold CO components. We fitted the regions 5.5–5.8 and 6.55–7.3 $`\mu \mathrm{m}`$, which do not contain lines from the lowest rotational levels and thus are particularly sensitive to warm $`\mathrm{H}_2\mathrm{O}`$ vapor along the line of sight (Helmich et al. held96 (1996); Dartois et al. 1998b ). The excitation temperature of this gas is $`T_{\mathrm{ex}}=500\pm 300`$ K, with column densities similar to that of the single component model. In the high temperature regime ($`T_{\mathrm{ex}}`$500 K), the modeled line depths in the 6.0–6.5 $`\mu \mathrm{m}`$ region, tracing colder gas, are significantly underestimated. To determine the temperature and column density of this possible cold component, we fitted the sum of a good fitting hot gas model ($`T_{\mathrm{ex}}=`$500 K, $`N=5\times 10^{17}`$ $`\mathrm{cm}^2`$, $`b_\mathrm{D}`$=5.0 $`\mathrm{km}\mathrm{s}^1`$) and a grid of models at a wide range of physical conditions to the spectrum of Elias~29. Thus, here we assume that the lines of the hot and cold gas have different radial velocities and the optical depth spectra can simply be added. We find that indeed a significant amount of ’cold’ $`\mathrm{H}_2\mathrm{O}`$ vapor, at $`T_{\mathrm{ex}}<200`$ K may be present (Fig. 14). At $`T_{\mathrm{ex}}<100`$ K the column density exceeds the assumed hot $`\mathrm{H}_2\mathrm{O}`$ column density of $`N=5\times 10^{17}`$ $`\mathrm{cm}^2`$. For a line width of $`b_\mathrm{D}`$=5.0 $`\mathrm{km}\mathrm{s}^1`$, we find that $`N<1\times 10^{18}`$ $`\mathrm{cm}^2`$. At $`b_\mathrm{D}`$=2.5 $`\mathrm{km}\mathrm{s}^1`$, the column density of this cold $`\mathrm{H}_2\mathrm{O}`$ gas cannot be constrained.
To summarize, the lines in the 5–7.3 $`\mu \mathrm{m}`$ range are reasonably fitted with $`\mathrm{H}_2\mathrm{O}`$ models at $`T_{\mathrm{ex}}350\pm 200`$ K, and $`N`$=$`(7\pm 4)\times 10^{17}`$ $`\mathrm{cm}^2`$ at low line optical depths. For narrower lines ($`b_\mathrm{D}`$$`<`$$`\mathrm{km}\mathrm{s}^1`$), the column density can be an order of magnitude larger. In accordance with the gaseous CO along the line of sight, equally good fits are obtained with a two component model, where the cool component has $`T_{\mathrm{ex}}<200`$ K, and the warmer component $`T_{\mathrm{ex}}>`$500 K. The cool component is then at least as abundant as the warm $`\mathrm{H}_2\mathrm{O}`$ gas.
## 4 Discussion
### 4.1 Gas and solid state abundances
We have calculated line of sight averaged gas and solid state abundances toward Elias~29, by dividing the column densities derived in this paper over the total hydrogen column density $`N_\mathrm{H}=1.2\times 10^{23}`$ $`\mathrm{cm}^2`$ (Sect. 3.2)<sup>1</sup><sup>1</sup>1For actual, local, abundances in high mass protostars we refer to Boonman et al. bood00 (2000). We compare these abundances with a sample of sight-lines, spanning the range from dark cloud core to fairly evolved protostars (Table 3). As a tracer of ices in quiescent dark cloud material, we chose the object Elias 16, an evolved star by chance located behind the Taurus molecular cloud (e.g. Whittet et al. whi98 (1998)). Gas phase CO and $`\mathrm{H}_2\mathrm{O}`$ abundances in dense clouds were taken from ISO–LWS studies (Caux et al. cau99 (1999); Liseau & Olofsson lis99 (1999)). The least evolved protostar in our comparison sample is NGC~7538~:~IRS9. The infrared spectrum of this deeply embedded object is characterized by cold ice (Whittet et al. whi96 (1996)), and the gas phase temperatures and abundances indicate a very modest hot core (Mitchell et al. mit90 (1990); Boonman et al. bood00 (2000)). W 33A is more embedded than NGC~7538~:~IRS9, but does have a significant amount of warm gas along the line of sight (Mitchell et al. mit90 (1990); Lahuis & van Dishoeck lah00 (2000); Boonman et al. bood00 (2000)), and has a lower abundance of volatile ices (Tielens et al. tie91 (1991)). The most evolved object in our sample is GL 2591. It is a typical high mass hot core source, with low ice abundances and high gas temperatures. All these protostars are associated with infrared reflection nebulae, and have well developed high velocity molecular outflows (Mitchell et al. mit91 (1991); Bontemps et al. bon96 (1996); van der Tak et al. tak00 (2000)). Finally, it is important to note that all the comparison protostars are at least three orders of magnitude more luminous than Elias~29. This allows an investigation of the effect of low and high mass star formation on the molecular envelopes. An extensive comparison with low luminosity embedded objects is at present not possible, because their infrared gas and solid state characteristics have not been studied in such great detail.
The $`\mathrm{H}_2\mathrm{O}`$ and CO ice abundances decrease for the sequence of quiescent dense cloud to NGC~7538~:~IRS9, W 33A and GL 2591 (Table 3). At the same time, the gas phase $`\mathrm{H}_2\mathrm{O}`$ abundance, the gas phase CO and $`\mathrm{H}_2\mathrm{O}`$ temperatures, as well as the gas-to-solid ratios (Table 4), increase for these objects. All these effects can be explained by evaporation of the ice mantles and heating of the hot core. It has been suggested that the observed $`\mathrm{H}_2\mathrm{O}`$ gas may also have been newly formed by reactions of atomic O and $`\mathrm{H}_2`$ in warm conditions ($`T>200`$ K) in the central hot core or in shocks created by the outflow (e.g., van Dishoeck & Blake disb98 (1998)). However, the total (gas plus ice) $`\mathrm{H}_2\mathrm{O}`$ abundance decreases for the more evolved objects, indicating that $`\mathrm{H}_2\mathrm{O}`$ is destroyed rather than being newly formed (van Dishoeck dis98 (1998)). The low gas phase $`\mathrm{CO}_2`$ abundance in all sources indicates that this molecule is destroyed even more efficiently after evaporation from the grains (Boonman et al. bood00 (2000); Charnley & Kaufman cha00 (2000)).
In the proposed heating sequence, Elias~29 is placed after W 33A, and before GL 2591. However, the various ice band profiles ($`\mathrm{H}_2\mathrm{O}`$, CO, $`\mathrm{CO}_2`$, and 6.85 $`\mu \mathrm{m}`$) in Elias~29, indicate little thermal processing, resembling very much NGC~7538~:~IRS9, rather than W 33A or GL 2591. The combination of high gas phase abundances and temperatures, together with a lack of signatures of thermal processing in the ice bands, as seen in Elias~29, is remarkable and is not seen in high mass protostars. Geometric effects may play an important role in the evolution of molecular envelopes around low mass protostars (see Sect. 4.2).
Whereas thermal evaporation can explain the abundance variations of volatiles such as $`\mathrm{H}_2\mathrm{O}`$, CO, $`\mathrm{CO}_2`$, $`\mathrm{NH}_3`$, and $`\mathrm{CH}_4`$, other mechanisms are needed to explain the variations of solid $`\mathrm{CH}_3\mathrm{OH}`$, and XCN abundances among the sources in our sample (Table 3). It has been widely considered that XCN molecules are formed by energetic processing of icy grain mantles by stellar or cosmic ray induced far-ultraviolet radiation, or by bombardment with highly energetic particles (e.g. Lacy et al. lac84 (1984), Grim & Greenberg gri87 (1987), Allamandola et al. all88 (1988)). The high $`\mathrm{CH}_3\mathrm{OH}`$ abundances toward sources with deep XCN bands, and the apparent absence of $`\mathrm{CH}_3\mathrm{OH}`$ toward low mass protostars and dark clouds might suggest that the energetics of nearby massive stars is needed to produce $`\mathrm{CH}_3\mathrm{OH}`$ (Gibb et al. gib00 (2000)).
### 4.2 The structure of Elias~29
The variety of dust, gas and ice absorption and emission components presented here, and in the literature, allows us to construct an overall view of the structure of Elias~29. The scale on which the detected hot CO gas is present can be constrained when one assumes that the pure rotational high-J CO emission lines detected toward Elias~29 with ISO–LWS (Ceccarelli et al., in prep.) are emitted by the same hot gas. We fit these observed line fluxes, by assuming spontaneous, optically thin emission from an LTE level distribution, and leaving the size of the emitting region as a free parameter. For the range of column densities and temperatures found to fit the CO absorption lines (Fig. 11), we find diameters in the range 85–225 AU. Thus, the observed hot CO gas may be present in a hot core region with the size of a circumstellar disk. The gas could be concentrated in a high density photospheric layer above the disk. To sufficiently heat it by radiation from the central star, the disk needs to flare outwards, rather than being flat (e.g. Chiang & Goldreich chi97 (1997)). We cannot exclude however that the gas is present more uniformly in the hot core, at lower densities. It might then also be partly heated by shocks from the outflow close to the star. A more detailed modeling of the CO emission lines, including departures from LTE, optical depth corrections, and taking into account excitation by radiative pumping, are needed to further confine the location of the gas phase CO and $`\mathrm{H}_2\mathrm{O}`$ components (Ceccarelli et al., in prep.).
The lack of signatures of thermal processing in the ice bands, locates the ice in a region shielded from the central heating source. The ices could be present in a foreground cloud, an extended envelope, or a circumstellar disk seen close to edge-on. Millimeter wave continuum observations indicate an extended envelope, concentrated on the infrared source (FWHM=17<sup>′′</sup>; 2600 AU; André & Montmerle and94 (1994); Motte et al. mot98 (1998)). For a wide range of power law model fits to the far-infrared SED, André & Montmerle find that the envelope mass is 0.1 $`M_{}`$, with a volume averaged dust temperature of typically $`T=`$35 K. This temperature is high enough to evaporate the most volatile, apolar ices, but too low to induce ice crystallization. Hence, indeed the observed $`\mathrm{H}_2\mathrm{O}`$, $`\mathrm{CO}_2`$, and probably “6.85” $`\mu \mathrm{m}`$ ices could be associated with this extended envelope. Some of the apolar CO ice has evaporated in the envelope after the formation of the low mass protostar, as indicated by the significantly lower CO/$`\mathrm{H}_2\mathrm{O}`$ ice ratio toward Elias~29, compared to other sight-lines with little thermal processing in the ices (NGC~7538~:~IRS9, Elias 16; Table 3). In this picture, the detected apolar CO ice could thus be spatially separate from the other ices, perhaps in foreground clouds, or well shielded in a very cold disk.
Knowledge of the source structure is essential to interpret the observed solid and gas phase species. For example, if the ice is present in the disk, rather than in the envelope, we must see the disk in a near edge-on configuration. Is there independent evidence for the presence of a disk surrounding Elias~29 and what would be its orientation? The most direct view is provided by lunar occultation observations. A central object with diameter of $``$1 AU emits 90% of the 2.2 $`\mu \mathrm{m}`$ continuum emission (Simon et al. sim87 (1987)). The remaining 10% comes primarily from an object of 60 AU in diameter, which could be the hot part of a disk ($`T1000`$ K). The strongest spectroscopic disk indicator would be the presence of emission or absorption of vibrational overtone band heads of CO (e.g. Carr car89 (1989); Najita et al. naj96 (1996)). The 2.0–2.5 $`\mu \mathrm{m}`$ spectrum of Elias~29 does not show these features, in contrast to other protostars in $`\rho `$ Oph, such as WL 16 (Greene & Lada gre96 (1996)). However, the absence of CO overtone bands does not prove the absence of an (inner) disk (Calvet et al. cal91 (1991)). For example, the Herbig Ae object AB Aur does not have detected CO overtone bands, while high spatial resolution radio continuum and emission line observations provide strong evidence for the presence of a circumstellar disk around this object (Mannings & Sargent man97 (1997)).
AB Aur is an interesting comparison source, since it has the same luminosity as Elias~29 ($``$40 $`L_{}`$), and the SEDs of both objects are remarkably similar (Fig. 2). The flatness of the SED in AB Aur is well reproduced in flaring disk models, where the dust in the outer parts of the disk is more efficiently heated than in flat disks (Chiang & Goldreich chi97 (1997)). The disk is optically thick up to 100 $`\mu \mathrm{m}`$, and becomes optically thin at longer wavelengths where the SED drops steeply (e.g. van den Ancker et al. anc00 (2000)). The similarity of the SEDs does however not necessarily imply that Elias~29 is dominated by an optically thick disk as well. A flat SED could also be produced by the envelope, if it has a shallow power law density profile (index $`0.5`$; André & Montmerle and94 (1994)). This density profile is remarkably flat compared to high mass protostars (van der Tak et al. tak00 (2000); Dartois et al. 1998a ), and other low mass protostars (e.g. Hogerheijde & Sandell hog00 (2000)). Finally, flat energy distributions are also created by the combination of a disk and envelope. Here, the heated envelope irradiates the outer parts of the disk (Natta nat93 (1993)).
Without direct high resolution imaging, it is difficult to discriminate between these models. Assuming a given model however, the present observations put some constraints. In the disk scenario, its orientation would have to be closer to edge-on than face-on to explain the absorption line spectrum of Elias~29 (Chiang & Goldreich chi99 (1999)). In these models, an inclination larger than $``$70<sup>o</sup> can however be excluded, because this would give an SED that peaks in the far-infrared, in contrast to what is observed for Elias~29. Also, if the disk were edge-on, a higher absorbing column, perhaps an order of magnitude larger than the observed $`N_\mathrm{H}1.2\times 10^{23}`$ $`\mathrm{cm}^2`$ (Sect. 3.2) would be expected (Sekimoto et al. sek97 (1997)). An independent measure for $`N_\mathrm{H}`$ and the disk inclination is provided by the hard X-ray flux and spectrum, arising from hot gas in the magnetosphere. For Elias~29, a high $`N_\mathrm{H}2\times 10^{23}`$ $`\mathrm{cm}^2`$ is observed during X-ray flares, but $`N_\mathrm{H}`$ is a factor of 5 lower in quiescent phases (Kamata et al. kam97 (1997)). Perhaps the X-ray flares are formed low in the magnetosphere, and in the relatively high inclination of the disk, they trace higher column densities compared to X-rays formed in quiescent phases higher in the magnetosphere.
## 5 Conclusions and future work
The 1.2–195 $`\mu \mathrm{m}`$ spectrum of the low mass protostellar object Elias~29 in the $`\rho `$ Ophiuchi molecular cloud shows a wealth of absorption lines of gas and solid state molecules. Hot CO and $`\mathrm{H}_2\mathrm{O}`$ gas are detected ($`T_{\mathrm{ex}}>`$300 K) at rather high abundances, on scales of not more than a few hundred AU. The ice abundances are relatively low. In this respect, Elias~29 resembles luminous protostars with significantly heated cores, such as GL 2591. However, none of the many ice bands that are detected, i.e. from $`\mathrm{H}_2\mathrm{O}`$, CO, $`\mathrm{CO}_2`$, and the 6.85 $`\mu \mathrm{m}`$ band, shows outspoken signs of thermal processing. Again in comparison with luminous protostars, Elias~29 now resembles less evolved objects, such as NGC~7538~:~IRS9. Our combined gas and solid state analysis thus shows that high and low mass protostars heat their molecular envelopes in different ways. This may be related to their different structure, such as the presence of a circumstellar disk in low mass protostars. The hot gas of Elias~29 could be present on the surface of a flaring disk, which is efficiently heated by the central star. The ices toward Elias~29 must be well shielded in a circumstellar disk seen close to edge-on, or far away in the envelope.
Does this imply that in general the ices in the disks or outer envelopes of low mass protostars remain unaltered, both in composition and structure, during the process of star formation? Are these ices the building blocks of the early solar system and are they preserved in present day observed cometary nuclei? To date, no Class I protostar has been found with strong signs of crystalline ices (Boogert et al. boo00 (2000)). On the other hand, the presence of crystalline ices and silicates has been reported in several isolated, less embedded Herbig AeBe objects (Malfait et al. mal99 (1999)). This research needs to be extended to a larger sample of low mass protostars, in a range of evolutionary stages and luminosities. Furthermore, it is essential for the interpretation of the gas and solid state characteristics toward Elias 29 that the presence of a circumstellar disk, and its inclination are determined by future high spatial resolution infrared or millimeter continuum observations.
###### Acknowledgements.
We thank Tom Greene (NASA/Ames Research Center) for providing us the 1.1–2.4 $`\mu \mathrm{m}`$ spectrum of Elias~29 in electronic format, Willem Schutte (Leiden Observatory) for providing the $`\mathrm{H}_2\mathrm{O}`$:$`\mathrm{NH}_3`$ laboratory ice mixtures, and T.Y. Brooke (NASA/JPL) for the 3 $`\mu \mathrm{m}`$ spectrum of NGC~7538~:~IRS9. The referee D. Ward-Thompson is thanked for a number of useful comments. D.C.B.W. is funded by NASA through JPL contract no.961624 and by the NASA Exobiology and Long-Term Space Astrophysics programs (grants NAG5-7598 and NAG5-7884, respectively).
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# Far-UV FUSE spectroscopy of the O vi resonance doublet in Sand 2 (WO)
## 1 Introduction
Wolf-Rayet (WR) stars provide keys to our understanding of massive stellar evolution, nucleosynthesis processes and chemical enrichment of the ISM. Of these, the oxygen-sequence (WO) introduced by Barlow & Hummer (1982) is by far the rarest. Their high excitation oxygen emission lines are widely interpreted as revealing the late core helium-burning or possibly carbon-burning stage (Smith & Maeder 1991), of importance for constraining the controversial <sup>12</sup>C($`\alpha ,\gamma )^{16}`$O reaction rate.
In spite of searches in all Local Group galaxies, only six massive WO stars are known to date, namely Sand 1 (Sk 188) in the SMC, Sand 2 (BAT99–123) in the LMC, Sand 4 (WR 102), Sand 5 (WR 142) and MS4 (WR 30a) in our Galaxy, and DR1 in IC1613. Since O vi 3811–34Å is a primary WO classification diagnostic, with an equivalent width of up to 1700Å (Kingsburgh, Barlow & Storey (1995, hereafter KBS), observations of O vi 1032–38Å are keenly sought. However, its location in the far-UV has ruled out such observations to date. This situation has changed following the successful launch of the Far-Ultraviolet Spectroscopic Explorer (FUSE, Moos et al. 2000), which permits routine high dispersion far-UV spectroscopy of massive stars in the Magellanic Clouds.
In this Letter, we analyse FUSE spectroscopy of Sand 2 (Sanduleak 1971), alias Sk$``$68 145 = Brey 93 = BAT99-123 (Breysacher, Azzopardi & Testor 1999), together with Hubble Space Telescope (HST) and ground-based datasets.
## 2 Observations
Previously unpublished far-UV, UV and optical/near-IR spectroscopy of Sand 2 have been obtained with FUSE, HST and the Mt Stromlo and Siding Spring Observatory (MSSSO) 2.3m, respectively.
### 2.1 Far-UV spectroscopy
Sand 2 was observed by FUSE as part of the Early Release Observation programme X018 on 1999 Oct 31. A 8134 sec exposure of Sand 2 with the $`30^{^{\prime \prime }}\times 30^{^{\prime \prime }}`$ (LWRS) aperture provided data at $`R`$12,000 with the two Lithium Fluoride (LiF) channels, covering $`\lambda \lambda `$979–1187Å, obtained in time-tag (TTAG) mode. At this epoch, the two Silicon Carbide (SiC) channels, covering $`\lambda \lambda `$905–1104Å, were badly aligned so that SiC data were of poor quality.
Sand 2 data were processed through the pipeline data reduction, CALFUSE (version 1.6.8), and are shown in Fig. 1. The pipeline extracted the 1D spectrum, removed the background, and corrected for grating wobble and detector drifts. No corrections for astigmatism or flat fielding have been applied. From Fig. 1, the far-UV continuum flux of Sand 2 is low ($`5\times 10^{14}`$ erg s<sup>-1</sup> cm<sup>-2</sup> Å<sup>-1</sup>), with two principal stellar features, O vi $`\lambda `$1032-38 and C iv $`\lambda `$1168. The latter is blended with C iii $`\lambda `$1175, while C iv $`\lambda `$1108 and O vi $`\lambda `$1125 are present, though weak. The FUSE spectrum is affected by a multitude of interstellar absorption features, principally H i and H<sub>2</sub>, plus airglow emission features due to Ly$`\beta `$, \[O i\] and \[N i\].
### 2.2 Near-UV spectroscopy
Sand 2 was observed with the HST Faint Object Spectrograph (FOS) instrument during 1996 March (PI: D.J. Hillier, Program ID 5460). Exposures totalling 5980, 3780 and 1090 sec were obtained with the G130H, G190H and G270H gratings, respectively. Many important spectral features are identified in the HST/FOS dataset (see also KBS), principally C iv $`\lambda `$1550, C iii $`\lambda `$2297, He ii $`\lambda `$1640, O iv $`\lambda `$1340, O v $`\lambda `$1371, $`\lambda `$2781–87. Following Prinja, Barlow & Howarth (1990) we find $`v_{\mathrm{}}`$=4100 km s<sup>-1</sup> from C iv $`\lambda `$1548–51 (KBS derived 4500 km s<sup>-1</sup> from IUE spectroscopy).
### 2.3 Optical spectroscopy
We have used the Double Beam Spectrograph (DBS) at the 2.3m MSSSO telescope to observe Sand 2 on 1997 Dec 24–27. Use of a dichroic and the 300B, 600B and 316R gratings permitted simultaneous spectroscopy covering 3620–6085Å and 6410–8770Å for 1200 sec, plus 3240–4480Å and 8640–11010Å for 2500 sec. A 2<sup>′′</sup> slit and 1752$`\times `$532 pixel SiTE CCD’s provided a 2 pixel spectral resolution of $``$5Å. A standard data reduction was carried out, including absolute flux calibration, using wide slit spectrophotometry of HD 60753 (B3 IV) and $`\mu `$ Col (O9 V), plus atmospheric correction using HR 2221 (B8 V). Convolving our dataset with Johnson broad-band filter profiles reveals V=15.1 and B-V=+0.54 mag. Since these are contaminated by emission lines, we have convolved our data with Smith $`ubvr`$ narrow-band filters to reveal $`v`$=16.1, $`ub=`$0.19, $`bv=0.07`$ and $`vr=0.14`$ mag.
Our optical dataset confirms the spectral morphology previously presented and discussed by KBS, with exceptionally strong O iv $`\lambda `$3400, O iv $`\lambda `$3818, C iv+He ii $`\lambda `$4660 and C iv $`\lambda `$5801. Depending on the WO selection criteria, its spectral type is WO3 (Crowther, De Marco & Barlow 1998) or WO4 (KBS).
## 3 Quantitative analysis of Sand 2
Wolf-Rayet winds are so dense that non-LTE effects, spherical geometry and an expanding atmosphere are minimum assumptions.
### 3.1 Analysis technique
We employ the code of Hillier & Miller (1998), cmfgen, which iteratively solves the transfer equation in the co-moving frame subject to statistical and radiative equilibrium in an expanding, spherically symmetric and steady-state atmosphere. These models account for clumping, via a volume filling factor, $`f`$, and line blanketing, both of which have a significant effect on the physical properties of WC stars (e.g. Hillier & Miller 1999). Through the use of ‘super-levels’, extremely complex atoms can be included. For the present application, a total of 3,552 levels (combined into 796 super-levels), 60 depth points and 63,622 spectral lines of He i-ii, C ii-iv, O ii-vi, Ne ii-iv, Si iv, S iv-vi, Ar iii-v, and Fe iv-viii are considered simultaneously (see Dessart et al. 2000 for the source of atomic data used).
We adopt a form for the velocity law (Eqn 8 from Hillier & Miller 1999) such that two exponents are considered ($`\beta _1`$=1, $`\beta _2`$=50, $`v_{\mathrm{ext}}`$=2900 km s<sup>-1</sup>, $`v_{\mathrm{turb}}`$=100 km s<sup>-1</sup>), with the result that acceleration is modest at small radii, but continues to large distance (0.9$`v_{\mathrm{}}`$ is reached at 100$`R_{}`$).
Our spectroscopic analysis derives $`\dot{M}/\sqrt{f}`$, rather than $`\dot{M}`$ and $`f`$ individually, since line blending is severe in Sand 2. A series of models were calculated in which stellar parameters, $`T_{}`$, log $`(L/L_{})`$, $`\dot{M}/\sqrt{f}`$, C/He and O/He, were adjusted until the observed line strengths and spectroscopic fluxes were reproduced. We adopt 0.4$`Z_{}`$ abundances for Ne, Si, S, Ar and Fe. The distance to the LMC was assumed to be 51.2 kpc (Panagia et al. 1991).
The wind ionization balance is ideally deduced using isolated spectral lines from adjacent ionization stages of carbon (C iii $`\lambda `$2297, C iv $`\lambda `$5801-12) or oxygen (O iv $`\lambda `$3404-14, O v $`\lambda `$3144, O vi $`\lambda `$5290). In practice this was difficult to achieve for Sand 2 because of the severe line blending.
### 3.2 Stellar parameters
Our initial parameter study of Sand 2 revealed a fairly similar spectral appearance spanning 120kK$`T_{}`$170kK, with log $`(L/L_{})`$=5.28 and log ($`\dot{M}/M_{}`$yr$`{}_{}{}^{1})=`$4.94 fixed. The principal differences between these models are that (i) the observed strength of the O vi 3811–34Å doublet, plus lines in the red such as C iv 7700Å favour a high $`T_{}`$; (ii) the weakness of the O vi 1032–38Å doublet, as revealed by FUSE, favours a lower $`T_{}`$. Since other parameters, in particular abundances, are largely unaffected by these discrepancies, we shall adopt $`T_{}`$=150kK (i.e. $`R_{}`$=0.65$`R_{}`$). Note that higher luminosity models do produce significant effects, such as a dramatic weakening of C iv 5801–12Å emission.
Fig. 2 compares our synthetic spectrum with the observed far-UV, UV and optical spectroscopy of Sand 2. Our model is reddened by $`E(BV)`$=0.08 mag due to our Galaxy, obtained from the reddening map of Burstein & Heiles (1982). An additional LMC component of 0.11 mag was required, such that $`M_v=3.0`$ mag. In the absence of a far-UV extinction law, standard UV laws (Seaton 1979; Howarth 1983) are extrapolated for $`\lambda `$1200Å with (variable) influence on the fit quality to FUSE data.
Overall the observed spectrum of Sand 2 is very well reproduced by the model, with most He ii, C iii-iv and O iv-vi lines matched in strength and shape, except for O vi $`\lambda \lambda `$3811–34 (model too weak) and O vi $`\lambda \lambda `$1032–1038 (model too strong) as discussed above. In particular, the flat-topped nature of C iii $`\lambda `$2297 is well matched. It is clear that although WO stars have little or no C iii $`\lambda `$5696 (a classification diagnostic), other C iii lines are indeed present (Hillier 1989). Many spectral features in Sand 2 are due to blends because of the very broad spectral lines and the fact that recombination lines of O vi and C iv overlap with He ii lines. For example, the spectral feature at $`\lambda \lambda `$4650–4686 has principal contributors He ii $`\lambda `$4686, C iii $`\lambda \lambda `$4647-50, C iv $`\lambda `$4658 and C iv $`\lambda `$4685, while minor contributors include C iv $`\lambda `$4646, $`\lambda `$4689 and O vi $`\lambda `$4678.
He ii $`\lambda `$5412/C iv $`\lambda `$5471 provides an excellent diagnostic of C/He (e.g. Hillier & Miller 1998) for spectroscopic studies of WC stars. However, the large line widths of WO stars, and the fact that C iv $`\lambda `$5412 (14–8) contributes to He ii $`\lambda `$5412 (7–4) hinders the use of these lines (see inset box in Fig. 2). Instead, we are able to derive C/He$`0.7\pm `$0.2 by number from He ii (6–4) $`\lambda `$6560/C iv $`\lambda `$7700, although C iv (12–8) contributes to the former. Other C iv and He ii recombination lines show excellent agreement, except that C iv $`\lambda `$1107 is predicted to be stronger than FUSE observations reveal.
In contrast to WC stars, numerous oxygen recombination lines are present in the UV and optical spectra of WO stars (e.g. O vi $`\lambda `$5290, $`\lambda `$2070). From their strength relative to He and C recombination lines, we estimate O/He$`0.15_{0.05}^{+0.10}`$ by number. The weakness of lower ionization oxygen features, such as O iv $`\lambda `$1400 and O v $`\lambda `$3150 also argue against higher O/He ratios, although O v $`\lambda `$5590 favours O/He$``$0.25. Poor fits to O vi $`\lambda \lambda `$1032–38 and $`\lambda \lambda `$3811–34 are discussed above, while O v $`\lambda \lambda `$2781–87 and O iv $`\lambda \lambda `$3063–71 are systematically too weak for all models.
We present the predicted temperature structure of our Sand 2 model in Fig. 3. Although Sand 2 is an extremely hot star, its very high content of efficient C and O coolants (see Hillier 1989) directly results in a very cool ($`<`$10kK) outer wind ($`r/R_{}>`$ 100). Consequently, the ionization structure of metal species is predicted to be very stratified, such that O<sup>6+</sup> is the dominant ionization stage for $`r/R_{}<`$5, yet O<sup>3+</sup> is dominant for $`r/R_{}>`$25.
The high stellar temperature of Sand 2 implies a high bolometric correction ($``$5.7 mag) and consequently hard ionizing spectrum, such that 50% of the emergent photons have energies greater than 13.6eV (912Å Lyman edge), and 30% greater than 24.6eV (504Å He i edge). The ionizing fluxes in the H i, He i, and He ii continua are 10<sup>49.1</sup> s<sup>-1</sup>, 10<sup>48.8</sup> s<sup>-1</sup> and 10<sup>40.3</sup> s<sup>-1</sup>, respectively. WO stars are known to produce strong nebular He ii $`\lambda `$4686 emission in associated H ii regions (e.g. Kingsburgh & Barlow 1995), although the predicted number of He ii continuum ionizing photons in our Sand 2 model is much lower than those inferred from other WO stars.
## 4 Discussion
Abundances derived here qualitatively support the results from KBS, who used recombination line theory to derive C/He=0.5 and O/He=0.1 for Sand 2. WO stars are well suited to recombination studies for carbon, although oxygen is somewhat more problematic for recombination line studies, since the ionization structure is more complex (see Fig. 3), and few lines have available coefficients.
Gräfener, Hamann & Koesterke (1999) have carried out a detailed non-LTE spectroscopic analysis of Sand 2 (see Gräfener et al. 1998). Overall, we confirm their $`\dot{M}/\sqrt{f}`$ determination, but derive a higher luminosity (by 0.2 dex) and temperature (they derived $`T_{}`$=101kK), attributable to the incorporation of line blanketing. More significantly, we obtain systematically lower metal abundances (they estimated C/He=1.3 and O/He=1.2), such that the oxygen mass fraction for Sand 2 is only 16%, versus 50% according to Gräfener et al. (1999). We attribute this major revision to improved spectroscopic and atomic datasets (Gräfener et al. used simple C and O model atoms plus low S/N optical data). Our higher luminosity and clumpy wind conspire to revise the wind performance ratio, $`\dot{M}v_{\mathrm{}}/(L/c)`$, from 56 to 12.
Fig. 4 compares the luminosities and (C+O)/He ratios for Sand 2 with six LMC WC4 stars, updated from Dessart (1999). He improved upon similar work by Gräfener et al. (1998) using line blanketed, clumped models, revealing a greater range of carbon abundances, 0.1$``$C/He$``$0.3 (due to improved spectroscopy), systematically higher luminosities (because of blanketing and improved reddenings), and lower mass-loss rates (due to clumping). The carbon enrichment of Sand 2 is substantially higher than the WC4 stars, although O/He does not differ so greatly from the WC4 sample, for which O/He$``$0.08. This supports the suggestion by Crowther (1999) that unusually high oxygen enrichment may not be a pre-requisite for a WO classification.
We have superimposed (non-rotating) evolutionary tracks from Meynet et al. (1994) at 0.4$`Z_{}`$ for $`M_{\mathrm{initial}}`$=60, 85 and 120$`M_{}`$ on Fig. 4. These evolutionary models predict C/O$``$2 when C/He$``$0.7, in conflict with our determination of C/O$``$4. Better agreement is expected for evolutionary models in which rotation is accounted for, since these predict higher C/O ratios during the WC/WO phase (Maeder & Meynet 2000). From interior models, $`M_{\mathrm{initial}}60M_{}`$ for Sand 2, with a corresponding age of $``$3–4.3 Myr, such that a supernova is expected within the next 0.1–5$`\times 10^4`$ years. The stellar luminosity implies a current mass of 10$`M_{}`$ (Schaerer & Maeder 1992), such that the mean post-main sequence mass-loss rate is 2.5–5$`\times 10^5`$ $`M_{}`$ yr<sup>-1</sup>.
Quantitative analysis of other WO stars suffer from either (i) high interstellar reddening (WR 102, WR 142), (ii) complications because of binarity (Sand 1), or (iii) large distances (DR1). Nevertheless, we expect similar C and O enrichment to that derived here for Sand 2 (KBS).
This work is based on data obtained: (i) for the Guaranteed Time Team by the NASA-CNES-CSA FUSE mission, operated by the Johns Hopkins University, (ii) with the NASA/ESA Hubble Space Telescope obtained from the data archive at the STScI, which is operated by AURA under the NASA contract NAS5-26555; (iii) with the 2.3m Mount Stromlo and Siding Spring Observatory. Financial support is acknowledged from the Royal Society (PAC), NASA contract NAS5-32985 (U.S. FUSE participants), STScI grant GO-04550.01-92A, NASA grant NAG5–8211 (both JDH), PPARC (OD) and the UCL Perren Fund (LD).
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# Vassiliev Invariants of Legendrian, Transverse, and Framed Knots in Contact 3-manifolds
## 1. Introduction
In this section we describe the main results of the paper. (In case any of the terminology appears to be new to the reader, the corresponding definitions are given in the next section.)
If a contact structure on a $`3`$-manifold is cooriented, then every Legendrian knot (i.e. a knot that is everywhere tangent to the contact distribution) has a natural framing (a continuous normal vector field). Hence when studying Legendrian knots in such contact manifolds the main question is to distinguish those of them that realize isotopic framed knots.
Similarly if the contact structure is parallelized, then every transverse knot (i.e. a knot that is everywhere transverse to the contact distribution) also has a natural framing, and when studying transverse knots in such contact manifolds again the main question is to distinguish those of them that realize isotopic framed knots.
Vassiliev invariants proved to be an extremely useful tool in the study of framed knots, and the conjecture is that they are sufficient to distinguish all the isotopy classes of framed knots. Vassiliev invariants can also be easily defined in the categories of Legendrian and of transverse knots. In this paper we study the relationship between the groups of Vassiliev invariants of these three categories of knots, and explore when these invariants can be used to distinguish Legendrian knots that realize isotopic framed knots.
Consider a contact manifold $`M`$ with a cooriented contact structure. Fix an Abelian group $`𝒜`$, a connected component $``$ of the space of framed immersions of $`S^1`$ into $`M`$, and a connected component $``$ of the space of Legendrian immersions of $`S^1`$ into $`M`$. We study the relation between the groups of $`𝒜`$-valued Vassiliev invariants of framed knots from $``$ and of $`𝒜`$-valued Vassiliev invariants of Legendrian knots from $``$. The main results obtained in this paper are described below.
###### Theorem 1.
The groups of $`𝒜`$-valued Vassiliev invariants of Legendrian knots from $``$ and of framed knots from $``$ are canonically isomorphic, provided that the Euler class of the contact bundle vanishes on every $`\alpha H_2(M,)`$ realizable by a mapping $`\mu :S^1\times S^1M`$.
(See Theorem 3.1.3 and Proposition 3.1.4.)
Using Theorem 1 we show that:
###### Theorem 2.
The groups of $`𝒜`$-valued Vassiliev invariants of Legendrian knots from $``$ and of framed knots from $``$ are canonically isomorphic, provided that one of the following conditions holds:
the contact structure is tight;
the Euler class of the contact bundle is in the torsion of $`H^2(M,)`$ (in particular if the Euler class is zero).
the contact manifold is closed and admits a metric of negative sectional curvature.
(See Sections 3.1.5 and 3.1.9 and Theorem 3.1.10.)
As a corollary, we get that for any surface $`F`$ the group of finite order Arnold’s $`J^+`$-type invariants of wave fronts on $`F`$ is isomorphic to the group of Vassiliev invariants of framed knots in the spherical cotangent bundle $`ST^{}F`$ of $`F`$.
Previously the isomorphism of the groups of Vassiliev invariants of Legendrian and of framed knots was known only in the case where $`𝒜=`$ and $`M`$ is the standard contact $`^3`$ (result of D. Fuchs and S. Tabachnikov ) or the standard contact solid-torus (result of J. Hill ). The proofs of these isomorphisms were based on the fact that for the $``$-valued Vassiliev invariants of framed knots in these manifolds there exists a universal Vassiliev invariant also known as the Kontsevich integral. (Currently the existence of the Kontsevich integral is known only for a total space of an $`^1`$-bundle over a compact oriented surface with boundary, see the paper of Andersen, Mattes, and Reshetikhin.)
Thus the approach used in and to show the isomorphism of the groups of Vassiliev invariants is not applicable for almost all contact $`3`$-manifolds and Abelian groups $`𝒜`$ and our results appear to be a strong generalization of the results of Fuchs, Tabachnikov and Hill.
We also construct the first examples where Vassiliev invariants can be used to distinguish Legendrian knots that realize isotopic framed knots and are homotopic as Legendrian immersions. These are also the first examples where the groups of Vassiliev invariants of Legendrian and of framed knots from the corresponding components of the spaces of Legendrian and of framed immersions are not canonically isomorphic.
###### Theorem 3.
The manifold $`S^1\times S^2`$ admits infinitely many cooriented contact structures for which there exist Legendrian knots that can be distinguished by $``$-valued Vassiliev invariants even though they realize isotopic framed knots and are homotopic as Legendrian immersions.
(See Theorem 4.1.1 and Theorem 4.2.2 in which the similar result is proved for any orientable total space of an $`S^1`$-bundle over a nonorientable surface of a sufficiently high genus.)
For transverse knots we obtain the following result (see Theorem 3.2.2):
###### Theorem 4.
Let $`(M,C)`$ be a contact manifold with a parallelized contact structure, then the groups of $`𝒜`$-valued Vassiliev invariants of transverse and of framed knots (from the corresponding components of the spaces of transverse and of framed immersions) are canonically isomorphic.
## 2. Conventions and definitions.
In this paper $`𝒜`$ is an Abelian group (not necessarily torsion free), and $`M`$ is a connected oriented $`3`$-dimensional Riemannian manifold (not necessarily compact).
A contact structure on a $`3`$-dimensional manifold $`M`$ is a smooth field $`\{C_xT_xM|xM\}`$ of tangent $`2`$-dimensional planes, locally defined as a kernel of a differential $`1`$-form $`\alpha `$ with non-vanishing $`\alpha d\alpha `$. A manifold with a contact structure possesses the canonical orientation determined by the volume form $`\alpha d\alpha `$. The standard contact structure in $`^3`$ is the kernel of the $`1`$-form $`\alpha =ydxdz`$.
A contact element on a manifold is a hyperplane in the tangent space to the manifold at a point. For a surface $`F`$ we denote by $`ST^{}F`$ the space of all cooriented (transversally oriented) contact elements of $`F`$. This space is the spherical cotangent bundle of $`F`$. Its natural contact structure is the distribution of tangent hyperplanes given by the condition that the velocity vector of the incidence point of a contact element belongs to the element.
A contact structure is cooriented if the $`2`$-dimensional planes defining the contact structure are continuously cooriented (transversally oriented). A contact structure is oriented if the $`2`$-dimensional planes defining the contact structure are continuously oriented. Since every contact manifold has a natural orientation we see that every cooriented contact structure is naturally oriented and every oriented contact structure is naturally cooriented. A contact structure is parallelizable (parallelized) if the $`2`$-dimensional vector bundle $`\{C_x\}`$ over $`M`$ is trivializable (trivialized). Since every contact manifold has a canonical orientation, one can see that every parallelized contact structure is naturally cooriented. A contact structure $`C`$ on a manifold $`M`$ is said to be overtwisted if there exists a $`2`$-disk $`D`$ embedded into $`M`$ such that the boundary $`D`$ is tangent to $`C`$ while the disk $`D`$ is transverse to $`C`$ along $`D`$. Not overtwisted contact structures are called tight.
A curve in $`M`$ is an immersion of $`S^1`$ into $`M`$. (All curves have the natural orientation induced by the orientation of $`S^1`$.) A framed curve in $`M`$ is a curve in $`M`$ equipped with a continuous unit normal vector field.
A Legendrian curve in a contact manifold $`(M,C)`$ is a curve in $`M`$ that is everywhere tangent to $`C`$. If the contact structure on $`M`$ is cooriented, then every Legendrian curve has a natural framing given by the unit normals to the planes of the contact structure that point in the direction specified by the coorientation.
To a Legendrian curve $`K_l`$ in a contact manifold with a parallelized contact structure one can associate an integer that is the number of revolutions of the direction of the velocity vector of $`K_l`$ (with respect to the chosen frames in $`C`$) under traversing $`K_l`$ according to the orientation. This integer is called the Maslov number of $`K_l`$. The set of Maslov numbers enumerates the set of the connected components of the space of Legendrian curves in $`^3`$ (cf. 2.0.1).
A transverse curve in a contact manifold $`(M,C)`$ is a curve in $`M`$ that is everywhere transverse to $`C`$. If the contact structure on $`M`$ is parallelized, then a transverse curve has a natural framing given by the unit normals corresponding to the projections of the first of the two coordinate vectors of the contact planes on the $`2`$-planes orthogonal to the velocity vectors of the curve. A transverse curve in a contact manifold with a cooriented contact structures is said to be positive if at every point the velocity vector of the curve points into the coorienting half-plane, and it is said to be negative otherwise. There are two connected components of the space of transverse curves in $`^3`$, they consist of positive and negative transverse curves respectively. In general if $`(M,C)`$ is a contact manifold with a cooriented contact structure, then every connected component of the space of unframed curves contains two connected components of the space of transverse curves. They consist of positive and negative transverse curves respectively.
A knot (framed knot) in $`M`$ is an embedding (framed embedding) of $`S^1`$ into $`M`$. In a similar way we define Legendrian and transverse knots in $`M`$.
A singular (framed) knot with $`n`$ double points is a curve (framed curve) in $`M`$ whose only singularities are $`n`$ transverse double points. An isotopy of a singular (framed) knot with $`n`$ double points is a path in the space of singular (framed) knots with $`n`$ double points under which the preimages of the double points on $`S^1`$ change continuously.
An $`𝒜`$-valued framed (resp. Legendrian, resp. transverse) knot invariant is an $`𝒜`$-valued function on the set of the isotopy classes of framed (resp. Legendrian, resp. transverse) knots.
A transverse double point $`t`$ of a singular knot can be resolved in two essentially different ways. We say that a resolution of a double point is positive (resp. negative) if the tangent vector to the first strand, the tangent vector to the second strand, and the vector from the second strand to the first form the positive $`3`$-frame. (This does not depend on the order of the strands). If the singular knot is Legendrian (resp. transverse), then these resolution can be made in the category of Legendrian (resp. transverse) knots.
A singular framed (resp. Legendrian, resp. transverse) knot $`K`$ with $`(n+1)`$ transverse double points admits $`2^{n+1}`$ possible resolutions of the double points. The sign of the resolution is put to be $`+`$ if the number of negatively resolved double points is even, and it is put to be $``$ otherwise. Let $`x`$ be an $`𝒜`$-valued invariant of framed (resp. Legendrian, resp. transverse) knots. The invariant $`x`$ is said to be of finite order (or Vassiliev invariant) if there exists a nonnegative integer $`n`$ such that for any singular knot $`K_s`$ with $`(n+1)`$ transverse double points the sum (with appropriate signs) of the values of $`x`$ on the nonsingular knots obtained by the $`2^{n+1}`$ resolutions of the double points is zero. An invariant is said to be of order not greater than $`n`$ (of order $`n`$) if $`n`$ can be chosen as the integer in the definition above. The group of $`𝒜`$-valued finite order invariants has an increasing filtration by the subgroups of the invariants of order $`n`$.
###### 2.0.1.
$`h`$-principle for Legendrian curves. For $`(M,C)`$ a contact manifold with a cooriented contact structure, we put $`CM`$ to be the total space of the fiberwise spherization of the contact bundle, and we put $`\mathrm{pr}:CMM`$ to be the corresponding locally trivial $`S^1`$-fibration. The $`h`$-principle proved for the Legendrian curves by M. Gromov (, pp.338-339) says that the space of Legendrian curves in $`(M,C)`$ is weak homotopy equivalent to the space of free loops $`\mathrm{\Omega }CM`$ in $`CM`$. The equivalence is given by mapping a point of a Legendrian curve to the point of $`CM`$ corresponding to the direction of the velocity vector of the curve at this point. In particular the $`h`$-principle implies that the set of the connected components of the space of Legendrian curves can be naturally identified with the set of the conjugacy classes of elements of $`\pi _1(CM)`$.
###### 2.0.2.
Description of Legendrian and of transverse knots in $`^3`$. The contact Darboux theorem says that every contact $`3`$-manifold $`(M,C)`$ is locally contactomorphic to $`^3`$ with the standard contact structure that is the kernel of the $`1`$-form $`\alpha =ydxdz`$. A chart in which $`(M,C)`$ is contactomorphic to the standard contact $`^3`$ is called a Darboux chart.
Transverse and Legendrian knots in the standard contact $`^3`$ are conveniently presented by the projections into the plane $`(x,z)`$. Identify a point $`(x,y,z)^3`$ with the point $`(x,z)^2`$ furnished with the fixed direction of an unoriented straight line through $`(x,z)`$ with the slope $`y`$. Then the curve in $`^3`$ is a one parameter family of points with (non-vertical) directions in $`^2`$.
A curve in $`^3`$ is transverse if and only if the corresponding curve in $`^2`$ is never tangent to the chosen directions along itself.
While a generic regular curve has a regular projection into the $`(x,z)`$-plane, the projection of a generic Legendrian curve into the $`(x,z)`$-plane has isolated critical points (since all the planes of the contact structure are parallel to the $`y`$-axis). Hence the projection of a generic Legendrian curve may have cusps. A curve in $`^3`$ is Legendrian if and only if the corresponding planar curve with cusps is everywhere tangent to the field of directions. In particular this field is determined by the curve with cusps.
## 3. Isomorphisms of the groups of Vassiliev invariants of Legendrian, of transverse, and of framed knots.
### 3.1. Isomorphism between the groups of order $`n`$ invariants of Legendrian and of framed knots
Let $`(M,C)`$ be a contact manifold with a cooriented contact structure. Let $``$ be a connected component of the space of Legendrian curves in $`M`$, and let $``$ be the connected component of the space of framed curves that contains $``$. (Such a component exists because a Legendrian curve in a manifold with a cooriented contact structure is naturally framed, and a path in the space of Legendrian curves corresponds to a path in the space of framed curves.) Let $`V_n^{}`$ (resp. $`W_n^{}`$) be the group of $`𝒜`$-valued order $`n`$ invariants of Legendrian (resp. framed) knots from $``$ (resp. from $``$). Clearly every invariant $`yW_n^{}`$ restricted to the category of Legendrian knots in $``$ is an element $`\varphi (y)V_n^{}`$. This gives a homomorphism $`\varphi :W_n^{}V_n^{}`$.
###### Theorem 3.1.1.
Let $`(M,C)`$ be a contact manifold with a cooriented contact structure. Let $``$ be a connected component of the space of Legendrian curves in $`M`$, and let $``$ be the connected component of the space of framed curves that contains $``$. Then the following two statements a and b are equivalent.
$`x(K_1)=x(K_2)`$ for any $`xV_n^{}`$ and any knots $`K_1,K_2`$ representing isotopic framed knots.
$`\varphi :W_n^{}V_n^{}`$ is a canonical isomorphism.
If the mapping from the isotopy classes of Legendrian knots in $``$ to the isotopy classes of framed knots in $``$ is surjective, then the proof of Theorem 3.1.1 is obvious. However in general this mapping is not surjective and the proof of Theorem 3.1.1 is given in Section 5.2. The famous Bennequin inequality shows that this mapping is not surjective even in the case where $`M`$ is the standard contact $`^3`$.
Theorem 3.1.1 implies that to obtain the isomorphism between the groups $`W_n^{}`$ and $`V_n^{}`$ it suffices to show that statement $`𝐚`$ of Theorem 3.1.1 is true for the connected components $``$ and $``$ of the spaces of Legendrian and of framed curves.
###### 3.1.2.
Condition $`()`$. In D. Fuchs and S. Tabachnikov showed that statement $`𝐚`$ holds for all the connected components of the space of Legendrian curves when the ambient manifold is the standard contact $`^3`$ and the group $`𝒜`$ is $``$. (One can verify that the proof of this Theorem of Fuchs and Tabachnikov goes through for $`𝒜`$ being any Abelian group.) They later observed that since their proof of this fact is mostly local, the similar fact should be true for a big class of contact manifolds.
However in fact the proof of their theorem is not completely local and is also based on the existence of well-defined Bennequin invariant and Maslov number for a Legendrian knot in $`^3`$. In general the Bennequin invariant is not well-defined unless the knot is zero-homologous and the Maslov number is not well-defined unless either the knot is zero-homologous or the contact structure is parallelizable. Thus the generalization of this Theorem to the case of manifolds other than $`^3`$ meets certain difficulties. (And in fact the corresponding result does not hold for a big class of contact manifolds, see Section 4.)
By analyzing the proof of the Theorem of Fuchs and Tabachnikov (see Section 5.3) we get that it can be generalized to the case of an arbitrary contact $`3`$-manifold with a cooriented contact structure, provided that the connected component $``$ (containing $``$) satisfies the following
Condition $`()`$: the connected component $``$ of the space of framed curves contains infinitely many components of the space of Legendrian curves. (See Proposition 3.1.4 for the homological interpretation of condition $`()`$.)
This generalization of the result of Fuchs and Tabachnikov and Theorem 3.1.1 imply the following Theorem.
###### Theorem 3.1.3.
Let $`(M,C)`$ be a contact manifold with a cooriented contact structure, and let $``$ be a connected component of the space of Legendrian curves in $`M`$. Let $``$ be the connected component of the space of framed curves that contains $``$. Let $`V_n^{}`$ (resp. $`W_n^{}`$) be the group of $`𝒜`$-valued order $`n`$ invariants of Legendrian (resp. framed) knots from $``$ (resp. from $``$). Then the groups $`V_n^{}`$ and $`W_n^{}`$ are canonically isomorphic, provided that $``$ satisfies condition $`()`$.
Now we give a homological interpretation of condition $`()`$.
###### Proposition 3.1.4.
Let $`(M,C)`$ be a contact manifold with a cooriented contact structure, let $`\chi _CH^2(M)`$ be the Euler class of the contact bundle, and let $``$ be a component of the space of framed curves in $`M`$. Then $``$ does not satisfy condition $`()`$ if and only if there exists $`\alpha H_2(M,)`$ such that $`\chi _C(\alpha )0`$ and $`\alpha `$ is realizable by a mapping $`\mu :S^1\times S^1M`$ with the property that $`\mu |_{1\times S^1}`$ is a loop free homotopic to loops realized by curves from $``$.
For the Proof of Proposition 3.1.4 see Subsection 5.4.
###### Remark 3.1.5.
Some immediate corollaries of Theorem 3.1.3 and the generalization of the Theorem of Fuchs and Tabachnikov about the isomorphism of the groups of the $``$-valued Vassiliev invariants in the case of $`M=^3`$. Proposition 3.1.4 implies that if the contact structure is parallelizable (and hence the Euler class of the contact bundle is zero) then all the connected components of the space of framed curves satisfy condition $`()`$. Applying Theorem 3.1.3 we conclude that for any Abelian group $`𝒜`$ and for every connected component of the space of Legendrian curves $``$ and for the containing it component of the space of framed curves $``$ the groups $`V_n^{}`$ and $`W_n^{}`$ of $`𝒜`$-valued Vassiliev invariants are canonically isomorphic.
Clearly the value of the Euler class of the contact bundle is zero if $`M`$ is an integer homology sphere. Hence for any Abelian group $`𝒜`$ we obtain the isomorphism of the groups $`V_n^{}`$ and $`W_n^{}`$ of $`𝒜`$-valued Vassiliev invariants. This generalizes the Theorem of D. Fuchs and S. Tabachnikov saying that for the standard contact $`^3`$ and for $`𝒜=`$ the quotient groups $`V_n^{}/V_{n1}^{}`$ and $`W_n^{}/W_{n1}^{}`$ are canonically isomorphic.
The proof of this Theorem of Fuchs and Tabachnikov was based on the fact that for the $``$-valued Vassiliev invariants of framed knots in $`^3`$ there exists the universal Vassiliev invariant constructed by T. Q. T. Le and J. Murakami . (For unframed knots in $`^3`$ the construction of the universal Vassiliev invariant is the classical result of M. Kontsevich , and the invariant itself is the famous Kontsevich integral.) The existence of the universal Vassiliev invariant is currently known only for a very limited collection of $`3`$-manifolds, and only for $`𝒜`$ being $``$, $``$, or $``$. (Andersen, Mattes, Reshetikhin proved its existence in the case where $`𝒜=`$ and $`M`$ is the total space of an $`^1`$-bundle over a compact oriented surface $`F`$ with $`F\mathrm{}`$.)
Thus the approach used in to show the isomorphism of the quotient groups is not applicable for almost all contact $`3`$-manifolds and Abelian groups $`𝒜`$, and Theorem 3.1.3 appears to be a strong generalization of the result of Fuchs and Tabachnikov.
###### Remark 3.1.6.
Let $`(M,C)`$ be a contact manifold with a cooriented contact structure, and let $``$ be a connected component of the space of framed curves in $`M`$. Theorem 3.1.3 implies that the group of $`𝒜`$-valued order $`n`$ invariants of Legendrian knots from a connected component $``$ of the space of Legendrian curves does not depend on the choice of a cooriented contact structure, provided that for this choice $``$ satisfies condition $`()`$. And hence in these cases the group can not be used to distinguish cooriented contact structures on $`M`$. (See Remark 3.1.5 and Theorems 3.1.8 and 3.1.10 for the list of cases when the connected components of the space of framed curves are known to satisfy condition $`()`$.)
###### 3.1.7.
Finite order Arnold’s $`J^+`$-type invariants of wave fronts on surfaces. A very interesting class of contact manifolds satisfying the conditions of Theorem 3.1.3 is formed by the spherical cotangent bundles $`ST^{}F`$ of surfaces $`F`$ with the natural contact structure on $`ST^{}F`$ (see 2). The theory of the invariants of Legendrian knots in $`ST^{}F`$ is often referred to as the theory of Arnold’s $`J^+`$-type invariants of fronts on a surface $`F`$. The natural contact structure on $`ST^{}F`$ is cooriented. (The coorientation is induced from the coorientation of the contact elements of $`F`$.) One can verify that for orientable $`F`$ the standard contact structure on $`ST^{}F`$ is parallelizable, and hence all the components of the space of framed curves satisfy condition $`()`$. If $`F`$ is not orientable, then the standard cooriented contact structure on $`ST^{}F`$ is not parallelizable, but one can still verify (cf. Proposition 8.2.4 ) that every connected component of the space of framed curves satisfies condition $`()`$. Hence for any Abelian group $`𝒜`$ and for any surface $`F`$ we obtain the canonical isomorphism of the groups of $`𝒜`$-valued order $`n`$ invariants of Legendrian and of framed knots (from the corresponding components of the spaces of Legendrian and of framed curves in $`ST^{}F`$ with the standard contact structure). Or equivalently we get that the groups of $`𝒜`$-valued order $`n`$ $`J^+`$-type invariants of fronts on $`F`$ and of $`𝒜`$-valued order $`n`$ invariants of framed knots in $`ST^{}F`$ (from the corresponding components of the two spaces) are canonically isomorphic.
Previously it was known that for $`F=^2`$ and $`𝒜=`$ the quotient groups $`V_n^{}/V_{n1}^{}`$ and $`W_n^{}/W_{n1}^{}`$ are canonically isomorphic. The proof of this result of J. W. Hill was based on the fact that for the $``$-valued Vassiliev invariants of framed knots in $`ST^{}^2`$ there exists the universal Vassiliev invariant constructed by V. Goryunov . (For unframed knots in $`^3`$ the existence of the universal Vassiliev invariant is the classical result of M. Kontsevich , and the invariant itself is the famous Kontsevich integral.) Our results generalize the result of J. W. Hill (even in the case of $`M=ST^{}^2)`$.
The following Theorem describes another big class of contact manifolds for which the groups of Vassiliev invariants of Legendrian and of framed knots (from the corresponding components of the two spaces of curves) are canonically isomorphic.
###### Theorem 3.1.8.
Let $`(M,C)`$ be a contact manifold (with a cooriented contact structure) such that $`\pi _2(M)=0`$, and for every mapping $`\mu :T^2M`$ of the two-torus the homomorphism $`\mu _{}:\pi _1(T^2)\pi _1(M)`$ is not injective. Then all the components of the space of framed curves in $`M`$ satisfy condition $`()`$, and hence the groups of $`𝒜`$-valued order $`n`$ invariants of Legendrian and of framed knots (from the corresponding components of the spaces of Legendrian and framed curves) are canonically isomorphic.
For the Proof of Theorem 3.1.8 see Subsection 5.5.
###### 3.1.9.
The isomorphism of the groups of Vassiliev invariants in the case of closed manifolds admitting a metric of negative sectional curvature and other corollaries of Theorem 3.1.8.
Let $`M`$ be a closed manifold admitting a metric of negative sectional curvature. A well-known Theorem by A. Preissman (see pp. 258-265) says that every nontrivial commutative subgroup of the fundamental group of a closed $`3`$-dimensional manifold of negative sectional curvature is infinite cyclic. Hence for every mapping $`\mu :T^2M`$ the kernel of $`\mu _{}:\pi _1(T^2)=\pi _1(M)`$ is nontrivial. It is also known that the universal covering of such $`M`$ is diffeomorphic to $`^3`$, and hence $`\pi _2(M)=0`$. Thus every closed manifold $`M`$ admitting a metric of negative sectional curvature satisfies all the conditions of Theorem 3.1.8 and for an arbitrary cooriented contact structure on such $`M`$ we obtain the isomorphism of the groups of $`𝒜`$-valued order $`n`$ invariants of Legendrian and of framed knots from the corresponding components of the spaces of Legendrian and of framed curves.
Another important class of contact manifolds for which every connected component of the space of framed curves satisfies condition $`()`$ is formed by contact manifolds with a tight contact structure.
The following Theorem appeared as a result of discussions of Stefan Nemirovski and the author.
###### Theorem 3.1.10.
Let $`(M,C)`$ be a contact manifold with a tight cooriented contact structure. Then all the components of the space of framed curves in $`M`$ satisfy condition $`()`$, and hence the groups of $`𝒜`$-valued order $`n`$ invariants of Legendrian and of framed knots (from the corresponding components of the spaces of Legendrian and framed curves) are canonically isomorphic.
For the Proof of Theorem 3.1.10 see Subsection 5.6.
### 3.2. Isomorphisms between the groups of Vassiliev invariants of transverse and of framed knots
Let $`M`$ be a contact manifold with a parallelized contact structure $`C`$. Let $`𝒯`$ be a connected component of the space of transverse curves in $`(M,C)`$, and let $``$ be the connected component of the space of framed curves that contains $`𝒯`$. (Such a component exists because a transverse curve in a manifold with a parallelized contact structure is naturally framed, and a path in the space of transverse curves corresponds to a path in the space of framed curves.) Let $`V_n^𝒯`$ (resp. $`W_n^{}`$) be the group of $`𝒜`$-valued order $`n`$ invariants of transverse (resp. framed) knots from $`𝒯`$ (resp. from $``$). Clearly every invariant $`yW_n^{}`$ restricted to the category of transverse knots in $`𝒯`$ is an element $`\varphi (y)V_n^𝒯`$. This gives a homomorphism $`\varphi :W_n^{}V_n^𝒯`$.
###### Theorem 3.2.1.
Let $`(M,C)`$ be a contact manifold with a parallelized contact structure. Let $`𝒯`$ be a connected component of the space of transverse curves in $`(M,C)`$, and let $``$ be the component of the space of framed curves that contains $`𝒯`$. Then the following two statements a and b are equivalent.
$`x(K_1)=x(K_2)`$ for any $`xV_n^𝒯`$ and any knots $`K_1,K_2𝒯`$ representing isotopic framed knots.
$`\varphi :W_n^{}V_n^𝒯`$ is a canonical isomorphism.
The proof of Theorem 3.2.1 is analogous to the Proof of Theorem 3.1.1.
Similar to the case of Theorem 3.1.1 the proof of Theorem 3.2.1 becomes obvious if the mapping from the isotopy classes of transverse knots in $`𝒯`$ to the isotopy classes of framed knots in $``$ is surjective. However in general this mapping is not surjective and to obtain the proof of Theorem 3.2.1 one follows the ideas of the proof of Theorem 3.1.1 (The famous Bennequin inequality shows that this mapping is not surjective even for the standard contact $`^3`$.)
Thus to obtain the isomorphism between the groups $`W_n^{}`$ and $`V_n^𝒯`$ it suffices to show that statement $`𝐚`$ of Theorem 3.2.1 is true for the connected components $`𝒯`$ and $``$ of the spaces of transverse and of framed curves.
In D. Fuchs and S. Tabachnikov showed that statement $`𝐚`$ holds for all the connected components of the space of transverse curves in the case where $`M`$ is the standard contact $`^3`$ and $`𝒜=`$. (One can verify that the proof of this Theorem of Fuchs and Tabachnikov goes through for $`𝒜`$ being any Abelian group.) They later observed that since their proof of this fact is mostly local, the similar fact should be true for a big class of contact manifolds. However in fact the proof of their theorem is not completely local and is based on the existence of a well-defined Bennequin invariant for a transverse knot in $`^3`$. Unfortunately the Bennequin invariant is not well-defined unless the knot is zero homologous. And the generalization of this Theorem to the case of manifolds other than $`^3`$ meets certain difficulties that are similar to the ones we meet when we generalize the analogous Theorem of Fuchs and Tabachnikov for Legendrian knots, see 5.3. We imitate the arguments we use in 5.3 and obtain that statement a of Theorem 3.2.1 is true for any contact $`3`$-manifold with a parallelized contact structure. (Observe that in the case of transverse knots, on the contrary to the case of Legendrian knots, no extra conditions on the contact manifold appear to be needed for the statement a to be true.) Thus we get the following Theorem.
###### Theorem 3.2.2.
Let $`(M,C)`$ be a contact manifold with a parallelized contact structure, and let $`𝒯`$ be a connected component of the space of transverse curves in $`M`$. Let $``$ be the connected component of the space of framed curves that contains $`𝒯`$. Let $`V_n^𝒯`$ (resp. $`W_n^{}`$) be the group of $`𝒜`$-valued order $`n`$ invariants of transverse (resp. framed) knots from $`𝒯`$ (resp. from $``$). Then the groups $`V_n^𝒯`$ and $`W_n^{}`$ are canonically isomorphic.
This generalizes the Theorem of D. Fuchs and S. Tabachnikov saying that for the standard contact $`^3`$ and for $`𝒜=`$ the quotient groups $`V_n^𝒯/V_{n1}^𝒯`$ and $`W_n^{}/W_{n1}^{}`$ are canonically isomorphic.
The proof of this Theorem of Fuchs and Tabachnikov was based on the fact that for the $``$-valued Vassiliev invariants of framed knots in $`^3`$ there exists the universal Vassiliev invariant constructed by T. Q. T. Le and J. Murakami . (For unframed knots in $`^3`$ the construction of the universal Vassiliev invariant is the classical result of M. Kontsevich , and the invariant itself is the well-known Kontsevich integral.) The existence of the universal Vassiliev invariant is currently known only for a very limited collection of $`3`$-manifolds, and only for $`𝒜`$ being $``$, $``$ or $``$. (Andersen, Mattes, Reshetikhin proved its existence in the case where $`𝒜=`$ and $`M`$ is the total space of an $`^1`$-bundle over a compact oriented surface $`F`$ with $`F\mathrm{}`$.)
Thus the approach used in to show the isomorphism of the quotient groups is not applicable for almost all contact $`3`$-manifolds and Abelian groups $`𝒜`$, and Theorem 3.2.2 appears to be a strong generalization of the result of Fuchs and Tabachnikov.
###### Remark 3.2.3.
Let $`(M,C)`$ be a contact manifold with a parallelized contact structure, and let $``$ be a connected component of the space of framed curves in $`M`$. Theorem 3.2.2 implies that for any $`n`$ the group of $`𝒜`$-valued order $`n`$ invariants of transverse knots from a connected component of the space of transverse curves contained in $``$ does not depend on the choice of a parallelized contact structure. Hence this group can not be used to distinguish parallelized contact structures on $`M`$.
## 4. Examples of Legendrian knots that are distinguishable by finite order invariants.
In this section we construct a big class of examples when Vassiliev invariants distinguish Legendrian knots that realize isotopic framed knots and are homotopic as Legendrian curves. Theorem 3.1.1 says that in these examples the groups of Vassiliev invariants of Legendrian and of framed knots are not canonically isomorphic, and we obtain the first known examples when these groups are not canonically isomorphic.
Theorem of R. Lutz says that for an arbitrary orientable $`3`$-manifold $`M`$ every homotopy class of distributions of $`2`$-planes tangent to $`M`$ contains a contact structure. (The Theorem of Ya. Eliashberg says even more that every homotopy class of the distributions of $`2`$-planes tangent to $`M`$ contains a positive overtwisted contact structure.)
However in our constructions we will use only the Euler classes of contact bundles. For this reason we start with the following Proposition.
###### Proposition 4.0.1.
Let $`M`$ be an oriented $`3`$-manifold and let $`e`$ be an element of $`H^2(M,)`$. Then $`eH^2(M,)`$ can be realized as the Euler class of a cooriented contact structure on $`M`$ if and only if $`e=2\alpha `$, for some $`\alpha H^2(M,)`$.
For the Proof of Proposition 4.0.1 see Subsection 5.7.
### 4.1. Examples of nonisotopic Legendrian knots in $`S^1\times S^2`$ that can be distinguished by Vassiliev invariants.
Let $`C`$ be a cooriented contact structure on $`M=S^1\times S^2`$ such that the Euler class of the contact bundle is nonzero. (The Euler class takes values in $`=H^2(S^1\times S^2)`$, and Proposition 4.0.1 says that for any even $`i`$ there exists a cooriented contact structure on $`S^1\times S^2`$ with the Euler class $`i`$.)
Let $`K`$ be a knot in $`S^1\times S^2`$ that crosses exactly once one of the spheres $`t\times S^2`$. The Theorem of Chow and Rashevskii says that there exists a Legendrian knot $`K_0`$ that is $`C^0`$-small isotopic to $`K`$ as an unframed knot. Let $`K_1`$ be the Legendrian knot that is the same as $`K_0`$ everywhere except of a small piece located in a chart contactomorphic to the standard contact $`^3`$ where it is changed as it is shown in Figure 1 (see 2.0.2).
###### Theorem 4.1.1.
Legendrian knots $`K_0`$ and $`K_1`$ belong to the same component of the space of Legendrian curves and realize isotopic framed knots.
There exists a $``$-valued order one invariant $`I`$ of Legendrian knots, such that $`I(K_0)I(K_1)`$.
For the Proof of Theorem 4.1.1 see Subsection 5.8.
###### Remark 4.1.2.
Let $`K_i`$, $`i`$, be the knot that is the same as $`K_0`$ everywhere except of a small piece located in a chart contactomorphic to the standard contact $`^3`$ where it is changed in the way described by the addition of $`i`$ zigzags shown in Figure 1. The Proof of Theorem 4.1.1 implies that all $`K_i`$’s are homotopic as Legendrian curves and realize isotopic framed knots, but for all $`i_1i_2`$ Legendrian knots $`K_{i_1}`$ and $`K_{i_2}`$ are not Legendrian isotopic. The order one invariant of Legendrian knots $`I`$ constructed in the Proof of Theorem 4.1.1 has the property that $`I(K_{i_1})=I(K_{i_2})+(i_2i_1)`$. Hence this $`I`$ distinguishes all the $`K_i`$’s.
###### 4.1.3.
Examples of nonisotopic Legendrian knots with overtwisted complements that realize isotopic framed knots and are homotopic as Legendrian immersions. Let $`\mathrm{\Delta }`$ be an embedded into $`M`$ disk centered at a point $`pM`$. The Theorem of Eliashberg says that every homotopy class of distributions of $`2`$-planes tangent to $`M`$ contains an overtwisted contact structure that has $`\mathrm{\Delta }`$ as the standard overtwisted disk. In the example of Theorem 4.1.1 we can start with an overtwisted contact structure that has $`\mathrm{\Delta }`$ as an overtwisted disk and with an unframed knot $`K`$ that is far away from $`\mathrm{\Delta }`$. Then since both $`K_0`$ and $`K_1`$ were constructed using a $`C^0`$-small approximation of $`K`$, we can assume that they are also far away from $`\mathrm{\Delta }`$. And we have constructed examples of nonisotopic Legendrian knots with overtwisted complements that realize isotopic framed knots and are homotopic as Legendrian immersions. Previously such examples were unknown and the Theorem of Ya. Eliashberg and M. Fraser says that such examples are impossible if the ambient manifold is $`S^3`$.
### 4.2. Examples of nonisotopic Legendrian knots in the total spaces of $`S^1`$-bundles over nonorientable surfaces that can be distinguished by Vassiliev invariants.
###### 4.2.1.
Below we describe another big family of examples where finite order invariants distinguish Legendrian knots that realize isotopic framed knots and are homotopic as Legendrian immersions.
Let $`F`$ be a nonorientable surface that can be decomposed as a connected sum of the Klein bottle $`K`$ and a surface $`F^{}S^2`$. Let $`M`$ be an orientable manifold that admits a structure of a locally trivial $`S^1`$-fibration $`p:MF`$. (For example one can take $`M`$ to be the spherical tangent bundle $`STF`$ of $`F`$.)
Consider an $`S^1`$-fibration $`\xi :NS^1`$ induced from $`p`$ by the mapping $`S^1F`$ that corresponds to the solid loop in Figure 2. (In this Figure the enumeration of the end points of the arcs indicates which pairs of points should be identified to obtain the loop.) Since the solid loop is an orientation preserving loop in $`F`$, we get that $`N=T^2`$ (torus). Put $`\mu :N=T^2M`$ to be the natural mapping of the total space of the induced fibration $`\xi :NS^1`$ into the total space of $`p:MF`$.
A homology class in $`H_1(M,)`$ projecting to the dashed loop in Figure 2 has intersection $`1`$ with the class $`[\mu (T^2)]H_2(M,)`$ realized by $`\mu (T^2)`$. Thus there exists $`\alpha H^2(M,)`$ such that $`\alpha ([\mu (T^2)])=1`$. Proposition 4.0.1 says that for every $`r`$ the class $`2r\alpha `$ is realizable as the Euler class of a cooriented contact structure on $`M`$. Thus for every $`r`$ there exists a cooriented contact structure on $`M`$ such that the value of the Euler class of the contact bundle on $`[\mu (T^2)]`$ is equal to $`2r`$.
Let $`C`$ be a cooriented contact structure on $`M`$ such that the Euler class $`eH^2(M,)`$ of the contact bundle satisfies $`e([\mu (T^2)])=2r`$, for some nonzero $`r`$.
Let $`K`$ be an arbitrary Legendrian knot such that its projection to $`F`$ (considered as a loop) is free homotopic to the solid loop in Figure 2. Let $`K_1,K_2`$ be Legendrian knots that are the same as $`K`$ everywhere except of a chart (contactomorphic to the standard contact $`^3`$) where $`K_1`$ and $`K_2`$ are different from $`K`$ as it is described in Figure 3, see 2.0.2. (The number of cusps in Figure 3 is $`e([\mu (T^2)])=2r0`$.)
###### Theorem 4.2.2.
The knots $`K_1`$ and $`K_2`$ described above belong to the same component $``$ of the space of Legendrian curves and realize isotopic framed knots. There exists a $``$-valued order one invariant $`I`$ of Legendrian knots from $``$ such that $`I(K_1)I(K_2)`$.
For the Proof of Theorem 4.2.2 see Subsection 5.9.
###### Remark 4.2.3.
Similarly to 4.1.3 one verifies that the contact structure and the knots $`K_1`$ and $`K_2`$ in the statement of Theorem 4.2.2 can be chosen so that the restrictions of the contact structure to the complements of $`K_1`$ and of $`K_2`$ are overtwisted.
Using the ideas of the Proof of Theorem 4.2.2 one can construct many other examples of Legendrian knots that can be distinguished by Vassiliev invariants of Legendrian knots even though they realize isotopic framed knots and are homotopic as Legendrian immersions. For example as a solid loop in Figure 2 we could take any loop $`\beta `$ such that the number of double points that separate $`\beta `$ into two orientation reversing loops is odd, and the value of the Euler class of the contact bundle on $`[\mu (T^2)]H_2(M,)`$ is nonzero.
## 5. Proofs
### 5.1. Useful facts, Lemmas, and some technical definitions
###### Proposition 5.1.1.
Let $`p:XY`$ be a locally trivial $`S^1`$-fibration of an oriented manifold $`X`$ over a (not necessarily orientable) manifold $`Y`$. Let $`f\pi _1(X)`$ be the class of an oriented $`S^1`$-fiber of $`p`$, and let $`\alpha `$ be an element of $`\pi _1(X)`$. Then:
$`\alpha f=f\alpha \pi _1(X)`$, provided that $`p(\alpha )`$ is an orientation preserving loop in $`Y`$.
$`\alpha f=f^1\alpha \pi _1(X)`$, provided that $`p(\alpha )`$ is an orientation reversing loop in $`Y`$.
###### 5.1.2.
Proof of Proposition 5.1.1. If we move an oriented fiber along the loop $`\alpha X`$, then in the end it comes to itself either with the same or with the opposite orientation. It is easy to see that it comes to itself with the opposite orientation if and only if $`p(\alpha )`$ is an orientation reversing loop in $`Y`$.∎
###### Proposition 5.1.3.
Let $`FS^2,T^2\text{ (torus), }P^2,K\text{ (Klein bottle)}`$ be a surface (not necessarily compact or orientable), and let $`G`$ be a nontrivial commutative subgroup of $`\pi _1(F)`$. Then $`G`$ is infinite cyclic.
###### 5.1.4.
Proof of Proposition 5.1.3. It is well-known that any closed $`F`$, other than $`S^2,T^2,P^2,K`$ admits a hyperbolic metric. (It is induced from the universal covering of $`F`$ by the hyperbolic plane.) The Theorem of A. Preissman (see pp. 258-265) says that if $`M`$ is a closed Riemannian manifold of negative sectional curvature, then any nontrivial Abelian subgroup $`G<\pi _1(M)`$ is isomorphic to $``$. Thus if $`FS^2,T^2,P^2,K`$ is closed, then any nontrivial commutative $`G<\pi _1(F)`$ is infinite cyclic. If $`F`$ is not closed, then the statement of the Proposition is also true because in this case $`F`$ is homotopy equivalent to a bouquet of circles. ∎
###### Proposition 5.1.5.
Let $`FS^2,P^2,T^2,K`$ (Klein bottle) be a surface not necessarily closed or orientable. Let $`M`$ be an orientable $`3`$-manifold, and let $`p:MF`$ be a locally trivial $`S^1`$-fibration. Let $`f\pi _1(M)`$ be the class of an oriented $`S^1`$-fiber of $`p`$, and let $`\alpha \pi _1(M)`$ be an element with $`p_{}(\alpha )1\pi _1(F)`$. Let $`\beta `$ be an element of the centralizer $`Z(\alpha )<\pi _1(M)`$ of $`\alpha `$. Then there exist $`i,j`$ and nonzero $`n`$ such that $`\beta ^n=\alpha ^if^j`$.
###### 5.1.6.
Proof of Proposition 5.1.5. Since $`\alpha `$ and $`\beta `$ commute in $`\pi _1(M)`$ we get that $`p_{}(\alpha )`$ and $`p_{}(\beta )`$ commute in $`\pi _1(F)`$. Proposition 5.1.3 and the fact that $`p_{}(\alpha )1\pi _1(F)`$ imply that there exist $`g\pi _1(M)`$ with $`p_{}(g)1\pi _1(F)`$, $`i`$, and nonzero $`n`$ such that $`p_{}(g)^n=p_{}(\alpha )`$ and $`p_{}(g)^i=p_{}(\beta )`$.
Hence (see Proposition 5.1.1), $`\alpha =g^nf^k\text{ and }\beta =g^if^l,\text{ for some }k,l`$. Using 5.1.1 we get that $`\beta ^n=\alpha ^if^j,\text{ for some }j`$. Since $`n`$ was initially chosen to be nonzero, we get the statement of the Proposition. ∎
###### 5.1.7.
An important homomorphism. Let $`X`$ be a manifold, let $`\mathrm{\Omega }X`$ be the space of free loops in $`X`$, and let $`\omega \mathrm{\Omega }X`$ be a loop. An element $`\alpha \pi _1(\mathrm{\Omega }X,\omega )`$ is realizable by a mapping $`\mu ^\alpha :T^2=S^1\times S^1X`$ with $`\mu ^\alpha |_{t\times S^1}=\alpha (t)`$. Let $`t(\alpha )=\mu ^\alpha |_{S^1\times 1}\pi _1(X,\omega (1))`$ be the element corresponding to the trace of the point $`1S^1`$ under the homotopy of $`\omega `$ described by $`\alpha `$. Let $`t:\pi _1(\mathrm{\Omega }X,\omega )\pi _1(X,\omega (1))`$ be the homomorphism that maps $`\alpha \pi _1(\mathrm{\Omega }X,\omega )`$ to $`t(\alpha )\pi _1(X,\omega (1))`$.
Since the $`2`$-cell of $`T^2`$ is glued to the $`1`$-skeleton along the commutation relation of the meridian and of the longitude of $`T^2`$, we get that $`t:\pi _1(\mathrm{\Omega }X,\omega )\pi _1(X,\omega (1))`$ is a surjective homomorphism of $`\pi _1(\mathrm{\Omega }X,\omega )`$ onto the centralizer $`Z(\omega )`$ of $`\omega \pi _1(X,\omega (1))`$.
If $`t(\alpha )=t(\beta )\pi _1(X,\omega (1))`$ for $`\alpha ,\beta \pi _1(\mathrm{\Omega }X,\omega )`$, then the mappings $`\mu ^\alpha `$ and $`\mu ^\beta `$ of $`T^2`$ corresponding to these loops can be deformed to be identical on the $`1`$-skeleton of $`T^2`$. Clearly the obstruction for $`\mu ^\alpha `$ and $`\mu ^\beta `$ to be homotopic as mappings of $`T^2`$ (with the mapping of the $`1`$-skeleton of $`T^2`$ fixed under homotopy) is an element of $`\pi _2(X)`$ obtained by gluing together the boundaries of the $`2`$-cells of the two tori. In particular we get the Proposition of V. L. Hansen saying that $`t:\pi _1(\mathrm{\Omega }X,\omega )Z(\omega )<\pi _1(X,\omega (1))`$ is an isomorphism, provided that $`\pi _2(X)=0`$.
###### 5.1.8.
$`h`$-principle for curves in $`M`$. For a $`3`$-dimensional manifold $`M`$ we put $`STM`$ to be the manifold obtained by the fiberwise spherization of the tangent bundle of $`M`$, and we put $`\mathrm{pr}^{}:STMM`$ to be the corresponding locally trivial $`S^2`$-fibration. The $`h`$-principle (that can be found in ) says that the space of curves in $`M`$ is weak homotopy equivalent to $`\mathrm{\Omega }STM`$ (the space of free loops in $`STM`$). The weak homotopy equivalence is given by mapping a curve $`K`$ to a loop $`\stackrel{}{K}\mathrm{\Omega }STM`$ that sends a point $`tS^1`$ to the point of $`STM`$ corresponding to the direction of the velocity vector of $`K`$ at $`K(t)`$.
###### Definition 5.1.9 (of $`m(K_1,K_2)`$ and of $`K^0,K^{\pm 1},K^{\pm 2}\mathrm{}`$).
Let $`K_1`$ and $`K_2`$ be two framed knots that coincide pointwise as embeddings of $`S^1`$. Then there is an integer obstruction $`m(K_1,K_2)`$ for them to be isotopic as framed knots with the embeddings of $`S^1`$ fixed under the isotopy. This obstruction is calculated as follows. Let $`K_1^{}`$ be the knot obtained by shifting $`K_1`$ along the framing and reversing the orientation on the shifted copy. Together $`K_1`$ and $`K_1^{}`$ bound a thin strip. We put $`m(K_1,K_2)`$ to be the intersection number of the strip with a very small shift of $`K_2`$ along its framing.
For a framed knot $`K^0`$ we denote by $`K^i`$, $`i`$, the isotopy class of a framed knot that coincides with $`K^0`$ as an embedding of $`S^1`$ and has $`m(K^0,K^i)=i`$.
For two singular framed knots $`K_{1s}`$ and $`K_{2s}`$ with $`n`$ transverse double points that coincide pointwise as immersions of $`S^1`$, we put $`m(K_{1s},K_{2s})`$ to be the value of $`m`$ on the nonsingular framed knots $`K_1`$ and $`K_2`$ that coincide pointwise as embeddings of $`S^1`$ and are obtained from $`K_{1s}`$ and $`K_{2s}`$ by resolving each pair of the corresponding double points of $`K_{1s}`$ and of $`K_{2s}`$ in the same way. (The value $`m(K_{1s},K_{2s})`$ does not depend on the resolution as soon as the corresponding double points of the two knots are resolved in exactly the same way.) As before $`m(K_{1s},K_{2s})`$ is the integer valued obstruction for $`K_{1s}`$ and $`K_{2s}`$ to be isotopic as singular framed knots with the immersion of $`S^1`$ corresponding to the two knots fixed under isotopy.
For a singular framed knot $`K_s^0`$ with $`n`$ transverse double points we denote by $`K_s^i`$, $`i`$, the isotopy class of a singular framed knot with $`n`$ transverse double points that coincides with $`K_s^0`$ as an immersion of $`S^1`$ and has $`m(K_s^0,K_s^i)=i`$.
###### Proposition 5.1.10.
Let $`K_1`$ and $`K_2`$ be framed knots (resp. singular framed knots with $`n`$ transverse double points) that coincide pointwise as embeddings (resp. immersions) of $`S^1`$. Then $`K_1`$ and $`K_2`$ are homotopic as framed knots (resp. singular framed knots with $`n`$ transverse double points) if and only if $`m(K_1,K_2)`$ is even.
###### 5.1.11.
Proof of Proposition 5.1.10. Clearly if $`m(K_1,K_2)`$ is even, then $`K_1`$ and $`K_2`$ are framed homotopic. (We can change the obstruction by two by creating a small kink and passing through a double point at its vertex.)
Every oriented $`3`$-dimensional manifold $`M`$ is parallelizable, and hence it admits a $`\mathrm{spin}`$-structure. A framed curve $`K`$ in $`M`$ represents a loop in the principal $`SO(3)`$-bundle of $`TM`$. (The $`3`$-frame corresponding to a point of $`K`$ is the velocity vector, the framing vector, and the unique third vector of unit length such that the $`3`$-frame defines the positive orientation of $`M`$.) One observes that the values of the $`\mathrm{spin}`$-structure on the loops in the principal $`SO(3)`$-bundle of $`TM`$ realized by $`K_1`$ and $`K_2`$ are different provided that $`m(K_1,K_2)`$ is odd. But these values do not change under homotopy of framed curves. Hence if $`m(K_1,K_2)`$ is odd, then $`K_1`$ and $`K_2`$ are not framed homotopic. ∎
###### Definition 5.1.12 (of the number of framings of a knot).
Using the self-linking invariant of framed knots one can easily show that if $`K_1`$ and $`K_2`$ in 5.1.9 are pointwise coinciding zero-homologous framed knots and $`m(K_1,K_2)0`$, then $`K_1`$ is not isotopic to $`K_2`$ in the category of framed knots. However for knots that are not zero-homologous this is not generally true, see 5.8.1. For this reason we introduce the following definitions.
If for an unframed knot $`K`$ there exist isotopic framed knots $`K_1`$ and $`K_2`$ that coincide with $`K`$ pointwise and have $`m(K_1,K_2)0`$, then we say that $`K`$ admits finitely many framings. For $`K`$ that admits finitely many framings we put the number of framings $`m_K`$ of $`K`$ to be the minimal positive integer $`l`$ such that there exist isotopic framed knots $`K_1`$ and $`K_2`$ that coincide with $`K`$ pointwise and have $`m(K_1,K_2)=l`$. One can easily show that if $`K`$ admits finitely many framings, then there are exactly $`m_K`$ isotopy classes of framed knots realizing the isotopy class of the unframed knot $`K`$. Proposition 5.1.10 implies that $`m_K`$ is even.
In a similar way we introduce the notion of the number of framings for unframed singular knots with $`n`$ double points.
###### Proposition 5.1.13.
Let $`(M,C)`$ be a contact $`3`$-manifold with a cooriented contact structure, let $``$ be a connected component of the space of framed curves, and let $``$ be a connected component of the space of Legendrian curves in $`(M,C)`$.
Let $`K`$ be an unframed knot obtained by forgetting the framing on a knot from $``$. Then there exists a Legendrian knot from $``$ realizing the isotopy class of $`K`$.
If $`K^0`$ is an isotopy class of framed knots in $``$ that is realizable by a Legendrian knot from $``$, then the isotopy class of $`K^2`$ (see 5.1.9) is also realizable by a Legendrian knot from $``$.
Let $`K_s`$ be an unframed singular knot with $`n`$ double points obtained by forgetting the framing on a singular knot from $``$. Then there exists a singular Legendrian knot from $``$ realizing the isotopy class of $`K_s`$.
If $`K_s^0`$ is an isotopy class of singular framed knots in $``$ that is realizable by a singular Legendrian knot from $``$, then the isotopy class of $`K_s^2`$ is also realizable by a singular Legendrian knot from $``$.
###### 5.1.14.
Proof of statement a of Proposition 5.1.13. Let $`CM`$ be the fiberwise spherization of the $`2`$-dimensional contact vector bundle, and let $`\mathrm{pr}:CMM`$ be the corresponding locally trivial $`S^1`$-fibration. We denote by $`f\pi _1(CM)`$ the class of an oriented $`S^1`$-fiber of $`\mathrm{pr}`$. For a Legendrian curve $`K_l:S^1M`$ denote by $`\stackrel{}{K}_l`$ the loop in $`CM`$ obtained by mapping a point $`tS^1`$ to the point of $`CM`$ corresponding to the direction of the velocity vector of $`K_l`$ at $`K_l(t)`$.
The $`h`$-principle 2.0.1 says that Legendrian curves $`K_1`$ and $`K_2`$ in $`M`$ belong to the same component of the space of Legendrian curves in $`M`$ if and only if $`\stackrel{}{K}_1`$ and $`\stackrel{}{K}_2`$ are free homotopic loops in $`CM`$.
W. L. Chow and P. K. Rashevskii showed that every unframed knot $`K`$ is isotopic to a Legendrian knot $`K_l`$ (and this isotopy can be made $`C^0`$-small). Deforming $`K`$ we can assume (see 5.1.8) that: 1: $`K`$ and $`K_l`$ coincide in the neighborhood of $`1S^1`$, 2: $`K`$ and $`K_l`$ realize the same element $`[K]\pi _1(M,K_l(1))`$, and 3: that liftings to $`CM`$ of Legendrian curves from $``$ are free homotopic to a loop $`\alpha `$ in $`CM`$ such that $`\alpha (1)=\stackrel{}{K}_l(1)`$ and $`\mathrm{pr}(\alpha )=[K]\pi _1(M,K_l(1))`$.
Proposition 5.1.1 says that $`f`$ is in the center of $`\pi _1(CM,\stackrel{}{K}_l(1))`$, since the contact structure is cooriented and hence oriented. Then $`\stackrel{}{K}_l=\alpha f^i\pi _1(CM,\stackrel{}{K}_l(1))`$, for some $`i`$.
Take a chart of $`M`$ (that is contactomorphic to the standard contact $`^3`$) containing a piece of the Legendrian knot. From the formula for the Maslov number deduced in it is easy to see that the modifications of the Legendrian knot corresponding to the insertions of two cusps shown in Figure 4 (see 2.0.2) induce multiplication by $`f^{\pm 1}`$ of the lifting of $`K_l`$ to an element of $`\pi _1(CM,\stackrel{}{K}_l(1))`$. (Here the sign depends on the choice of an orientation of the fiber used to define $`f`$.) Performing this operation sufficiently many times we obtain the Legendrian knot from $``$ realizing the isotopy class of the unframed knot $`K`$.
One easily modifies the arguments above to obtain the proof of statement c of Proposition 5.1.13.
Proof of statement b of Proposition 5.1.13. Take a chart of $`M`$ (that is contactomorphic to the standard contact $`^3`$) containing a piece of the knot $`K^0`$ and perform the homotopy in $``$ shown in Figure 5, see 2.0.2. (Observe that a self-tangency point of the projection of a Legendrian curve in $`^3`$ to the $`(x,z)`$-plane corresponds to a double point of the Legendrian curve.) Straightforward verification (cf. the formula for the Bennequin invariant deduced in ) shows that the Legendrian knot we obtain in the end of the homotopy realizes $`K^2`$.
One easily modifies these arguments to obtain the proof of statement d of Proposition 5.1.13. This finishes the proof of Proposition 5.1.13.∎
### 5.2. Proof of Theorem 3.1.1
The fact that statement b of Theorem 3.1.1 implies statement a is clear. Thus we have to show that statement a implies statement b. This is done by showing that there exists a homomorphism $`\psi :V_n^{}W_n^{}`$ such that $`\varphi \psi =\mathrm{id}_{V_n^{}}`$ and $`\psi \varphi =\mathrm{id}_{W_n^{}}`$.
Let $`xV_n^{}`$ be an invariant. In order to construct $`\psi (x)W_n^{}`$ we have to specify the value of $`\psi (x)`$ on every framed knot $`K`$.
###### 5.2.1.
Definition of $`\psi (x)`$. If the isotopy class of the knot $`K`$ is realizable by a Legendrian knot $`K_l`$, then put $`\psi (x)(K)=x(K_l)`$. The value $`\psi (x)(K)`$ is well-defined because if $`K_l^{}`$ is another knot realizing $`K`$, then $`x(K_l)=x(K_l^{})`$ by statement a of Theorem 3.1.1.
Let $`𝒞`$ be the component of the space of unframed curves that corresponds to forgetting framings on the curves from $``$. Propositions 5.1.13 and 5.1.10 imply that if an unframed knot $`K_u𝒞`$ admits finitely many framings (see 5.1.12), then all the isotopy classes of framed knots from $``$ realizing the isotopy class of the unframed knot $`K_u`$ are realizable by Legendrian knots from $``$. Thus we have defined the value of $`\psi (x)`$ on all the framed knots from $``$ that realize unframed knots admitting finitely many framings.
If $`K_u𝒞`$ admits infinitely many framings, then either 1) all the isotopy classes of framed knots from $``$ realizing the isotopy class of $`K_u`$ are realizable by Legendrian knots from $``$ or 2) there exists a knot $`K^0`$ realizing the isotopy class of $`K_u`$ such that $`K^0`$ is realizable by a Legendrian knot from $``$ and $`K^{+2}`$ (see 5.1.9) is not realizable by a Legendrian knot from $``$. (In this case $`K^{+4},K^{+6}`$ etc. also are not realizable by Legendrian knots from $``$, see 5.1.13.) In the case 1) the value of $`\psi (x)`$ is already defined on all the framed knots from $``$ realizing $`K_u`$. In the case 2) put
(1)
$$\psi (x)(K^{+2})=\underset{i=1}{\overset{n+1}{}}\left((1)^{i+1}\frac{(n+1)!}{i!(n+1i)!}\psi (x)(K^{+22i})\right).$$
(Proposition 5.1.13 implies that the sum on the right hand side is well-defined.) Similarly put
$$\begin{array}{c}\psi (x)(K^{+4})=_{i=1}^{n+1}\left((1)^{i+1}\frac{(n+1)!}{i!(n+1i)!}\psi (x)(K^{+42i})\right),\hfill \\ \psi (x)(K^{+6})=_{i=1}^{n+1}\left((1)^{i+1}\frac{(n+1)!}{i!(n+1i)!}\psi (x)(K^{+62i})\right)\text{ etc.}\hfill \end{array}$$
Now we have defined $`\psi (x)`$ on all the framed knots (from $``$) realizing $`K_u`$. Doing this for all $`K_u`$ for which case 2) holds we define the value of $`\psi (x)`$ on all the knots from $``$.
Below we show that $`\psi (x)`$ is an order $`n`$ invariant of framed knots from $``$. We start by proving the following Proposition.
###### Proposition 5.2.2.
Let $`K^0`$ be a framed knot from $``$, then $`\psi (x)`$ defined as above satisfies identity (1).
###### 5.2.3.
Proof of Proposition 5.2.2. If $`K^{+2}`$ is not realizable by a Legendrian knot from $``$, then the statement of the proposition follows from the formula we used to define $`\psi (x)(K^{+2})`$.
If $`K^{+2}`$ is realizable by a Legendrian knot $`K_l^0`$, then consider a singular Legendrian knot $`K_{ls}`$ with $`(n+1)`$ double points that are vertices of $`(n+1)`$ small kinks such that we get $`K_l`$ if we resolve all the double points positively staying in the class of the Legendrian knots. (To create $`K_{ls}`$ we perform the first half of the homotopy shown in Figure 5 in $`n+1`$ places on $`K_l^0`$.)
Let $`\mathrm{\Sigma }`$ be the set of the $`2^{n+1}`$ possible resolutions of the double points of $`K_{ls}`$. For $`\sigma \mathrm{\Sigma }`$ put $`\mathrm{sign}(\sigma )`$ to be the sign of the resolution, and put $`K_{ls}^\sigma `$ to be the nonsingular Legendrian knot obtained via the resolution $`\sigma `$. Since $`x`$ is an order $`n`$ invariant of Legendrian knots we get that
(2)
$$0=\underset{\sigma \mathrm{\Sigma }}{}\left(\mathrm{sign}(\sigma )x(K_{ls}^\sigma )\right)=\psi (x)(K_l^0)+\underset{i=1}{\overset{n+1}{}}(1)^i\frac{(n+1)!}{i!(n+1i)!}\psi (x)(K_l^{2i}).$$
(Observe that if we resolve $`i`$ double points of $`K_{ls}`$ negatively, then we get the isotopy class of $`K_l^{2i}`$.) This finishes the proof of the Proposition. ∎
###### 5.2.4.
Let $`K_s`$ be a singular framed knot with $`(n+1)`$ double points. Let $`\mathrm{\Sigma }`$ be the set of the $`2^{n+1}`$ possible resolutions of the double points of $`K_s`$. For $`\sigma \mathrm{\Sigma }`$ put $`\mathrm{sign}(\sigma )`$ to be the sign of the resolution, and put $`K_s^\sigma `$ to be the isotopy class of the knot obtained via the resolution $`\sigma `$.
In order to prove that $`\psi (x)`$ is an order $`n`$ invariant of framed knots from $``$, we have to show that
(3)
$$0=\underset{\sigma \mathrm{\Sigma }}{}\left(\mathrm{sign}(\sigma )\psi (x)(K_s^\sigma )\right),$$
for every $`K_s`$.
First we observe that the fact whether identity (3) holds or not depends only on the isotopy class of the singular knot $`K_s`$ with $`(n+1)`$ double points.
If the isotopy class of $`K_s`$ is realizable by a singular Legendrian knot from $``$, then identity (3) holds for $`K_s`$, since $`x`$ is an order $`n`$ invariant of Legendrian knots (and the value of $`\psi (x)`$ on a framed knot $`K`$ realizable by a Legendrian knot $`K_l`$ was put to be $`x(K_l)`$).
Proposition 5.1.13 says that the isotopy class of the singular unframed knot $`K_{us}`$ obtained by forgetting the framing on $`K_s`$ is realizable by a singular Legendrian knot from $``$.
If $`K_{us}`$ admits finitely many framings, then all the isotopy classes of singular framed knots from $``$ realizing the isotopy class of $`K_{us}`$ are realizable by singular Legendrian knots from $``$, and we get that identity (3) holds for $`K_s`$.
If $`K_{us}`$ admits infinitely many framings and all the isotopy classes of singular framed knots from $``$ realizing $`K_{us}`$ are realizable by singular Legendrian knots from $``$, then (3) automatically holds for $`K_s`$. If $`K_{us}`$ admits infinitely many framings but not all the isotopy classes of singular framed knots from $``$ realizing $`K_{us}`$ are realizable by singular Legendrian knots from $``$. Then put $`K_{us}^0`$ to be the framed knot realizing $`K_{us}`$ that is realizable by a singular Legendrian knot from $``$ and such that $`K_{us}^{2i}`$, $`i>0`$, are not realizable by singular Legendrian knots from $``$.
Proposition 5.1.13 says that $`K_{us}^{2i}`$, $`i>0`$, are realizable by singular Legendrian knots from $``$ and hence identity (3) holds for $`K_{us}^{2i}`$, $`i0`$. Using Proposition 5.2.2 and the fact that identity (3) holds for $`K_{us}^{2i}`$, $`i0`$, we show that (3) holds for $`K_{us}^{+2}`$. Namely,
(4)
$$\begin{array}{c}\underset{\sigma \mathrm{\Sigma }}{}\mathrm{sign}(\sigma )\psi (x)(K_{us}^{+2}{}_{}{}^{\sigma })\hfill \\ \hfill =\underset{\sigma \mathrm{\Sigma }}{}\left(\mathrm{sign}(\sigma )\underset{i=1}{\overset{n+1}{}}(1)^{i+1}\frac{(n+1)!}{i!(n+1i)!}\psi (x)(K_{us}^{(+22i)}{}_{}{}^{\sigma })\right)\\ \hfill =\underset{i=1}{\overset{n+1}{}}\left((1)^{i+1}\frac{(n+1)!}{i!(n+1i)!}\times \left(\underset{\sigma \mathrm{\Sigma }}{}\mathrm{sign}(\sigma )\psi (x)(K_{us}^{(+22i)}{}_{}{}^{\sigma })\right)\right)\\ \hfill =\underset{i=1}{\overset{n+1}{}}\left((1)^{i+1}\frac{(n+1)!}{i!(n+1i)!}\times \left(0\right)\right)=0.\end{array}$$
Similarly we show that (3) holds for $`K_{us}^{+4},K_{us}^{+6},\text{ etc}\mathrm{}`$.
###### 5.2.5.
Clearly $`\varphi \psi =\mathrm{id}_{V_n^{}}`$.
Considering the values of $`yW_n^{}`$ on the $`2^{n+1}`$ possible resolutions of a singular framed knot with $`n+1`$ singular small kinks we get that $`y`$ should satisfy identity (1). Hence $`\psi \varphi =\mathrm{id}_{W_n^{}}`$ and this finishes the proof of Theorem 3.1.1. ∎
### 5.3. The reasons for the condition $`()`$ to appear
The proof of the Theorem of Fuchs and Tabachnikov that says that statement $`𝐚`$ of Theorem 3.1.1 is true for all the connected components of the space of Legendrian curves when the ambient contact manifold is the standard contact $`^3`$ is based on the following three observations:
There are two types of cusps arising under the projection of the part of a Legendrian knot that is contained in a Darboux chart to the $`(x,z)`$-plane (see 2.0.2). They are formed by cusps for which the branch of the projection of the knot going away from the cusp is locally located respectively above or below the tangent line at the cusp point, see Figure 4. For a Legendrian knot $`K`$ and $`i,j`$ we denote by $`K^{i,j}`$ the Legendrian knot obtained from $`K`$ by the modification corresponding to an addition of $`i`$ cusp pairs of the first type and $`j`$ cusp pairs of the second type to the projection of the part of $`K`$ located in a Darboux chart.
Let $`K_1`$ and $`K_2`$ be Legendrian knots in the standard contact $`^3`$ that realize isotopic unframed knots. Then for any $`n_1`$ and $`n_2`$ large enough there exist $`n_3,n_4`$ such that the Legendrian knot $`K_1^{n_1,n_2}`$ is Legendrian isotopic to $`K_2^{n_3,n_4}`$.
If there exists $`n`$ such that Legendrian knots $`K_1^{n,n}`$ and $`K_2^{n,n}`$ are Legendrian isotopic, then every Vassiliev invariant of Legendrian knots takes equal values on $`K_1`$ and on $`K_2`$.
The number $`n`$ from the previous observation exists if the ambient contact manifold is $`^3`$ and the Legendrian knots $`K_1`$ and $`K_2`$ belong to the same component of the space of Legendrian curves and realize isotopic framed knots.
The first two observation are true for any contact $`3`$-manifold (since the proof of the corresponding facts is local). But the number $`n`$ from the statement of the third observation does not exist in general. In the case of the ambient manifold being $`^3`$ Fuchs and Tabachnikov showed the existence of such $`n`$ using the explicit calculation involving the Maslov classes and Bennequin invariants of Legendrian knots. However in order for the Bennequin invariant to be well-defined the knots have to be zero-homologous, and in order for the Maslov class to be well-defined the knots have to be zero-homologous or the contact structure has to be parallelizable.
Below we show that such $`n`$ exists for any $`K_1`$ and $`K_2`$ that realize isotopic framed knots and belong to the same component of the space of Legendrian curves, provided that the connected component $``$ of the space of framed curves that contains $`K_1`$ and $`K_2`$ satisfies condition $`()`$. (We assume that the contact structure on $`M`$ is cooriented.)
Let $`K_1`$ and $`K_2`$ be Legendrian knots as above, and let $`n_1,n_2,n_3,n_4`$ be such that $`K_1^{n_1,n_2}`$ and $`K_2^{n_3,n_4}`$ are Legendrian isotopic.
We start by showing that $`n_1,n_2,n_3,n_4`$ can be chosen so that $`n_1+n_2=n_3+n_4`$, and that if $``$ satisfies condition $`()`$, then $`n_1n_2=n_3n_4`$.
###### 5.3.1.
Proof of the fact that $`n_1,n_2,n_3,n_4`$ can be chosen so that $`n_1+n_2=n_3+n_4`$. Let $`\mu :S^1\times [0,1]M`$ be the isotopy changing $`K_1`$ to $`K_2`$ in the category of framed knots. Analyzing the proof of Fuchs and Tabachnikov one verifies that for $`n_1,n_2`$ large enough the Legendrian isotopy $`\overline{\mu }`$ changing $`K_1^{n_1,n_2}`$ to $`K_2^{n_3,n_4}`$ can be chosen so that for every $`t[0,1]`$ the Legendrian knot $`\overline{\mu }_t:S^1\times tM`$ is contained in a thin tubular neighborhood $`T_t`$ of $`\mu _t:S^1\times tM`$ and is isotopic (as an unframed knot) to $`\mu _t`$ inside $`T_t`$.
For two framed knots $`\mu _t`$ and $`\overline{\mu }_t`$ realizing unframed knots that are isotopic inside $`T_t`$ there is a well-defined $``$-valued obstruction to be isotopic inside $`T_t`$ in the category of framed knots. This obstruction is the difference of the self-linking numbers of the inclusions of $`\mu _t`$ and $`\overline{\mu }_t`$ into $`^3`$ induced by an identification of $`T_t`$ with the standard solid torus in $`^3`$. (One verifies that for $`\mu _t`$ and $`\overline{\mu }_t`$ that are isotopic as unframed knots inside $`T_t`$ this difference does not depend on the choice of the identification of $`T_t`$ with the standard solid torus in $`^3`$.)
From the formula for the Bennequin invariant stated in one gets that the value of the obstruction for $`K_1^{n_1,n_2}`$ to be isotopic as a framed knot to $`K_1`$ inside $`T_0`$ is equal to $`n_1+n_2`$. Similarly the value of the obstruction for $`K_2^{n_3,n_4}`$ to be isotopic as a framed knot to $`K_2`$ inside $`T_1`$ is equal to $`n_3+n_4`$. Clearly the value of the obstruction for $`\overline{\mu }_t`$ to be isotopic to $`\mu _t`$ inside $`T_t`$ does not depend on $`t`$ (for the isotopy $`\mu `$ changing $`K_1`$ to $`K_2`$ in the category of framed knots), and we get that $`n_1+n_2=n_3+n_4`$.
###### 5.3.2.
Proof of the fact that if $``$ satisfies condition $`()`$, then $`n_1n_2=n_3n_4`$. Let $`f\pi _1(CM)`$ be the class of an oriented $`S^1`$-fiber of $`\mathrm{pr}:CMM`$. From the $`h`$-principles for Legendrian and for unframed curves (see 2.0.1 and 5.1.8) one obtains that every component of the space of Legendrian curves contained in $``$ corresponds to the conjugacy class of $`\stackrel{}{K}_1f^l\pi _1(CM)`$, for some $`l`$. (Connected components of the space of free loops in $`CM`$ are naturally identified with the conjugacy classes of the elements of $`\pi _1(CM)`$.)
Using Proposition 5.1.1 one verifies that if $``$ satisfies condition $`()`$, then for every nonzero $`l`$ the elements $`\stackrel{}{K}_1`$ and $`\stackrel{}{K}_1f^l`$ are not conjugate in $`\pi _1(CM)`$.
From the formula for the Maslov number deduced in and the $`h`$-principle for Legendrian curves one gets that $`K_1^{n_1,n_2}`$ is contained in the component of the space of Legendrian curves that corresponds to the conjugacy class of $`\stackrel{}{K}_1f^{n_1n_2}\pi _1(CM)`$. Using the fact that $`\stackrel{}{K}_1`$ and $`\stackrel{}{K}_2`$ are conjugate in $`\pi _1(CM)`$ (since $`K_1`$ and $`K_2`$ are Legendrian homotopic) and the fact that since the contact structure is cooriented $`f`$ is in the center of $`\pi _1(CM)`$ (see 5.1.1), we get that $`K_2^{n_3,n_4}`$ is contained in the component that corresponds to the conjugacy class of $`\stackrel{}{K}_1f^{n_3n_4}\pi _1(CM)`$. Since $`K_1^{n_1,n_2}`$ and $`K_2^{n_3,n_4}`$ are Legendrian isotopic (and hence Legendrian homotopic) we get that $`\stackrel{}{K}_1f^{n_1n_2}`$ and $`\stackrel{}{K}_1f^{n_3n_4}`$ are conjugate in $`\pi _1(CM)`$, and using 5.1.1 we get that $`\stackrel{}{K}_1`$ is conjugate to $`\stackrel{}{K}_1f^{(n_1n_2)(n_3n_4)}`$. But since $``$ satisfies condition $`()`$ we have $`(n_1n_2)(n_3n_4)=0`$, and hence $`n_1n_2=n_3n_4`$.
From the identities $`n_1+n_2=n_3+n_4`$ and $`n_1n_2=n_3n_4`$ one gets that $`n_1=n_3`$ and $`n_2=n_4`$. Assume that $`n_1n_2`$. (The case where $`n_2>n_1`$ is treated similarly.) Put $`k=n_1n_2`$. It is easy to show that since $`K_1^{n_1,n_2}`$ and $`K_2^{n_3,n_4}`$ are Legendrian isotopic, then $`K_1^{n_1,n_2+k}`$ and $`K_2^{n_3,n_4+k}`$ are also Legendrian isotopic. (Basically one can keep the $`k`$ extra cusp pairs close together on a small piece of the projection of the part of the knot contained in a Darboux chart during the whole isotopy process.) But $`K_1^{n_1,n_2+k}`$ and $`K_2^{n_3,n_4+k}`$ are obtained from $`K_1`$ and $`K_2`$ by the modification corresponding to the addition of $`n_1=n_2+k=n_3=n_4+k`$ pairs of cusps of each of the two types, and we can take $`n`$ from the observation 2 to be $`n_1=n_2+k=n_3=n_4+k`$.
This shows that $`K_1`$ and $`K_2`$ can not be distinguished by the Vassiliev invariants of Legendrian knots provided that $``$ satisfies condition $`()`$, and that $`K_1`$ and $`K_2`$ realize isotopic framed knots and are homotopic as Legendrian immersions. Hence statement a of Theorem 3.1.1 is true provided that $``$ satisfies condition $`()`$.
### 5.4. Proof of Proposition 3.1.4.
The $`h`$-principle for curves 5.1.8 says that the set $`𝒞`$ of the connected components of the space of curves in $`M`$ is naturally identified with the set of the connected components of the space of free loops in the spherical tangent bundle $`STM`$ of $`M`$. Hence it is also naturally identified with the set of conjugacy classes of the elements of $`\pi _1(STM)`$. (From the long homotopy sequence of the fibration $`\mathrm{pr}^{}:STMM`$ we see that it is also naturally identified with the set of conjugacy classes of the elements of $`\pi _1(M)`$.) Choose a $`\mathrm{spin}`$-structure on $`M`$. It is easy to see (cf. 5.1.10 and 5.1.11) that the set $`𝒞_{}`$ of the connected components of the space of framed curves in $`M`$ is identified with the product $`_2\times 𝒞`$. Here the $`_2`$-factor is the value of the $`\mathrm{spin}`$-structure on the loop in the principal $`SO(3)`$-bundle of $`TM`$ that corresponds to a framed curve from the connected component, see 5.1.11. (This value does not depend on the choice of the framed curve in the component.)
The $`h`$-principle for the Legendrian curves says that the set of the connected components of the space of Legendrian curves is naturally identified with the set of homotopy classes of free loops in $`CM`$ (the spherical contact bundle of $`M`$). Hence it is also naturally identified with the set of conjugacy classes of the elements of $`\pi _1(CM)`$. Since every contact manifold is oriented and the contact structure was assumed to be cooriented, we get that the planes of the contact structure are naturally oriented. This orientation induces the orientation of the $`S^1`$-fibers of $`\mathrm{pr}:CMM`$. Put $`f\pi _1(CM)`$ to be the class of the oriented $`S^1`$-fiber of $`\mathrm{pr}:CMM`$.
The Theorem of Chow and Rashevskii says that every connected component of the space of curves contains a Legendrian curve. Straightforward verification shows that the insertion of the zig-zag into the Legendrian curve $`K`$ (see Figure 4) changes the value of the $`\mathrm{spin}`$-structure on the corresponding framed curve. It is easy to verify (see ) that the two connected components of the space of Legendrian curves that contain respectively $`K`$ and $`K`$ with the extra zigzag correspond to the conjugacy classes of $`\stackrel{}{K}`$ and of $`\stackrel{}{K}f`$ (or of $`\stackrel{}{K}f^1`$) in $`\pi _1(CM)`$. (We obtain $`\stackrel{}{K}f`$ or $`\stackrel{}{K}f^1`$ depending on which of the two possible zig-zags we insert.)
Let $``$ be a connected component of the space of Legendrian curves in $`(M,C)`$ that corresponds to the conjugacy class of $`\stackrel{}{K}\pi _1(CM)`$. Then every connected component $`^{}`$ of the space of Legendrian curves corresponds to the conjugacy class of $`\stackrel{}{K}f^{2n}\pi _1(CM)`$, for some $`n`$.
Hence $``$ satisfies condition $`()`$ if and only if for every $`n0`$ the elements $`\stackrel{}{K}`$ and $`\stackrel{}{K}f^{2n}`$ are not conjugate in $`\pi _1(CM)`$.
Assume that $``$ does not satisfy condition $`()`$, then there exists a nonzero $`n`$ and $`\beta \pi _1(CM)`$ such that
(5)
$$\beta \stackrel{}{K}\beta ^1=\stackrel{}{K}f^{2n}\pi _1(CM,\stackrel{}{K}(1)).$$
This implies that $`\mathrm{pr}_{}(\beta )`$ and $`\mathrm{pr}_{}(\stackrel{}{K})`$ commute in $`\pi _1(M,K(1))`$. The commutation relation gives a mapping $`\mu :T^2M`$ of the two-torus $`T^2=S^1\times S^1`$ such that $`\mu |_{(1\times S^1)}=K`$ and $`\mu |_{(S^1\times 1)}=\mathrm{pr}(\beta )`$.
Put $`e`$ to be the Euler class of the contact bundle. Consider the locally-trivial $`S^1`$-fibration $`p:M^{}T^2`$ induced by $`\mu `$ from the $`S^1`$-fibration $`\mathrm{pr}:CMM`$. One can verify that $`2n=H^2(T^2,)`$ is the Euler class of $`p`$. On the other hand the Euler class of $`p`$ is $`\mu ^{}(e)`$ and it is naturally identified with the value of $`e`$ on the homology class realized by $`\mu (T^2)`$. This implies that if $``$ does not satisfy condition $`()`$, then there exists a homology class $`\alpha `$ from the statement of the Proposition.
On the other hand the existence of the class $`\alpha `$ from the statement of the Proposition implies that there exists a Legendrian curve $`K`$ such that $`\stackrel{}{K}`$ is conjugate to $`\stackrel{}{K}f^n`$, for $`n`$ being the value of $`e`$ (the Euler class of the contact bundle) on the homology class realized by $`\mu (T^2)`$. (Proposition 4.0.1 says that $`e=2\alpha `$, for some $`\alpha H^2(M,)`$, and hence $`n`$ is even.) This means that $``$ does not satisfy condition $`()`$ and we have proved Proposition 3.1.4. ∎
### 5.5. Proof of Theorem 3.1.8
###### 5.5.1.
Let $`\mathrm{pr}:CMM`$ be the locally trivial $`S^1`$-fibration introduced in 2.0.1 and let $`f\pi _1(CM)`$ be the class of an oriented $`S^1`$-fiber of $`\mathrm{pr}`$.
Similar to 5.4 we get that to prove that all the components of the space of framed curves satisfy condition $`()`$, it suffices to show that $`\stackrel{}{K}`$ and $`\stackrel{}{K}f^{2n}`$ are not conjugate in $`\pi _1(CM)`$, for all $`0n`$ and $`\stackrel{}{K}\pi _1(CM)`$. Let $`\stackrel{}{K},\beta \pi _1(CM)`$ and $`n`$ be such that
(6)
$$\beta \stackrel{}{K}\beta ^1=\stackrel{}{K}f^{2n}\pi _1(CM,\stackrel{}{K}(1)).$$
We have to show that $`n=0`$.
Proposition 5.1.1 says that $`f`$ is in the center of $`\pi _1(CM,\stackrel{}{K}(1))`$. Hence
(7)
$$\beta \stackrel{}{K}=\stackrel{}{K}\beta f^{2n}.$$
Identity (7) implies that $`\mathrm{pr}_{}(\beta )`$ and $`\mathrm{pr}_{}(\stackrel{}{K})`$ commute in $`\pi _1(M)`$. Hence there exists a mapping of the two-torus $`\mu :T^2=S^1\times S^1M`$ such that $`\mu (S^1\times 1)=\mathrm{pr}(\stackrel{}{K})`$ and $`\mu (1\times S^1)=\mathrm{pr}(\beta )`$. By the assumption of the Theorem $`\mu :\pi _1(T^2)\pi _1(M)`$ has a nontrivial kernel. Thus there exist $`i,j`$ with at least one of $`i`$ and $`j`$ being nonzero such that $`\mathrm{pr}(\stackrel{}{K})^i=\mathrm{pr}(\beta )^j\pi _1(M)`$, and hence
(8)
$$\stackrel{}{K}^i=\beta ^jf^l,\text{ for some }l.$$
Since the situation is symmetric, we assume that $`j0`$.
Thus $`\stackrel{}{K}^i\stackrel{}{K}^i=\stackrel{}{K}^i\beta ^jf^l=\beta ^jf^l\stackrel{}{K}^i`$. Applying (7) to the last identity we get that $`f^{2nij}=1`$. Since $`\pi _2(M)=0`$ we see that $`f`$ has infinite order in $`\pi _1(CM)`$, and hence $`2nij=0`$. If $`n`$ is zero, then we are done. Hence we have to look at the case of $`i=0`$. (We assumed that $`j0`$.) From (8) we get that $`\beta ^j=f^l`$, and hence by Proposition 5.1.1 $`\beta ^j`$ is in the center of $`\pi _1(CM)`$. Thus $`\beta ^j\stackrel{}{K}=\stackrel{}{K}\beta ^j`$. On the other hand using (7) we get that $`\beta ^j\stackrel{}{K}=\stackrel{}{K}\beta ^jf^{2nj}`$. Since $`f`$ has infinite order in $`\pi _1(CM)`$ we get that $`2nj=0`$. By our assumptions $`j0`$ and we have $`n=0`$. This finishes the proof of Theorem 3.1.8.∎
### 5.6. Proof of Theorem 3.1.10
By Theorem 3.1.1 it suffices to show that all the connected components of the space of framed curves in $`M`$ satisfy condition $`()`$.
Let $`eH^2(M,)`$ be the Euler class of the contact bundle of $`(M,C)`$. Proposition 3.1.4 implies that it suffices to show that $`e(\alpha )=0`$, for every homology class $`\alpha H_2(M)`$ realizable by a mapping $`\mu :T^2M`$. The result of D. Gabai (see Corollary 6.18 ) implies that every $`\alpha H^2(M,)`$ realizable by a mapping of $`T^2`$ can be realized by a collection of spheres and of a torus that are embedded into $`M`$. Finally the result of Ya. Eliashberg (see Theorem $`\mathrm{2.2.1}`$) says that for a tight contact structure the value of $`e`$ on any embedded torus or sphere is zero. Hence $`e(\alpha )=0`$. This finishes the Proof of Theorem 3.1.10.∎
### 5.7. Proof of Proposition 4.0.1
First we show that if $`eH^2(M,)`$ can be realized as the Euler class of the contact structure, then $`e=2\alpha `$ for some $`\alpha H^2(M,)`$.
Since the contact structure is cooriented we get that the tangent bundle $`TM`$ is isomorphic to the sum $`Cϵ`$ of the oriented contact bundle $`C`$ and the trivial oriented line bundle $`ϵ`$. The tangent bundle of every orientable $`3`$-manifold is trivializable and we get that the second Stiefel-Whitney class of the contact bundle is zero. But the second Stiefel-Whitney of $`C`$ is the projection of the Euler class of $`C`$ under the natural mapping $`H^2(M,)H^2(M,_2)`$, and we get that $`e=2\alpha `$ for some $`\alpha H^2(M,)`$.
Now we show that if $`e=2\alpha `$, for some $`\alpha H^2(M,)`$, then $`e`$ can be realized as the Euler class of a cooriented contact structure on $`M`$.
Consider an oriented $`2`$-dimensional vector bundle $`\xi `$ over $`M`$ with the Euler class $`e(\xi )=e=2\alpha H^2(M,)`$. The second Stiefel-Whitney class $`w_2(\xi )`$ of $`\xi `$ is zero, since it is the projection of $`e(\xi )=2\alpha H^2(M,)`$. Since $`\xi `$ is an oriented vector bundle we have $`w_1(\xi )=0`$.
Consider the sum $`\xi ϵ`$ of $`\xi `$ with the trivial oriented $`1`$-dimensional vector bundle $`ϵ`$. Clearly the total Stiefel-Whitney class of the $`3`$-dimensional oriented vector bundle $`\xi ϵ`$ is equal to $`1`$, and the Euler class of $`\xi ϵ`$ is equal to $`0`$. Using the interpretation of the Stiefel-Whitney and the Euler classes of $`\xi ϵ`$ as obstructions for the trivialization of $`\xi ϵ`$, we get that $`\xi ϵ`$ is trivializable. Since the tangent bundle of an oriented $`3`$-dimensional manifold is trivializable, we see that $`\xi `$ is isomorphic to an oriented sub-bundle of $`TM`$. Since $`M`$ is oriented this sub-bundle of $`TM`$ is also cooriented. Now the Theorem of Lutz , that says that every homotopy class of distributions of $`2`$-planes tangent to $`M`$ contains a contact structure, implies the existence of a cooriented contact structure with the Euler class $`e`$.∎
### 5.8. Proof of Theorem 4.1.1
###### 5.8.1.
Proof of statement a of Theorem 4.1.1. Clearly (see Figure 6) the two Legendrian knots $`K_0`$ and $`K_1`$ belong to the same component of the space of Legendrian curves. It is easy to see that if $`K_0`$ realizes the isotopy class of a framed knot $`\stackrel{~}{K}^0`$, then $`K_1`$ realizes the isotopy class of $`\stackrel{~}{K}^2`$ (see 5.1.9 for the definition of $`\stackrel{~}{K}^2`$). Below we show that $`\stackrel{~}{K}^0`$ and $`\stackrel{~}{K}^2`$ are isotopic framed knots.
Let $`t\times S^2S^1\times S^2`$ be the sphere that crosses $`\stackrel{~}{K}^0`$ at exactly one point, and let $`N=[0,1]\times S^2`$ be a thin tubular neighborhood of $`t\times S^2`$. Fix $`xS^2`$ (below called the North pole) and the direction in $`T_xS^2`$ (below called the zero meridian). We can assume that the knot $`\stackrel{~}{K}^0`$ inside $`N=[0,1]\times S^2`$ looks as follows: it intersects each $`y\times S^2N=[0,1]\times S^2`$ at the North pole of the corresponding sphere, and the framing of the knot is parallel to the zero meridian.
Consider an automorphism $`\nu :S^1\times S^2S^1\times S^2`$ that is identical outside of $`N=[0,1]\times S^2`$ such that it rotates each $`y\times S^2[0,1]\times S^2`$ by $`4\pi y`$ around the North pole in the clockwise direction. Clearly under this automorphism $`\stackrel{~}{K}^0`$ gets two extra negative twists of the framing and $`\nu (\stackrel{~}{K}^0)=\stackrel{~}{K}^2`$. On the other hand it is easy to see that $`\nu `$ is diffeotopic to the identity, since it corresponds to the contractible loop in $`SO(3)=P^3`$. Hence we see that $`\stackrel{~}{K}^0`$ and $`\stackrel{~}{K}^2`$ are isotopic framed knots. This finishes the proof of statement a of Theorem 4.1.1.
To prove statement b of the Theorem we need the following Proposition.
###### Proposition 5.8.2.
Let $`C`$ be a cooriented contact structure on $`M=S^1\times S^2`$ with a nonzero Euler class $`e`$ of the contact bundle. Let $`CM`$ be the spherical contact bundle, let $`\mathrm{pr}:CMM`$ be the corresponding locally trivial $`S^1`$-fibration, and let $`f\pi _1(CM)`$ be the class of an oriented $`S^1`$-fiber of $`\mathrm{pr}`$. Then $`f`$ is of finite order in $`\pi _1(CM)`$ and $`\pi _2(CM)=0`$.
###### 5.8.3.
Proof of Proposition 5.8.2. Consider the oriented $`2`$-plane bundle $`p:\xi S^2`$ that is the restriction of the contact bundle over $`M`$ to the sphere $`1\times S^2S^1\times S^2`$. The Euler class of $`p`$ is the value of $`e`$ on the homology class realized by $`1\times S^2`$, and hence is nonzero. Let $`S\xi `$ be the manifold obtained by the fiberwise spherization of $`p`$, and let $`\overline{p}:S\xi S^2`$ be the corresponding locally trivial $`S^1`$-fibration. Since the Euler class of $`p`$ is nonzero we get that a certain multiple of the class of the fiber of $`\overline{p}`$ is homologous to zero. But $`\pi _1(S\xi )`$ is generated by the class of the fiber, and hence the class of the fiber of $`\overline{p}`$ is of finite order in $`\pi _1(S\xi )`$. This implies that $`f\pi _1(CM)`$ is of finite order.
The statement that $`\pi _2(CM)=0`$ follows from the exact homotopy sequence of $`\mathrm{pr}:CMM`$ and the fact that $`f\pi _1(CM)`$ is of finite order. ∎
###### 5.8.4.
Proof of statement b of Theorem 4.1.1. Let $``$ be the connected component of the space of Legendrian curves that contains $`K_0`$ and $`K_1`$. Figure 6 shows that $`K_0`$ can be changed to $`K_1`$ (in the space of Legendrian curves) by a sequence of isotopies and one passage through a transverse double point. Hence if there exists a $``$-valued invariant $`I`$ of Legendrian knots from $``$ that increases by one under every positive passage through a transverse double point of a Legendrian knot, then it distinguishes $`K_0`$ and $`K_1`$. (Clearly if such $`I`$ does exist, then it is an order one invariant of Legendrian knots.) Below we show the existence of such $`I`$ in the connected component $``$.
Put $`I(K_0)=0`$. Let $`K^{}`$ be a Legendrian knot, and let $`\gamma `$ be a generic path in $``$ connecting $`K_0`$ and $`K^{}`$. Let $`J_\gamma `$ be the set of moments when $`\gamma `$ crosses the discriminant (i.e. the subspace of singular knots) in $``$, and let $`\sigma _j`$, $`jJ_\gamma `$, be the signs of these crossings. For a generic path $`\gamma `$ put $`\mathrm{\Delta }_I(\gamma )=_{jJ_\gamma }\sigma _j`$. It is clear that if $`I`$ (with $`I(K_0)=0`$) does exist, then $`I(K^{})=\mathrm{\Delta }_I(\gamma )`$. To show that $`I`$ does exist we have to verify that for every Legendrian knot $`K^{}`$ and for a generic path $`\gamma `$ connecting $`K^{}`$ to $`K_0`$ the value of $`\mathrm{\Delta }_I(\gamma )`$ does not depend on the choice of a generic path $`\gamma `$ connecting $`K_0`$ and $`K^{}`$, or equivalently we have to show that $`\mathrm{\Delta }_I(\gamma )=0`$ for every generic closed loop $`\gamma `$ connecting $`K_0`$ to itself.
There are two codimension two strata of the discriminant of $``$. They are formed respectively by singular Legendrian knots with two transverse double points, and by Legendrian knots with one double point at which the two intersecting branches are tangent of order one. Straightforward verification shows that $`\mathrm{\Delta }_I(\beta )=0`$, for every small closed loop $`\beta `$ going around a codimension two stratum of $``$.
This implies that for every generic loop $`\gamma `$ connecting $`K_0`$ to itself the value of $`\mathrm{\Delta }_I(\gamma )`$ depends only on the element of $`\pi _1(,K_0)`$ realized by $`\gamma `$. Hence to prove the existence of $`I`$ it suffices to show that $`\mathrm{\Delta }_I(\gamma )=0`$ for every $`\gamma \pi _1(,K_0)`$.
Clearly $`\mathrm{\Delta }_I(\gamma ^p)=p\mathrm{\Delta }_I(\gamma )`$ and since $``$ is torsion free, we get that to prove Theorem 4.1.1 it suffices to show that for every $`\gamma \pi _1(,K_0)`$ there exists a nonzero $`p`$ such that $`\mathrm{\Delta }_I(\gamma ^p)=0`$.
The $`h`$-principle says that the space of Legendrian curves in $`(M,C)`$ is weak homotopy equivalent to the space of free loops $`\mathrm{\Omega }CM`$ in the spherical contact bundle $`CM`$ of $`M=S^1\times S^2`$. (The mapping giving the equivalence lifts a Legendrian curve $`K`$ to a loop $`\stackrel{}{K}`$ in $`CM`$ by sending $`tS^1`$ to the point in $`CM`$ that corresponds to the direction of the velocity vector of $`K`$ at $`K(t)`$.)
Thus $`\pi _1(,K_0)`$ is naturally isomorphic to $`\pi _1(\mathrm{\Omega }CM,\stackrel{}{K}_0)`$. Proposition 5.8.2 says that $`\pi _2(CM)=0`$ and from 5.1.7 we get that $`\pi _1(\mathrm{\Omega }CM,\stackrel{}{K}_0)`$ is isomorphic to the centralizer $`Z(\stackrel{}{K}_0)`$ of $`\stackrel{}{K}_0\pi _1(CM,\stackrel{}{K}_0(1))`$. Using Propositions 5.1.1 and 5.8.2 we see that either $`\pi _1(CM)=`$ or $`\pi _1(CM)=_p`$, for some nonzero $`p`$. Hence there exists $`n`$ and nonzero $`m`$ such that $`\gamma ^m=\stackrel{}{K}_0^n\pi _1(CM,\stackrel{}{K}_0(1))`$. (One should take $`n`$ and $`m`$ to be divisible by $`p`$ if $`\pi _1(CM)=_p`$.) But the loop $`\alpha `$ in $`\pi _1(,K_0)`$ corresponding to $`\stackrel{}{K}_0^n`$ is just the sliding of $`K_0`$ $`n`$ times along itself according to the orientation. (This deformation is induced by the rotation of the parameterizing circle.) This loop does not intersect the discriminant, and hence $`\mathrm{\Delta }_I(\alpha )=0`$. This finishes the proof of statement b of Theorem 4.1.1.∎
### 5.9. Proof of Theorem 4.2.2
###### 5.9.1.
$`K_1`$ and $`K_2`$ are homotopic Legendrian curves and they realize isotopic framed knots. Let $`f_1\pi _1(CM)`$ be the class of the $`S^1`$-fiber of the fibration $`\mathrm{pr}:CMM`$. The $`h`$-principle says that the connected component of the space of Legendrian curves that contains $`K`$ corresponds to the conjugacy class of $`\stackrel{}{K}\pi _1(CM)`$. From the formula for the Maslov number deduced in it is easy to see that the connected components containing $`K_1`$ and $`K_2`$ correspond to the conjugacy classes of $`\stackrel{}{K}_1=\stackrel{}{K}f_1^r`$ and of $`\stackrel{}{K}_2=\stackrel{}{K}f_1^r`$. Let $`f_2\pi _1(CM)`$ be an element projecting to the class $`f\pi _1(M)`$ of the $`S^1`$-fiber of $`p:MF`$. The value of the Euler class of the contact bundle on the homology class realized by $`\mu (T^2)`$ is equal to $`2r`$. (Here $`\mu `$ is the mapping from the description of the Euler class of the contact bundle.) And because of the reasons explained in the proof of Proposition 3.1.4 we get that $`\stackrel{}{K}f_2=f_2\stackrel{}{K}f_1^{2r}`$, for the Legendrian knot $`K`$ used to construct $`K_1`$ and $`K_2`$. Now Proposition 5.1.1 implies that $`\stackrel{}{K}_1`$ and $`\stackrel{}{K}_2`$ are conjugate in $`\pi _1(CM)`$ and hence $`K_1`$ and $`K_2`$ are in the same component of the space of Legendrian curves.
The fact that $`K_1`$ and $`K_2`$ realize isotopic framed knots is clear, because as unframed knots they are the same, and as it is shown in every pair of extra cusps corresponds to the negative extra twist of the framing.
###### 5.9.2.
The idea of the proof of the fact that $`K_1`$ and $`K_2`$ can be distinguished by an order one invariant of Legendrian knots. Let $`d`$ be a point in $`M`$. Let $`K_s`$ be a singular unframed knot with one double point. The double point separates $`K_s`$ into two oriented loops. Deform $`K_s`$ preserving the double point, so that the double point is located at $`d`$. Choosing one of the two loops of $`K_s`$ we obtain an ordered set of two elements $`\delta _1,\delta _2\pi _1(M,d)`$, or which is the same an element $`\delta _1\delta _2\pi _1(M,d)\pi _1(M,d)`$. Clearly there is a unique element of the set $`B`$ that corresponds to the original singular unframed knot $`K_s`$, where $`B`$ is is the quotient set of $`\pi _1(M,d)\pi _1(M,d)`$ modulo the consequent actions of the following groups:
$`\pi _1(M)`$ whose element $`\xi `$ acts on $`\delta _1\delta _2\pi _1(M)\pi _1(M)`$ by sending it to $`\xi \delta _1\xi ^1\xi \delta _2\xi ^1\pi _1(M)\pi _1(M)`$. (This corresponds to the ambiguity in deforming $`K_s`$, so that the double point is located at $`d`$.)
$`_2`$ that acts via the cyclic permutation of the two summands. (This corresponds to the ambiguity in the choice of one of the two loops of $`K_s`$.)
Thus we have a mapping $`\nu `$ from the set of singular unframed knots with one double point to $`B`$. Let $`\alpha :B`$ be the function such that
$`\alpha (b)=0`$, provided that $`b`$ contains the class of $`1\delta \pi _1(M)\pi _1(M)`$ for some $`\delta \pi _1(M)`$,
$`\alpha (b)=1`$ otherwise.
Assume that $`I^{}`$ is an invariant of Legendrian knots from $``$ such that under every (generic transverse) positive passage through a discriminant in $``$ it increases by $`\alpha \nu (K_s)`$, where $`K_s`$ is the unframed singular knot corresponding to the crossing of the discriminant. Clearly such $`I^{}`$ is an order one invariant of framed knots from $``$. To prove the Theorem we show the existence of such $`I^{}`$, and then we show that it distinguishes $`K_1`$ and $`K_2`$.
###### 5.9.3.
The existence of $`I^{}`$. Let $`\gamma `$ be a generic path in $``$ that starts with $`K_1`$. Let $`J_\gamma `$ be the set of instances when $`\gamma `$ crosses the discriminant (i.e. the subspace of singular knots) in $``$, and let $`\sigma _j`$, $`jJ_\gamma `$, be the signs of these crossings. Let $`J_\gamma ^{}J_\gamma `$ be those instances for which the value of $`\alpha \nu `$ on the corresponding singular unframed knots is $`1`$. For a generic path $`\gamma `$ put $`\mathrm{\Delta }_I^{}(\gamma )=_{j^{}J_\gamma ^{}}\sigma _j^{}`$.
Similarly to 5.8 we get that to prove the existence of $`I^{}`$ it suffices to show that $`\mathrm{\Delta }_I^{}(\gamma )=0`$, for every generic closed loop $`\gamma `$.
Let $`𝒞`$ be the connected component of the space of unframed curves obtained by forgetting the framings on curves from $``$, and let $`K_1^{}`$ be the unframed knot obtained by forgetting the framing on $`K_1`$. Similarly to the above for a generic path $`\gamma `$ in $`𝒞`$ starting with $`K_1^{}`$ we put $`\mathrm{\Delta }_I^𝒞(\gamma )=_{j^{}J_\gamma ^{}}\sigma _j^{}`$. (As above $`J_\gamma ^{}`$ is the set of instances when the value of $`\alpha \nu `$ on the singular unframed knots obtained under $`\gamma `$ is equal to $`1`$, and $`\sigma _j^{}`$, $`j^{}J_\gamma ^{}`$, are the signs of the corresponding crossings of the discriminant.)
The codimension two stratum of the discriminant of $`𝒞`$ consists of singular curves whose only singularities are two distinct transverse double points. Straightforward verification shows that $`\mathrm{\Delta }_I^𝒞(\beta )=0`$ for any small loop $`\beta `$ going around the codimension two stratum. Hence $`\mathrm{\Delta }_I^𝒞:\pi _1(𝒞,K_1^{})`$ is a homomorphism.
There are two codimension two strata in $``$. They consist of respectively singular Legendrian curves whose only singularities are two distinct transverse double points and of singular Legendrian curves whose only singularity is one double point at which the two branches are tangent. Considerations similar to the ones above show that $`\mathrm{\Delta }_I^{}:\pi _1(,K_1)`$ is a homomorphism.
It is clear that if $`\gamma ^{}\pi _1(𝒞)`$ is the element corresponding to $`\gamma \pi _1()`$, then we have
(9)
$$\mathrm{\Delta }_I^{}(\gamma )=\mathrm{\Delta }_I^𝒞(\gamma ^{}).$$
The $`h`$-principle says that the space of Legendrian curves in $`(M,C)`$ is weak homotopy equivalent to the space of free loops in the spherical contact bundle $`CM`$ of $`M`$. (The mapping that gives the equivalence lifts a Legendrian curve $`K`$ in $`(M,C)`$ to a loop $`\stackrel{}{K}`$ in $`CM`$ by mapping $`tS^1`$ to the point of $`CM`$ that corresponds to the velocity vector of $`K`$ at $`K(t)`$.) Since $`\pi _2(CM)=0`$ for $`M`$ from the statement of the Theorem, we obtain (see 5.1.7) the natural isomorphism $`t:\pi _1(,K_1)Z(\stackrel{}{K}_1)<\pi _1(CM,\stackrel{}{K}_1(1))`$. Since $`\mathrm{\Delta }_I^{}(\gamma ^p)=p\mathrm{\Delta }_I^{}(\gamma )`$ and $``$ is torsion free, we get that to show the existence of $`I^{}`$ it suffices to show that for every $`\beta Z(\stackrel{}{K}_1)<\pi _1(CM,\stackrel{}{K}_1(1))`$ there exist $`0n`$ and $`\gamma \pi _1(,K_1)`$ such that $`t(\gamma )=\beta ^n\pi _1(CM,\stackrel{}{K}_1(1))`$ and $`\mathrm{\Delta }_I^{}(\gamma )=0`$.
Let $`f\pi _1(M,K_1(1))`$ be the class of the $`S^1`$-fiber of $`p:MF`$. Let $`f_1\pi _1(CM)`$ be the class of an oriented $`S^1`$-fiber of $`pr:CMM`$, and let $`f_2`$ be an element of $`\pi _1(CM)`$ such that $`pr_{}(f_2)=f\pi _1(M,K_1(1))`$.
Take $`\beta Z(\stackrel{}{K}_1)`$, then $`\mathrm{pr}_{}(\beta )Z(K_1)`$. Proposition 5.1.5 implies that there exist $`0n`$ and $`i,j`$ such that $`K_1^if^j=(\mathrm{pr}_{}(\beta ))^n\pi _1(M,K_1(1))`$. Using Proposition 5.1.1 we get that
(10)
$$\beta ^n=\stackrel{}{K}_1^if_2^jf_1^l\text{ for some }i,j,l.$$
As it was explained in 5.9.1 we have
(11)
$$\stackrel{}{K}_1f_2=f_2\stackrel{}{K}_1f_1^{2r}.$$
Since $`\beta Z(\stackrel{}{K}_1)`$ we get that $`\stackrel{}{K}_1\beta ^n=\beta ^n\stackrel{}{K}_1`$, and using (10) we see that $`\stackrel{}{K}_1\stackrel{}{K}_1^if_2^jf_1^l=\stackrel{}{K}_1^if_2^jf_1^l\stackrel{}{K}_1`$. Using (11), Proposition 5.1.1, and the fact that $`f_1`$ has infinite order in $`\pi _1(CM)`$, we see that $`j=0`$ in (10).
Hence
(12)
$$\beta ^n=\stackrel{}{K}_1^if_1^l,\text{ for some }i,l.$$
Clearly $`f_1Z(\stackrel{}{K}_1)`$ and since $`t:\pi _1(,K_1)Z(\stackrel{}{K}_1)`$ is surjective there exists a loop $`\gamma _3\pi _1(,K_1)`$ such that $`t(\gamma _3)=f_1`$. Let $`\gamma _2\pi _1(,K_1)`$ be the loop corresponding to the deformation under which $`K_1`$ slides once around itself according to the orientation of $`K_1`$. (This deformation is induced by the rotation of the circle parameterizing $`K_1`$.) Clearly $`\gamma _2`$ does not cross the discriminant and hence $`\mathrm{\Delta }_I^{}(\gamma _2)=0`$.
To prove the existence of $`I^{}`$ it suffices to show that $`\mathrm{\Delta }_I^{}(\gamma )=0`$, for $`\gamma \pi _1(,K_1)`$ such that $`t(\gamma )=\stackrel{}{K}_1^if_1^l`$. But this $`\gamma `$ is $`\gamma _2^i\gamma _3^l`$. Thus it suffices to show that $`0=\mathrm{\Delta }_I^{}(\gamma )=i\mathrm{\Delta }_I^{}(\gamma _2)+l\mathrm{\Delta }_I^{}(\gamma _3)`$. Since $`\mathrm{\Delta }_I^{}(\gamma _2)=0`$ we get that $`\mathrm{\Delta }_I^{}(\gamma )=l\mathrm{\Delta }_I^{}(\gamma _3)`$ and thus it suffices to show that $`\mathrm{\Delta }_I^{}(\gamma _3)=0`$.
Let $`\gamma _3^{}\pi _1(𝒞,K_1^{})`$ be the loop corresponding to $`\gamma _3\pi _1(,K_1)`$. Identity (9) says that $`\mathrm{\Delta }_I^{}(\gamma _3)=\mathrm{\Delta }_I^𝒞(\gamma _3^{})`$. Hence we have to show that $`\mathrm{\Delta }_I^𝒞(\gamma _3^{})=0`$.
In 5.9.4 we show that $`\gamma _3^{}\pi _1(𝒞,K_1^{})`$ can be realized as a power of the loop described by the deformation shown in Figure 7. Then since one of the loops of the only singular knot arising under this deformation is contractible and the value of $`\alpha \nu `$ on such a singular knot is zero we get that $`\mathrm{\Delta }_I^𝒞(\gamma _3^{})=0`$. This finishes the Proof of the existence of $`I^{}`$ modulo the explanations given in 5.9.4.
###### 5.9.4.
Now we show that $`\gamma _3^{}\pi _1(𝒞,K_1^{})`$ can be realized as a sequence of loops described by the deformation shown in Figure 7. The $`h`$-principle for curves 5.1.8 says that $`𝒞`$ is weak homotopy equivalent to the space of free loops $`\mathrm{\Omega }STM`$ in the spherical tangent bundle $`STM`$ of $`M`$. In particular $`\pi _1(𝒞,K_1^{})=\pi _1(\mathrm{\Omega }STM,\stackrel{}{K}_1^{})`$. In Subsubsection 5.1.7 we introduced a surjective homomorphism $`t`$ from $`\pi _1(\mathrm{\Omega }STM,\stackrel{}{K}_1^{})`$ onto $`Z(\stackrel{}{K}_1^{})<\pi _1(STM,\stackrel{}{K}_1^{}(1))`$.
Let $`\alpha ,\beta \pi _1(𝒞,K_1^{})`$ be loops such that $`t(\alpha )=t(\beta )`$. As it was explained in 5.1.7 the obstruction for $`\alpha `$ and $`\beta `$ to be homotopic is an element of $`\pi _2(STM)`$. Since every orientable $`3`$-manifold is parallelizable we get that $`STM=S^2\times M`$. Clearly $`\pi _2(M)=0`$ for $`M`$ from the statement of the Theorem and hence $`\pi _2(STM)=\pi _2(S^2)=`$.
Consider the loop $`\alpha ^{}`$ that looks the same as $`\alpha `$ except for a small period of time when we perform the deformation shown in Figure 7. Clearly $`t(\alpha ^{})=t(\alpha )=t(\beta )\pi _1(STM)`$, and straightforward verification show that the $``$-valued obstruction for $`\alpha ^{}`$ and $`\beta `$ to be homotopic differs by one from the obstruction for $`\alpha `$ and $`\beta `$ to be homotopic. Hence performing this operation (or its inverse) sufficiently many times we can change $`\alpha `$ to be homotopic to $`\beta `$.
Since $`t(\gamma _3^{})=t(1)=1\pi _1(STM)`$ we get that $`\gamma _3^{}\pi _1(𝒞,K_1^{})`$ can be realized as a power of the deformation shown in Figure 7.
###### 5.9.5.
Let us show that $`I^{}`$ distinguishes $`K_1`$ and $`K_2`$. Let $`\rho :[0,1]`$ be a generic path connecting $`K_1`$ and $`K_2`$. To prove the Theorem we have to show that $`\mathrm{\Delta }_I^{}(\rho )0`$.
Let $`K_1^{}`$ (resp $`K_2^{}`$) be the unframed knot obtained by forgetting the framing on $`K_1`$ (resp $`K_2`$). Let $`\rho ^{}:[0,1]𝒞`$ be the isotopy that deforms $`K_1^{}`$ into $`K_2^{}`$ in the category of unframed curves under which $`K_1^{}`$ all the time stays in a thin tubular neighborhood of $`K_2^{}`$. Consider a homotopy $`\overline{\rho }=S^1𝒞`$ that corresponds to a product of paths $`\rho ^{}\rho `$. ($`\overline{\rho }`$ connects $`K_1^{}`$ to itself.) Clearly $`\mathrm{\Delta }_{I^𝒞}(\overline{\rho })=\mathrm{\Delta }_I^{}(\rho )`$.
For a loop $`\alpha :S^1𝒞`$ (that connects $`K_1^{}`$ to itself) put $`\lambda ^\alpha :T^2=S^1\times S^1M`$ to be a mapping such that for every $`tS^1`$ the mappings $`\lambda ^\alpha |_{(t\times S^1)}`$ and $`\alpha (t):S^1M`$ are the same. The value of the Euler class of the contact bundle on the homology class realized by $`\lambda ^{\overline{\rho }}(T^2)`$ is equal to $`2r0`$.
Using the usual arguments we get that to prove the Theorem it suffices to show that for every $`\gamma \pi _1(𝒞,K_1^{})`$ there exists $`n0`$ such that either
the value of the Euler class of the contact bundle on the homology class realized by $`\lambda ^{\gamma ^n}:(T^2)M`$ is zero, or
$`\mathrm{\Delta }_{I^𝒞}(\gamma ^n)0`$.
Consider the following loops $`\gamma _1`$ and $`\gamma _2`$.
Loop $`\gamma _1`$. Since $`p(K_1^{})`$ is an orientation preserving loop and $`M`$ is orientable, we get that the $`S^1`$-fibration over $`S^1`$ (parameterizing the knots) induced from $`p:MF`$ by $`pK_1^{}:S^1F`$ is trivializable. Hence we can coherently orient the fibers of this fibration. The orientation of the $`S^1`$-fiber over $`tS^1`$ induces the orientation of the $`S^1`$-fiber of $`p`$ that contains $`K_1^{}(t)`$. The loop $`\gamma _1`$ is the deformation of $`K_1^{}`$ under which every point of $`K_1^{}`$ slides once around the fiber of $`p`$ that contains this point (staying inside the fiber) in the direction specified by the orientation of the fiber corresponding to this point.
The loop $`\gamma _2`$ is the sliding of $`K_1^{}`$ along itself according to the orientation. (This deformation is induced by the rotation of the circle that parameterizes $`K_1^{}`$.)
The $`h`$-principle for curves says that $`𝒞`$ is weak homotopy equivalent to the space of free loops $`\mathrm{\Omega }STM`$ in the spherical tangent bundle $`STM`$ of $`M`$. (The mapping that gives the equivalence lifts a curve $`K`$ in $`M`$ to a loop $`\stackrel{}{K}`$ in $`STM`$ by mapping $`tS^1`$ to the point of $`STM`$ that corresponds to the velocity vector of $`K`$ at $`K(t)`$.) Let $`t:\pi _1(𝒞,K_1^{})Z(\stackrel{}{K}_1^{})<\pi _1(\mathrm{\Omega }STM,\stackrel{}{K}_1^{})`$ be the surjective homomorphism described in 5.1.7, and let $`\overline{f}\pi _1(STM)`$ be the element that projects to the class $`f\pi _1(M)`$ of the $`S^1`$-fiber of $`p:MF`$.
Using Proposition 5.1.5 one verifies that for every $`\gamma \pi _1(𝒞,K_1^{})`$ there exist $`0n`$ such $`t(\gamma ^n)=\overline{f}^i(\stackrel{}{K}_1^{})^j=t(\gamma _1)^it(\gamma _2)^j`$, for some $`i,j`$. Let $`\gamma _4`$ be the loop described in Figure 7. (It is easy to see that $`\gamma _4`$ is in the center of $`\pi _1(𝒞,K_1^{})`$.) Similar to 5.9.4 we get that $`\gamma ^n=\gamma _1^i\gamma _2^j\gamma _4^k`$ for some $`k`$.
It is easy to see that if $`i=0`$ then the value of the Euler class on the homology class realized by $`\lambda ^{\gamma ^n}:T^2M`$ is zero, and hence a holds for $`\gamma `$.
On the other hand if $`i0`$ then as we show below in 5.9.6 $`\mathrm{\Delta }_{I^𝒞}(\gamma ^n)0`$ and b holds for $`\gamma `$. (This finishes the proof of the Theorem modulo the explanation below.)
###### 5.9.6.
If $`i0`$ then $`\mathrm{\Delta }_{I^𝒞}(\gamma ^n)0`$ and b holds for $`\gamma `$. The loop $`\gamma _1`$ crosses the discriminant twice, both crossings occur with the same sign and the values of $`\alpha \nu `$ on the corresponding singular knots are equal to one. These crossings occur in the fiber over the double point of $`p(K_1^{})`$. (Since the double point of $`p(K_1^{})`$ separates it into two orientation reversing loops, the two points of $`K_1^{}`$ contained in this fiber induce opposite orientations of it, and the two branches of $`K_1^{}`$ that intersect the fiber slide in the opposite directions under $`\gamma _1`$.) Hence $`\mathrm{\Delta }_{I^𝒞}(\gamma _1)=\pm 20`$. (The sign depends on the orientation of the $`S^1`$-fibers of $`T^2S^1`$ used to induce the orientations of the fibers containing the points of $`K_1^{}`$.)
Clearly $`\mathrm{\Delta }_{I^𝒞}(\gamma _2)=0`$.
Since one of the two loops of the only singular knot appearing in $`\gamma _4`$ is contractible, we have that the value of $`\alpha \nu `$ on the singular knot is $`0`$ and thus $`\mathrm{\Delta }_{I^𝒞}(\gamma _4)=0`$.
Hence if $`i0`$, then $`\mathrm{\Delta }_{I^𝒞}(\gamma ^n)=i\mathrm{\Delta }_{I^𝒞}(\gamma _1)+j\mathrm{\Delta }_{I^𝒞}(\gamma _2)+k\mathrm{\Delta }_{I^𝒞}(\gamma _4)=i\mathrm{\Delta }_{I^𝒞}(\gamma _1)=i(\pm 2)0`$ and hence b holds for $`\gamma `$.
This finishes the proof of Theorem 4.2.2. ∎
Acknowledgments. I am very grateful to Stefan Nemirovski, Serge Tabachnikov and Oleg Viro for the valuable discussions and suggestions. I am deeply thankful to H. Geiges and A. Stoimenow for the valuable suggestions, and to O. Baues, M. Bhupal, N. A’ Campo, A. Cattaneo, J. Fröhlich, J. Latschev, A. Shumakovich, and V. Turaev for many valuable discussions.
This paper was written during my stay at the Max-Planck-Institut für Mathematik (MPIM), Bonn, and it is a continuation of the research conducted at the ETH Zurich . I would like to thank the Directors and the staff of the MPIM and the staff of the ETH for hospitality and for providing the excellent working conditions.
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# Lineshape predictions via Bethe ansatz for the one-dimensional spin-1/2 Heisenberg antiferromagnet in a magnetic field
## I Introduction
Advances in experimental techniques combined with improvements in sample preparation make it possible to produce data of ever increasing resolution for the quantum fluctuations and the underlying collective excitations in quasi-one-dimensional (1D) magnetic compounds. Advances in the theoretical analysis of relevant model systems combined with progress in the computational treatment of aspects that remain elusive to exact analysis make it possible to gain an ever more profound understanding of the observable collective excitations in terms of a small number of constituent quasi-particles.
There is scarcely a better case for illustrating this multi-track advancement of understanding quantum fluctuations than the 1D $`s=\frac{1}{2}`$ Heisenberg antiferromagnet and the growing number of materials that have been discovered to be physical realizations of this model system. The Hamiltonian for $`N`$ spins $`\frac{1}{2}`$ arranged in a cyclic chain with isotropic exchange coupling $`J`$ between nearest neighbors and a uniform magnetic field $`h`$,
$$H=\underset{n=1}{\overset{N}{}}\left[J𝐒_n𝐒_{n+1}hS_n^z\right],$$
(1)
is amenable to exact analysis via Bethe ansatzBeth31 ; FT81 and displays dynamical properties of intriguing complexity. The field $`h`$ is a controllable continuous parameter, which leaves the eigenvectors unaltered, but changes the nature of the ground state via level crossings and thus has a strong impact on the dynamical properties, in particular at low temperatures.
At $`hh_S=2J`$ the ground state of $`H`$ has all spins aligned in field direction: $`|F|\mathrm{}`$ is the reference state of the coordinate Bethe ansatz, and all eigenstates are described as excitations of interacting magnons, a species of spin-1 quasi-particles. Hence $`|F`$ plays the role of the magnon vacuum. The ground state $`|A`$ of $`H`$ at $`h=0`$ contains $`N/2`$ magnons. The Bethe ansatz enables us to reconfigure this state as the physical vacuum for a different species of quasi-particles – the spinons, which have spin $`\frac{1}{2}`$. The entire spectrum of the Heisenberg model (1) can also be generated as composites of interacting spinons.note1
Both descriptions are valid throughout the spectrum, but the magnon interpretation is more useful near the magnon vacuum, and the spinon picture is more useful near the spinon vacuum. The interaction energy of magnon scattering states or spinon scattering states is of O($`N^1`$) as long as the number of quasi-particles in the collective excitations is of O(1).note2 In a macroscopic system, the spectrum of such states is thus indistinguishable from the corresponding free quasi-particle states. Even under these simplifying circumstances, however, the interaction of the quasi-particles remains important in the make-up of collective wave functions, and is likely to strongly affect the transition rates and lineshapes.
At intermediate values $`0<h<h_S`$ of the magnetic field, the number of magnons or spinons contained in the ground state of $`H`$ is of O($`N`$), implying that the interaction energy for either quasi-particle species remains nonzero for $`N\mathrm{}`$ in the ground state and in all low-lying excitations. This obscures the role of individual magnons or spinons in the collective excitations and obstructs the interpretation of spectral data obtained by experimental or computational probes.
We can circumvent this problem by configuring the ground state $`|G`$ at $`0<M_z/N<\frac{1}{2}`$ as the physical vacuum for yet a different species of quasi-particles. From the new vantage point, the dynamically relevant collective excitations are then again scattering states of few quasi-particles with an interaction energy of O($`N^1`$), which greatly facilitates the interpretation of the spectra probed experimentally or computationally.
## II Dynamic Structure Factor
In an inelastic neutron scattering experiment performed at low temperature, the observable scattering events predominantly involve transitions from the ground state to a subset of collective excitations filtered from the rest by selection rules and transition rates. Under idealized circumstances, the scattering cross section is proportional to the $`T=0`$ dynamic spin structure factor
$$S_{\mu \mu }(q,\omega )=2\pi \underset{\lambda }{}|G|S_q^\mu |\lambda |^2\delta \left(\omega \omega _\lambda \right),$$
(2)
where $`S_q^\mu =N^{1/2}_ne^{iqn}S_n^\mu `$, $`\mu =x,y,z`$ is the spin fluctuation operator. In a macroscopic system, the aggregate of spectral lines in (2) pertaining to scattering events with energy transfer $`\omega _\lambda E_\lambda E_G`$, momentum transfer $`qk_\lambda k_G`$, and transition rate $`|G|S_q^\mu |\lambda |^2`$ form characteristic patterns of spectral weight in $`(q,\omega )`$-space. The shape of the spectral weight distribution provides key information on how the dynamically relevant collective excitations are composed of quasi-particles with specific energy-momentum relations.
Experimentally it is possible, at least in principle, to separate the information contained in the dynamic structure factors of the spin components parallel and perpendicular to the field direction, i.e. the functions $`S_{zz}(q,\omega )`$ and $`S_{xx}(q,\omega )=\frac{1}{4}[S_+(q,\omega )+S_+(q,\omega )]`$, respectively, for the fluctuation operators $`S_q^z`$ and $`S_q^\pm =S_q^x\pm iS_q^y`$. At $`h=0`$ the additional symmetry of $`H`$ dictates that $`S_{zz}(q,\omega )=\frac{1}{2}S_+(q,\omega )=\frac{1}{2}S_+(q,\omega )`$.
An anchor point for the new results presented in the following is the exact 2-spinon dynamic spin structure factor at $`T=0`$, which was determined recently via algebraic analysis and shown to contribute 73% of the total intensity in $`S_{zz}(q,\omega )`$ at $`h=0`$.KMB+97 Given the energy-momentum relationFT81
$$ϵ_{sp}(p)=\frac{\pi }{2}J\mathrm{sin}p,0p\pi ,$$
(3)
of the spinon quasi-particle, the 2-spinon states with wave numbers $`q=p_1+p_2`$ and energy $`\omega =ϵ_{sp}(p_1)+ϵ_{sp}(p_2)`$ form a continuum confined by the boundariesDP62 ; Yama69
$$ϵ_L(q)=\frac{\pi }{2}J|\mathrm{sin}q|,ϵ_U(q)=\pi J\left|\mathrm{sin}\frac{q}{2}\right|,$$
(4)
as illustrated in Fig. 1 (inset). The main plot shows the exact 2-spinon lineshapes of $`S_{zz}(q,\omega )`$ at $`q=\pi /2,3\pi /4,\pi `$. The most detailed experimental data available for testing these results pertain to KCuF<sub>3</sub>.TCNT95
We shall see that the magnetic field causes dramatic changes in both the spectrum and the lineshapes. At the root of these changes is a change in the nature of the relevant quasi-particles. Two compounds suitable for studying magnetic-field effects on spectrum and lineshapes are Cu(C<sub>6</sub>D<sub>5</sub>COO)$`{}_{2}{}^{}`$3D<sub>2</sub>O and Cu(C<sub>4</sub>H<sub>4</sub>N<sub>2</sub>)(NO<sub>3</sub>)<sub>2</sub>.DHR+97 ; HSR+99
## III Bethe Ansatz Equations
The Bethe ansatzBeth31 is an exact method for the calculation of eigenvectors of integrable quantum many-body systems. The Bethe wave function of any eigenstate of (1) in the invariant subspace with $`r=N/2M_z`$ reversed spins relative to the magnon vacuum,
$$|\psi =\underset{1n_1<\mathrm{}<n_rN}{}a(n_1,\mathrm{},n_r)S_{n_1}^{}\mathrm{}S_{n_r}^{}|F,$$
(5)
has coefficients of the form
$`a(n_1,\mathrm{},n_r)={\displaystyle \underset{𝒫S_r}{}}\mathrm{exp}`$ $`(i{\displaystyle \underset{j=1}{\overset{r}{}}}k_{𝒫j}n_j+{\displaystyle \frac{i}{2}}{\displaystyle \underset{i<j}{\overset{r}{}}}\theta _{𝒫i𝒫j})`$ (6)
determined by $`r`$ magnon momenta $`k_i`$ and one phase angle $`\theta _{ij}=\theta _{ji}`$ for each magnon pair. The sum $`𝒫S_r`$ is over the permutations of the labels $`\{1,2,\mathrm{},r\}`$.
The consistency requirements for the coefficients $`a(n_1,\mathrm{},n_r)`$ inferred from the eigenvalue equation $`H|\psi =E|\psi `$ and the requirements imposed by translational invariance lead to a set of coupled nonlinear equations for the $`k_i`$ and $`\theta _{ij}`$. A computationally convenient rendition of the Bethe ansatz equations has the form
$$N\varphi (z_i)=2\pi I_i+\underset{ji}{}\varphi \left[(z_iz_j)/2\right],i=1,\mathrm{},r,$$
(7)
where $`\varphi (z)2\mathrm{arctan}z`$, $`k_i=\pi \varphi (z_i)`$ and $`\theta _{ij}=\pi \mathrm{sgn}[\mathrm{}(z_iz_j)]\varphi \left[(z_iz_j)/2\right]`$. Every solution of (7) is specified by a set of Bethe quantum numbers $`I_1<I_2<\mathrm{}<I_r`$, which assume integer values for odd $`r`$ and half-integer values for even $`r`$. The energy and wave number of the eigenvector thus determined are
$$\frac{EE_F}{J}=\underset{i=1}{\overset{r}{}}\frac{2}{1+z_i^2},k=\pi r\frac{2\pi }{N}\underset{i=1}{\overset{r}{}}I_i,$$
(8)
where $`E_F=JN/4`$ is the energy of the magnon vacuum.
We consider the class $`K_r`$ of eigenstates whose Bethe quantum numbers comprise, for $`0rN/2`$ and $`0mN/2r`$, all configurations
$$\frac{r}{2}+\frac{1}{2}mI_1<I_2<\mathrm{}<I_r\frac{r}{2}\frac{1}{2}+m.$$
(9)
Here we employ the solutions $`\{z_i\}`$ of the Bethe ansatz equations not only to generate spectral data via (8), which is standard practice, but also to evaluate transition rates $`|G|S_q^\mu |\lambda |^2`$ for the dynamic structure factor (3) directly from the normalized Bethe wave functions $`|\lambda |\psi /\psi `$. The computational aspects of this method are discussed elsewhere.KHM00
## IV Physical Vacuum and Quasi-Particles
The ground-state wave function $`|G`$ at $`0M_zN/2`$ is specified by the set of $`r=N/2M_z`$ Bethe quantum numbersYY66a
$$\{I_i\}_G=\{\frac{N}{4}+\frac{M_z}{2}+\frac{1}{2},\mathrm{},\frac{N}{4}\frac{M_z}{2}\frac{1}{2}\}.$$
(10)
As the magnetic field increases from $`h=0`$ to $`h_S=2J`$, the magnetization $`M_z`$ increases in units of one from zero to $`N/2`$. A sequence of level crossings produces a magnetization curve $`(M_z/N`$ versus $`h)`$ in the form of a staircase with $`N/2`$ steps of height $`1/N`$, which converges toward a smooth line as $`N\mathrm{}`$.Grif64 ; BF64 ; KHM98
Depending on the reference state used for the characterization of the ground state $`|G`$, it can be regarded as a scattering state of $`N/2M_z`$ magnons excited from the magnon vacuum $`|F`$ or as a scattering state of $`2M_z`$ spinons excited from the spinon vacuum $`|A`$. To illustrate the distinct roles played by the two species of quasi-particles in the class-$`K_r`$ states, we show in Fig. 2 the configuration of Bethe quantum numbers for $`|G`$ in a system with $`N=8`$ and all values of $`M_z`$ realized between $`h=0`$ and $`h=h_S`$. The positions of the magnons ($``$) are determined by the set (10) of $`I_i`$’s and the positions of the spinons ($``$) by the vacancies across the full range of the $`I_i`$’s allowed by (9) for class $`K_r`$ states.
Henceforth we treat $`|G`$ as the new physical vacuum. At $`h=0`$ (top row in Fig. 2) it coincides with the spinon vacuum, a state with $`N/2`$ magnons. At $`h=h_S`$ (bottom row) it coincides with the magnon vacuum, a state containing $`N`$ spinons. All states within class $`K_r`$ are generated from $`|G`$ by rearranging the magnons or (equivalently) the spinons into all allowed configurations.
For $`r=N/2`$ (top row) and $`r=0`$ (bottom row) the state shown is the only possible configuration within class $`K_r`$. In the fourth row, the lone magnon can be moved across the array of spinons, generating a branch (one-parameter set) of 1-magnon excitations for $`N\mathrm{}`$. In the second row, the two spinons can be moved independently across the array of magnons, generating a continuum (two-parameter set) of 2-spinon excitations for $`N\mathrm{}`$ with boundaries (4) as shown in Fig. 1. The center row in Fig. 2 pertains to the field at half the saturation magnetization $`(M_z=r=N/4)`$, the case we shall investigate extensively for various system sizes. Here $`|G`$ contains twice as many spinons as it contains magnons.
The integer $`m`$ with range $`0<mM_z`$ used in (9) is a convenient quantum number for the subdivision of the classes $`K_r`$. Every state of $`K_r`$ at fixed $`m`$ can then be regarded as a scattering state of $`m`$ pairs of spinon-like quasi-particles. To distinguish them from the spinons, we name the new quasi-particles psinons.
The ground state $`|G`$, the only state with $`m=0`$, is the psinon vacuum. Here the magnons form a single array flanked by two arrays of spinons (see Fig. 2). Relaxing the constraint in (9) from $`m=0`$ to $`m=1`$ yields a two-parameter set of states – the 2-psinon excitations. Here the array of magnons breaks into three clusters separated by the two innermost spinons, which now assume the role of psinons. The remaining $`2M_z2`$ spinons stay sidelined. In the 4-psinon states $`(m=2)`$, two additional spinons have been mobilized into psinons. By this prescription, we can systematically generate sets of $`2m`$-psinon excitations for $`0mM_z`$.
To illustrate the quasi-particle role of the psinons in the class-$`K_r`$ collective states we have plotted in Fig. 3 energy versus wave number of all 2-psinon states (circles) and 4-psinon states (squares) at $`M_z=N/4`$ for $`N=16`$. Also shown are the spectral boundaries of 2-psinon and 4-psinon states for $`N\mathrm{}`$ as inferred from solutions of (7) for $`N=2048`$. The 2-psinon continuum, outlined by thick lines, is confined to the interval $`|q|q_s`$, where
$$q_s\pi (12M_z/N)$$
(11)
denotes the wave number of an incommensurate soft mode. The lower 4-psinon spectral boundary is the same as the 2-psinon lower boundary but extended periodically over the entire Brillouin zone. The upper 4-psinon boundary is related to the upper 2-psinon boundary by a scale transformation $`(q2q,\omega 2\omega )`$.
The relationship between the ranges in $`(q,\omega )`$-space of the 2-psinon states and the 4-psinon states does indeed reflect the fact that they are scattering states of two or four quasi-particles, respectively, of the same species. Like the spinon, the psinon is not observable in isolation via neutron scattering, but its energy momentum relation $`ϵ_\psi (p),\pi /4p\pi /4`$, can be inferred from the data of Fig. 3 (see inset).
If there were no psinon interaction, the wave number and energy of a $`2m`$-psinon state would be $`q=_{i=1}^{2m}p_i`$, $`\omega =_{i=1}^{2m}ϵ_\psi (p_i)`$. The $`N=16`$ data make it quite clear that the finite-size energy correction caused by the psinon interaction is stronger in the 4-psinon states than in the 2-psinon states. In both sets of collective states, the interaction energy goes to zero as the scattering events become less and less frequent in a chain of increasing length. However, it takes longer chains for finite-$`N`$ 4-psinon data to reach comparable convergence toward the spectral boundaries predicted for $`N\mathrm{}`$, because for fixed $`N`$, the scattering events between psinons are more numerous in a typical 4-psinon state than in a typical 2-psinon state.
If instead of the psinon vacuum we had used the spinon vacuum as the reference state at $`M_z=N/4`$, then both the 2-psinon states and the 4-psinon states would have to be described as scattering states of $`N/2`$ spinons. Although we know the energy-momentum relation of a spinon, Eq. (3), it is of little use to determine the spectral threshold in Fig. 3. Since the 2-psinon and 4-psinon states maintain a finite density of spinons in the limit $`N\mathrm{}`$, the spinon interaction energy remains significant. This problem does not arise at $`M_z=0`$. In the 2-spinon scattering states depicted in Fig. 1, the spinon interaction energy vanishes for $`N\mathrm{}`$ just as the psinon interaction energy does in the 2-psinon and 4-psinon scattering states depicted in Fig. 3.note7
## V Dynamically relevant excitations
At $`M_z=0`$ the spectral weight in the dynamic spin structure factor $`S_{zz}(q,\omega )`$ is dominated by the 2-spinon excitations.KMB+97 Our task here is to determine how the spectral weight of $`S_{zz}(q,\omega )`$ at $`M_z0`$ is distributed among the $`2m`$-psinon excitations. In investigating this question, we follow the strategy of an older studyMTBB81 but with vastly improved conceptual and numerical tools.
We begin by exploring, in a chain of $`N=16`$ spins at $`M_z/N=\frac{1}{4}`$, the transition rates between the ground state $`|G`$ and all $`2m`$-psinon excitations for $`m=0,1,2,3,4`$. The Bethe quantum numbers of the states with $`m=0,1`$ are shown in Fig. 4. The first row represents the psinon vacuum with its four magnons sandwiched by two sets of four spinons. The two innermost spinons (marked grey) become psinons when at least one of them is moved to another position. In the rows underneath, the psinons are moved systematically across the array of magnons while the remaining spinons stay frozen in place. These eight configurations describe all 2-psinon states with $`q0`$.
The wave numbers, energies, and transition rates of the states shown in Fig. 4 are listed in Table 1. Remarkably, almost the entire 2-psinon spectral weight is concentrated in the lowest excitation for any given $`q`$. The dynamically dominant 2-psinon states are marked by solid circles in Fig. 3. In a macroscopic system, they form the lower boundary of the 2-psinon continuum.
Next we calculate the transition rates $`|G|S_q^z|\lambda |^2`$ for the complete set of 4-psinon states. Interestingly, we observe that most of the 4-psinon spectral weight is again carried by a single branch of excitations. The dynamically dominant 4-psinon states for $`N=16`$ are shown as full squares in Fig. 3. For large $`N`$ they form a branch adjacent to the 2-psinon spectral threshold.
An investigation of the remaining $`2m`$-psinon states shows that there exists one dynamically dominant branch of $`2m`$-psinon excitations for $`0<mM_z`$. The configurations of Bethe quantum numbers pertaining to the four branches for $`N=16`$, each consisting of $`N/2M_z=4`$ states (at $`q>0`$), are shown in Fig. 5. The energies, wave numbers, and transition rates of these excitations are listed in Table 2. All other $`2m`$-psinon excitations have transition rates that are smaller by at least two orders of magnitude at $`q<\pi /2`$, and still by more than one order of magnitude at $`q\pi /2`$.
Inspection of Fig. 5 reveals an interesting pattern, indicative of the composition of the dynamically relevant collective excitations. They form a two-parameter set. The two parameters are highlighted by grey circles. Hitherto we have interpreted each group of four configurations as a branch of $`2m`$-psinon excitations, which are seemingly arbitrary one-parameter subsets taken from $`2m`$-parameter sets of states. In a macroscopic system, all but the lowest such branches contain a macroscopic number of psinons. Hence the range of the dynamically relevant excitations in $`(q,\omega )`$-space cannot be inferred from the psinon energy-momentum relation alone as was possible for the 2-psinon and 4-psinon continua, because the psinon interaction energy will remain non-negligible in most of these states for $`N\mathrm{}`$, just as the spinon interaction energy was non-negligible in the 2-psinon and 4-psinon scattering states at $`M_z0`$.
A more natural interpretation of the pattern on display in Fig. 5 identifies one of the two parameters as a psinon (large grey circle) as before and the other parameter as a new quasi-particle (small grey circle). The latter is represented by a hole in what was one of two spinon arrays in the psinon vacuum. Instead of focusing on the cascade of psinons (mobile spinons) which this hole has knocked out of the vacuum, we focus on the hole itself, which has properties commonly attributed to antiparticles. The psinon $`(\psi )`$ and the antipsinon $`(\psi ^{})`$ exist in disjunct parts of the psinon vacuum, namely in the magnon and spinon arrays, respectively. When they meet at the border of the two arrays, they undergo a mutual annihilation, represented by the step from the second row to the top row in Fig. 5.
We could have interpreted the small grey circle as a magnon (spin-1 quasi-particle), but when we do that we must take into account that it then coexists in the magnon vacuum with a macroscopic number of fellow magnons (small black circles). From this perspective, the collective excitation must be viewed as containing a finite density of magnons (for $`N\mathrm{}`$), in which the magnon interaction remains energetically significant for scattering states. The nonzero interaction energy obscures the role of individual magnons.
On the other hand, when the small grey circle is interpreted as an antipsinon, then it lives in the psinon vacuum, i.e. almost in isolation. The only other particle present is a psinon (large grey circle). In a macroscopic system, the interaction energy in a psinon-antipsinon $`(\psi \psi ^{})`$ scattering state becomes negligible. Therefore, the identity of both quasi-particles is easily recognizable in the spectrum.
The energies versus the wave numbers of the 16 $`\psi \psi ^{}`$ states listed in Table 2 are shown in Fig. 6(a) as large symbols. The four branches from bottom to top pertain to $`m=1,\mathrm{},4`$. Also shown in the same plot are the $`\psi \psi ^{}`$ states for $`N=256`$. The lower boundary of the $`\psi \psi ^{}`$continuum emerging in the limit $`N\mathrm{}`$ touches down to zero frequency at $`q=0`$ and $`q=q_s=\pi /2`$. Between $`q_s`$ and $`\pi `$, it rises monotonically and reaches the value $`EE_G=h`$. A direct observation of the incommensurate soft mode at $`q_s`$ was made in a neutron scattering experiment on Cu(C<sub>6</sub>D<sub>5</sub>COO)$`{}_{2}{}^{}`$3D<sub>2</sub>O (copper benzoate).DHR+97
Figure 6(b) shows the relative integrated intensity of the $`\psi \psi ^{}`$ excitations for various $`N`$ at fixed $`M_z/N=\frac{1}{4}`$. At $`qq_s=\pi /2`$, virtually all spectral weight of $`S_{zz}(q,\omega )`$ originates from $`\psi \psi ^{}`$ fluctuations. An extrapolation of the data points at $`q=\pi /2`$ suggests that the relative $`\psi \psi ^{}`$ spectral weight is in excess of 93%.
At $`qq_s`$ the $`\psi \psi ^{}`$ contribution to the integrated intensity decreases monotonically but stays dominant over more than half the distance to the zone boundary. The width of the $`\psi \psi ^{}`$ continuum vanishes linearly on approach of $`q=\pi `$, and the relative spectral weight more slowly: $`S_{zz}(q)(\pi q)^\gamma `$, $`\gamma 0.3`$. This enhances the observability of the $`\psi \psi ^{}`$ excitations in the narrow energy range near the Brillouin zone in spite of the low absolute intensity. Finite-$`N`$ data for the integrated intensity $`S_{zz}(q)`$ are shown in the inset to Fig. 6(b). This function is peaked at $`q=q_s`$, where the $`\psi \psi ^{}`$ spectral weight is overwhelmingly predominant.
When we lower $`M_z`$, the soft mode at $`q_s`$ moves to the right, the number of $`2m`$-psinon branches that contribute to the $`\psi \psi ^{}`$ continuum shrinks but each branch gains additional states. At $`M_z=1`$ we are left with one 2-psinon branch extending over the interior of the entire Brillouin zone. This branch is equal to the lowest branch of 2-spinon states with dispersion $`ϵ_L(q)`$, Eq. (4). However, even for this case the psinon vacuum is different from the spinon vacuum. The former is the lowest-energy 2-spinon state (with $`M_z=1`$), whereas the latter is a state with $`M_z=0`$. The wave number of the two vacua differ by $`\pi `$. At $`M_z=0`$ the $`\psi \psi ^{}`$ excitations disappear altogether. The limit $`h0`$ of the infinite chain is very subtle and will be discussed elsewhere.note3
When we increase $`M_z`$ toward the saturation value, the soft mode moves to the left, and the number of $`2m`$-psinon branches increases, but each branch becomes shorter. At $`M_z=N/21`$, the two-parameter set collapses into a one-parameter set consisting of one $`2m`$-psinon state each for $`m=1,2,\mathrm{},N/21`$. These states are more naturally interpreted as a branch of 1-magnon excitations with dispersion $`ϵ_1(q)=J(1\mathrm{cos}q)`$. Their relative spectral weight in $`S_{zz}(q,\omega )`$ is now 100%, but the absolute intensity for $`q0`$ is only of O($`N^1`$).
To further illustrate the roles of the psinon and the antipsinon as the relevant quasi-particles in the collective excitations dominating the spectral weight in $`S_{zz}(q,\omega )`$, we compare in Fig. 7 the energies between the $`\psi \psi ^{}`$ scattering states for $`N=64`$ and the corresponding (fictitious) free $`\psi \psi ^{}`$ superpositions. The vertical displacement of any $`()`$ from the associated $`(+)`$ reflects the interaction energy between the two quasi-particles. This energy approaches zero for all states of this class as $`N\mathrm{}`$.
The energy-momentum relations of the two quasi-particles can be accurately inferred from $`N=2048`$ data for the spectral thresholds of the $`\psi \psi ^{}`$ states as illustrated in the inset to Fig. 7. The psinon dispersion $`ϵ_\psi (p)`$ is confined to the interval at $`0|p|\pi /4`$ (solid line) and the antipsinon dispersion $`ϵ_\psi ^{}(p)`$ to $`\pi /4|p|3\pi /4`$ (dashed line). The different ranges of momentum which the two quasi-particles are allowed to have correspond to the different regions in Fig. 5 across which the circles pertaining to them can be varied.
The lower boundary of the $`\psi \psi ^{}`$ continuum is defined by collective states in which one of the two particles has zero energy: the psinon for $`0|q|\pi /2`$ and the antipsinon for $`\pi /2|q|\pi `$. The upper boundary consists of three distinct segments.
For $`0q0.3935`$ the highest-energy $`\psi \psi ^{}`$ state is made up of a zero-energy psinon with momentum $`p_\psi =\pi /4`$ and an antipsinon with momentum $`p_\psi ^{}=\pi /4+q`$. Here the shape of the continuum boundary is that of the psinon dispersion. Likewise, for $`3\pi /4q\pi `$, the states along the upper continuum boundary are made up of a maximum-energy antipsinon (with momentum $`p_\psi ^{}=3\pi /4`$ and a psinon with momentum $`p_\psi =3\pi /4+q`$. Here the shape of the continuum boundary is that of the psinon dispersion.
When these two delimiting curves are extended into the middle segment, $`0.3935q3\pi /4`$, they join in a cusp singularity at $`q=\pi /2`$. Here the highest $`\psi \psi ^{}`$ state does not involve any zero-energy quasi-particles. The maximum of $`ϵ_\psi (p_\psi )+ϵ_\psi ^{}(p_\psi ^{})`$ subject to the constraint $`p_\psi +p_\psi ^{}=q`$ does not occur at the endpoint of any quasi-particle dispersion curve. Consequently, the $`\psi \psi ^{}`$ continuum is partially folded about the upper continuum boundary along the middle segment.
## VI Lineshapes
To calculate the lineshapes relevant for fixed-$`q`$ scans in an inelastic neutron scattering experiment from the spectrum and matrix elements obtained via Bethe ansatz, we exploit key properties of transition rates and densities of states of sets of excitations that form two-parameter continua in $`(q,\omega )`$ space for $`N\mathrm{}`$. The $`\psi \psi ^{}`$ transition rates (scaled by $`N`$) form a continuous function $`M_{zz}^{\psi \psi ^{}}(q,\omega )`$ for $`N\mathrm{}`$. The $`\psi \psi ^{}`$ density of states (scaled by $`N^1`$) becomes a continuous function $`D^{\psi \psi ^{}}(q,\omega )`$ for $`N\mathrm{}`$. The $`\psi \psi ^{}`$ spectral-weight distribution is then the product $`S_{zz}^{\psi \psi ^{}}(q,\omega )=D^{\psi \psi ^{}}(q,\omega )M_{zz}^{\psi \psi ^{}}(q,\omega )`$.note4 In the following, we consider three wave numbers at $`M_z=N/4`$.
At $`q=\pi /2`$, the $`\psi \psi ^{}`$ continuum is gapless and the relative $`\psi \psi ^{}`$ spectral weight in $`S_{zz}(q,\omega )`$ has a maximum. The scaled density of $`\psi \psi ^{}`$ states is generated from $`N=2048`$ data of the set of points
$$D^{\psi \psi ^{}}(q,\omega _\nu ^{})\frac{2\pi /N}{\omega _{\nu ^{}+1}\omega _\nu ^{}},$$
(12)
where $`\nu ^{}=m`$ marks the antipsinon quantum number in the $`\psi \psi ^{}`$ continuum and picks the dynamically relevant branch from the set of $`2m`$-psinon states. The psinon quantum number $`\nu `$ is adjusted to keep the wave number $`q`$ of the $`\psi \psi ^{}`$ state fixed. This choice of labels produces an ordered sequence of levels. Starting at $`\omega =0`$, the graph of $`D^{\psi \psi ^{}}(\pi /2,\omega _\nu ^{})`$ rises from a nonzero value very slowly up to near the upper band edge, where it bends into a square-root divergence as shown in Fig. 8(a). The divergence is produced by a maximum of the sequence $`\omega _\nu ^{}`$ at the fold of the $`\psi \psi ^{}`$-continuum.note5
In Fig. 8(b) we show finite-$`N`$ data at $`q=\pi /2`$ for the scaled transition rates
$$M_{zz}^{\psi \psi ^{}}(q,\omega _\nu ^{})N|G|S_q^z|\nu ^{}|^2.$$
(13)
These data compellingly suggest the existence of a smooth function $`M_{zz}^{\psi \psi ^{}}(\pi /2,\omega )`$ for the $`\psi \psi ^{}`$ transition rates in the limit $`N\mathrm{}`$, which further highlights the physical significance of the psinon and the antipsinon as relevant quasi-particles in this situation. The function $`M_{zz}^{\psi \psi ^{}}(\pi /2,\omega )`$ is monotonically decreasing with a divergence at $`\omega =0`$ and a cusp singularity at the upper band edge $`\omega _U1.679J`$.
The product of the transition rate function and the (interpolated) density of states is shown in Fig. 8(c).note8 The curve fitted through the data points represents the $`\psi \psi ^{}`$ lineshape at $`q=\pi /2`$ in $`S_{zz}(q,\omega )`$. Its most distinctive feature is the double peak due to apparent divergences at both band edges.
The divergence at $`\omega =0`$, which is caused by the matrix elements, is a power law, $`\omega ^\alpha `$, with an exponent that is exactly known from field theoretic studies of the Heisenberg model.Hald80 ; BIR87 ; FGM+96 For the situation at hand, the value is $`\alpha =0.4688\mathrm{}`$. The divergence at $`\omega _U`$ is caused by the diverging density of states but is weakened if the cusp singularity of $`M_{zz}^{\psi \psi ^{}}(\pi /2,\omega )`$ starts from zero at $`\omega =\omega _U`$. The expectation is a power-law singularity, $`(\omega _U\omega )^\beta `$ with an exponent $`0\beta \frac{1}{2}`$.
It is interesting to compare the $`\psi \psi ^{}`$ transition rate function $`M_{zz}^{\psi \psi ^{}}(\pi /2,\omega )`$ at $`M_z=N/4`$ inferred from the Bethe ansatz with the 2-spinon transition rate function $`M_{zz}^{(2)}(\pi ,\omega )`$ at $`M_z=0`$ calculated via algebraic analysis.KMB+97 The shape of both functions is similar, but there are some differences: $`M_{zz}^{(2)}(\pi ,\omega )`$ has a stronger power-law divergence at $`\omega =0`$ and it approaches zero more rapidly at the upper band edge. As a result it produces a monotonically decreasing spectral-weight distribution $`S_{zz}^{(2)}(\pi ,\omega )`$ (see Fig. 1) notwithstanding the fact that the 2-spinon density of states is also a monotonically increasing function terminating in a square-root divergence.
At $`q=\pi /4`$ the integrated intensity $`S_{zz}(q)`$ is only a third of what it was at $`q=\pi /2`$, but spread over a narrower range of frequencies (see Fig. 6). The bandwidth has shrunk to less than a third of the value it had at $`q=\pi /2`$. The relative $`\psi \psi ^{}`$ contribution to the intensity is even larger than at $`q=\pi /2`$, almost 100%. In this application, the method of analysis is stretched more closely to its limits because $`q=\pi /4`$ exists in fewer manageable system sizes. However, the data still make reliable lineshape predictions possible.
The density of states $`D_{zz}^{\psi \psi ^{}}(\pi /4,\omega )`$, plotted in Fig. 9(a), rises discontinuously from zero to a finite value at the spectral threshold, $`\mathrm{\Delta }E0.379J`$. From there it increases gradually with gradually increasing slope and ends in a cusp singularity at the upper band edge.note6 The finite-$`N`$ data for the scaled transition rates shown in Fig. 9(b) again suggest a smooth $`\omega `$-dependence in the form of a monotonically decreasing curve with enhanced steepness near both band edges. However, the countertrend of the density of states at the upper band edge is of sufficient strength to produce a second maximum in the line shape again.
Also shown in Fig. 9 are the corresponding data for the $`\psi \psi ^{}`$ density of states, transition rates, and lineshape at $`q=3\pi /4`$. Here the relative spectral weight carried by the $`\psi \psi ^{}`$ excitations is only 83% of the value at $`q=\pi /2`$, but that fraction is concentrated over a frequency band that has shrunk to 65% of the width at $`q=\pi /2`$, while the absolute intensity remains fairly high (87% of the value at $`q=\pi /2`$). Both quantities, which determine the $`\psi \psi ^{}`$ lineshape, exhibit similar frequency dependences as we have already observed for the other two fixed-$`q`$ scans. The density of states is divergent again at the upper boundary. The energy gap is now much larger, $`\mathrm{\Delta }E0.899J`$. The fact that the lower continuum boundary at $`q=3\pi /4`$ coincides with the upper continuum boundary at $`q=\pi /4`$ is a consequence of the quasi-particle dispersions as discussed previously.
## VII Conclusion
The spectrum of the completely integrable 1D $`s=\frac{1}{2}`$ Heisenberg antiferromagnet (1) can be generated in more than one way from multiple excitations of quasi-particles. The external magnetic field controls the nature of the ground state. In strong fields, it becomes the vacuum of magnons and in zero field the vacuum of spinons. The dynamically relevant collective excitations of specific quantum fluctuations in the two cases are then naturally described as composites of quasi-particles from the respective species and are likely to involve only a small number of quasi-particles.
In intermediate magnetic fields, neither the magnons nor the spinons provide a useful interpretation of dynamically relevant collective excitations for the same fluctuation operators. The ground state itself contains a macroscopic number of quasi-particles from one or the other of the two species. However, when it is reconfigured as the physical vacuum for psinons and antipsinons, then it turns out that the spin fluctuation operator $`S_q^z`$ induces predominantly transitions to $`\psi \psi ^{}`$ states, which contain just one particle from each kind.
Like the magnon and the spinon, the psinon and the antipsinon are interacting quasi-particles in the Heisenberg model (1). In the $`\psi \psi ^{}`$ scattering states, the interaction energy of the psinon and the antipsinon is of order O($`N^1`$) whereas the interaction energy among magnons or spinons is of order O(1). Hence, for $`N\mathrm{}`$, the $`\psi \psi ^{}`$ states join up in $`(q,\omega )`$-space to form a two-parameter continuum whose spectral boundaries and density of states are fully determined by the energy-momentum relations of the psinon and the antipsinon. Moreover, the scaled $`\psi \psi ^{}`$ transition rates converge for $`N\mathrm{}`$ toward a smooth function of $`q`$ and $`\omega `$.
We have exploited these asymptotic quasi-particle properties to extract lineshape information for the dynamic structure factor $`S_{zz}(q,\omega )`$, which probes the spin fluctuations parallel to the applied magnetic field. The same quasi-particles will also play a dominant role in the spin fluctuations perpendicular to the field, but here different combinations of them make up the composition of the dynamically relevant collective excitations. In the dynamic spin structure factor $`S_+(q,\omega )`$, for example, the spectral weight is almost completely carried by 2-psinon excitations.note3
In all likelihood, the psinon quasi-particles will also be useful for the analysis of thermal spin fluctuations in this model system. The peculiar spectral weight distributions found in recent complete diagonalization studiesFLS97 ; FL98 of $`S_{zz}(q,\omega )`$ at $`h=0`$ and $`T>0`$, for example, indicate the presence of stringent selection rules between collective states coupled by the spin fluctuation operator $`S_q^z`$. In zero field, psinon vacua are densely spread across the entire energy range of the model. Each psinon vacuum can be used as the reference state of a $`2m`$-psinon expansion (9). If there are general selection rules related to psinon quasi-particles among transition rates $`|\lambda ^{}|S_q^z|\lambda |^2`$ within a given class $`K_r`$ of Bethe ansatz solutions, they will have a strong impact on the spectral weight distribution in $`S_{zz}(q,\omega )`$ at all temperatures.
###### Acknowledgements.
We thank Andreas Klümper, Klaus Fabricius, and Alexander Meyerovich for interesting and useful discussions. Financial support from the URI Research Office (for G.M.) and from the DFG Schwerpunkt Kollektive Quantenzustände in elektronischen 1D Übergangsmetallverbindungen (for M.K.) is gratefully acknowledged.
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# Untitled Document
Spacetime Measurements in Kaluza-Klein Gravity
Hongya Liu
Department of Physics, Dalian University of Technology, Dalian, 116024, P.R. China
Bahram Mashhoon
Department of Physics and Astronomy, University of Missouri-Columbia
Columbia, Missouri 65211, USA
## Abstract
We extend the classical general relativistic theory of measurement to include the possibility of existence of higher dimensions. The intrusion of these dimensions in the spacetime interval implies that the inertial mass of a particle in general varies along its worldline if the observations are analyzed assuming the existence of only the four spacetime dimensions. The variations of mass and spin are explored in a simple 5D Kaluza-Klein model.
* PACS numbers: 04.20.Cv, 04.50.+h
1 Introduction
The most basic measurements of a physical observer are those of time and space. In principle, such observations may involve an atomic clock for local temporal measurements and three independent spatial axes for the characterization of space in the neighborhood of the observer. In standard general relativity, the observer carries an orthonormal tetrad frame along its worldline and physical observables are scalars that are obtained as the projections of tensors that correspond to various physical variables upon the local tetrad frame of the observer. In general relativity, just as in Newtonian physics, the observer can determine, via local measurements, the acceleration of its local frame. This acceleration could be in the form of the translational acceleration of the observer as well as the rotation of its local spatial frame. Theoretically, a set of three ideal orthogonal torque-free gyroscopes can provide a nonrotating (i.e. Fermi-Walker transported) spatial frame along the path of the observer. Thus in 4D spacetime a free test observer can carry a nonrotating orthonormal frame that is parallel transported along its geodesic worldline. In view of the possibility of existence of higher dimensions, it is worthwhile to examine how the theory of measurement in general relativity would have to be extended if higher dimensions intrude into the spacetime arena. This intrusion is expected in any realistic higher-dimensional physical theory ; nevertheless, the interpretation of observational data currently involves only the standard four spacetime dimensions.
In this paper, we explore the extension of the classical general relativistic theory of measurement to the Kaluza-Klein theory by studying some of the main observational consequences of the dependence of the spacetime metric upon extra dimensions. For instance, we show in section 2 that the mass of a test particle in general varies along its worldline if the motion of matter is not wholly confined to the four spacetime dimensions. This circumstance is expected in realistic higher-dimensional theories . Our physical considerations are motivated by the fact that experimental data are routinely analyzed assuming the existence of only the four spacetime dimensions. In section 3, we show that an initially orthonormal frame does not in general remain orthonormal along the worldline of a test observer. To render these results explicit, we consider a concrete 5D model in section 4 and explore its physical consequences. Section 5 contains a brief discussion of our results.
2 Variation of the Inertial Mass
Imagine that the universe can be described in terms of $`4+N`$ dimensions with $`N1`$ and the 4D spacetime part, which is embedded in the $`4+N`$ manifold, has a metric of the form
$$ds^2=g_{\mu \nu }(x,y)dx^\mu dx^\nu .$$
(1)
Here $`x`$ stands for the spacetime coordinates and $`y`$ stands for the extra dimensions $`(y^1,\mathrm{},y^N)`$ that are in general reflected in the spacetime metric. Greek indices run from 0 to 3. The path of a test particle in the $`4+N`$ dimensional manifold involves the 4D velocity $`u^\alpha =dx^\alpha /ds`$, where $`s`$ is the proper time along the path such that
$$g_{\mu \nu }(x,y)u^\mu u^\nu =1.$$
(2)
Differentiating (2) with respect to $`s`$, we find that
$$g_{\mu \nu ,\alpha }u^\alpha u^\mu u^\nu +\underset{i=1}{\overset{N}{}}\frac{g_{\mu \nu }}{y^i}\frac{dy^i}{ds}u^\mu u^\nu +2g_{\mu \nu }u^\mu \frac{du^\nu }{ds}=0.$$
(3)
This relation may be written in the form
$$\left(g_{\mu \nu ,\alpha }+g_{\mu \alpha ,\nu }g_{\nu \alpha ,\mu }\right)u^\alpha u^\mu u^\nu +2g_{\mu \nu }u^\mu \frac{du^\nu }{ds}=\underset{i=1}{\overset{N}{}}\frac{g_{\mu \nu }}{y^i}\frac{dy^i}{ds}u^\mu u^\nu .$$
(4)
Using the fact that the 4D connection is given by
$$\mathrm{\Gamma }_{\alpha \beta }^\nu (x,y)=\frac{1}{2}g^{\nu \eta }\left(g_{\eta \alpha ,\beta }+g_{\eta \beta ,\alpha }g_{\alpha \beta ,\eta }\right),$$
(5)
we can write equation (4) as
$$2g_{\mu \nu }u^\mu \left(\frac{du^\nu }{ds}+\mathrm{\Gamma }_{\alpha \beta }^\nu u^\alpha u^\beta \right)=\underset{i=1}{\overset{N}{}}\frac{g_{\mu \nu }}{y^i}\frac{dy^i}{ds}u^\mu u^\nu .$$
(6)
The 4D acceleration of the particle is defined by
$$A^\mu =\frac{du^\mu }{ds}+\mathrm{\Gamma }_{\alpha \beta }^\mu u^\alpha u^\beta ,$$
(7)
so that equation (6) implies
$$u_\mu A^\mu =\frac{1}{2}\underset{i=1}{\overset{N}{}}\frac{g_{\mu \nu }}{y^i}\frac{dy^i}{ds}u^\mu u^\nu .$$
(8)
If $`g_{\mu \nu }`$ does not depend on the extra dimensions, then equation (8) becomes $`u_\mu A^\mu =0`$, as expected. However, in the higher-dimensional theories $`g_{\mu \nu }`$ may depend on $`y`$ and $`y`$ may vary with respect to $`s`$. So the right-hand side of equation (8) may not vanish. Let us note that $`u^\mu `$ is a timelike vector, so $`u_\mu A^\mu 0`$ indicates that there may be a timelike component of acceleration in $`A^\mu `$ in higher-dimensional theories. This is an extraordinary result, since all known basic 4D forces are spacelike and lead to accelerations that are orthogonal to the 4D velocity of the particle . It turns out that the most natural way to incorporate a timelike acceleration into 4D physics is to assume that the “invariant” inertial mass of the test particle varies along its worldline.
Experimental data are reduced and interpreted at present assuming that the most general force law for the motion of a test particle may be written classically as
$$\frac{Dp^\mu }{ds}\frac{dp^\mu }{ds}+\mathrm{\Gamma }_{\alpha \beta }^\mu u^\alpha p^\beta =F^\mu ,$$
(9)
where $`p^\mu mu^\mu `$ is the momentum of the particle and $`F^\mu `$ consists of all forces acting on the particle arising from the known fundamental interactions. In the rest frame of the particle, we expect all forces acting on the particle to be 3D vectors; therefore, $`u_\mu F^\mu =0`$. This relation together with the force law (9), $`Dp^\mu /ds=\left(dm/ds\right)u^\mu +mA^\mu =F^\mu `$, implies that
$$\frac{1}{m}\frac{dm}{ds}=\frac{1}{2}\underset{i=1}{\overset{N}{}}\frac{g_{\mu \nu }}{y^i}\frac{dy^i}{ds}u^\mu u^\nu ,$$
(10)
where we have used equation (8). That is, the simplest interpretation of equation (8) in terms of 4D physics is to assume that the invariant “rest” mass of the particle may vary with respect to its proper time $`s`$ due to the existence of higher dimensions. Conversely, the observation of such a basic variation would indicate the presence of an extra timelike acceleration and this could come about precisely because of the intrusion of the extra dimensions into the 4D physics as indicated by equations (1) and (8).
Let us note that equation (10) may be expressed as
$$\delta \left(m^2\right)=\left[\underset{i=1}{\overset{N}{}}\frac{g_{\mu \nu }(x,y)}{y^i}\delta y^i\right]p^\mu p^\nu ,$$
(11)
where
$$m^2=g_{\mu \nu }(x,y)p^\mu p^\nu .$$
(12)
Equations (11) and (12) can be used to determine the variable inertial mass in higher-dimensional theories.
The acceleration $`A^\mu `$ and the variation of the extra coordinates of the particle along its worldline are determined by the equation of motion of the theory. In principle, it is possible that $`y`$ can vary in just such a way as to render $`dm/ds=0`$ in equation (10). This would, of course, require rather special circumstances; therefore, we pursue in this paper the general situation in which $`dm/ds0`$.
3 Variation of the Spin
Let us first consider an ideal gyroscope represented by the spin vector $`\sigma ^\mu `$ within the context of classical general relativity. We may imagine in our classical model that we are dealing here with the limiting case of an ideal gyroscope with its magnitude of spin given by $`\sigma =I\omega `$ in a certain “rest” frame of the system. Here $`I`$ is the proper moment of inertia of the gyroscope and $`\omega `$ is its angular speed of rotation with respect to the proper time of the “pointlike” spinning particle. The general relativistic theory is based on the Mathisson-Papapetrou equations for a “pole-dipole” test particle. It follows from detailed considerations of the classical theory of such ideal spinning “point” particles that $`\sigma _\mu u^\mu =0`$. Moreover, $`\sigma ^\mu `$ is nonrotating, i.e.
$$\frac{D\sigma ^\mu }{ds}=(\sigma _\alpha A^\alpha )u^\mu ,$$
(13)
so that $`\sigma _\alpha u^\alpha `$ and $`\sigma ^2=g_{\mu \nu }\sigma ^\mu \sigma ^\nu `$ are constants along the worldline of the particle.
Consider next the possibility that the spacetime metric could depend upon higher dimensions. In this case, one can extend the treatment of section 2 to demonstrate that
$$Dg_{\mu \nu }=\underset{i=1}{\overset{N}{}}\frac{g_{\mu \nu }}{y^i}dy^i.$$
(14)
Using this result, it is straightforward to show that
$$\frac{d}{ds}(\sigma _\mu u^\mu )=T_\mu u^\mu +A_\mu \sigma ^\mu +\underset{i=1}{\overset{N}{}}\frac{g_{\mu \nu }}{y^i}\frac{dyi}{ds}\sigma ^\mu u^\nu $$
(15)
and
$$\frac{d}{ds}(\sigma ^2)=2T_\mu \sigma ^\mu \underset{i=1}{\overset{N}{}}\frac{g_{\mu \nu }}{y^i}\frac{dyi}{ds}\sigma ^\mu \sigma ^\nu ,$$
(16)
where $`T^\mu =D\sigma ^\mu /ds`$ is the torque and $`A^\mu `$ is the translational acceleration (7) of the particle. The vectors $`A^\mu `$ and $`T^\mu `$ must be determined from the higher-dimensional theory under consideration. It follows from the extension of these results to the axes of a spatial frame carried by an observer that an initially orthonormal frame will not in general remain orthonormal in the course of time. In particular, the magnitude of spin can change with time. To illustrate how this could come about, we now turn to a simple Kaluza-Klein gravitational model.
4 A 5D Model
We consider a Kaluza-Klein model with one extra noncompactified spacelike dimension $`y`$ such that the 5D metric is given by
$$dS^2=\widehat{g}_{AB}dx^Adx^B=ds^2dy^2,$$
(17)
where $`x^A=(x^\mu ,y)`$ and $`\widehat{g}_{AB}`$ satisfies the 5D vacuum field equations $`\widehat{R}_{AB}=0`$. It is natural to assume in this theory that a test particle follows the 5D geodesic equation
$$\frac{d^2x^A}{dS^2}+\widehat{\mathrm{\Gamma }}_{BC}^A\frac{dx^B}{dS}\frac{dx^C}{dS}=0,$$
(18)
so that the 5D velocity vector $`U^A=dx^A/dS`$ is parallel transported along the path. Moreover, we expect that the 5D spin vector $`\mathrm{\Sigma }^A`$ of an ideal test gyroscope would also be parallel transported along its path
$$\frac{d\mathrm{\Sigma }^A}{dS}+\widehat{\mathrm{\Gamma }}_{BC}^AU^B\mathrm{\Sigma }^C=0.$$
(19)
To connect equations (18)-(19) with our spacetime variables in previous sections, we note that
$$u^\mu =\left(\frac{dS}{ds}\right)U^\mu ,\sigma ^\mu =\left(\frac{dS}{ds}\right)\mathrm{\Sigma }^\mu .$$
(20)
The relationship between 4D and 5D velocities follows simply from the definitions, while the corresponding relation for spin is expected by analogy.
The Kaluza-Klein field equations $`\widehat{R}_{AB}=0`$ can be reduced to certain constraint equations together with the 4D gravitational field equations of standard general relativity with an effective energy-momentum tensor as the source of the gravitational field. This is consistent with recent investigations based on Campbell’s theorem that an $`n`$-dimensional Riemannian space can be locally embedded in a Ricci-flat $`(n+1)`$-dimensional Riemannian space .
To proceed further, we need an explicit solution of the field equations. It is possible to show that with
$$g_{\mu \nu }(x,y)=\frac{y^2}{L^2}\stackrel{~}{g}_{\mu \nu }(x)$$
(21)
in the spacetime interval (1), $`\widehat{g}_{AB}`$ in equation (17) is a solution of $`\widehat{R}_{AB}=0`$, provided that $`\stackrel{~}{g}_{\mu \nu }(x)`$ is any source-free solution of general relativity with a cosmological constant $`\stackrel{~}{\mathrm{\Lambda }}=3/L^2`$. Here $`L`$ is a constant length. If $`\stackrel{~}{g}_{\mu \nu }(x)`$ is the de Sitter solution, then $`\widehat{R}_{ABCD}=0`$ and the 5D metric is flat. However, for more complicated Einstein spaces such as the Kerr-de Sitter solution the 5D manifold is curved. Let us note that for $`g_{\mu \nu }(x,y)`$, the 4D Ricci tensor is then given by $`R_{\alpha \beta }=3y^2g_{\alpha \beta }`$. Thus the spacetime metric can be interpreted as a source-free solution of general relativity with a cosmological “constant” $`\mathrm{\Lambda }=3/y^2`$. Various aspects of the gravitational model under consideration here have been explored in a number of publications . A recent discussion of variable cosmological “constant” is contained in .
With a spacetime metric of the form (21), equation (18) reduces to $`Du^\alpha /ds=A^\alpha `$ with
$$A^\alpha =\frac{1}{y}\frac{dy}{ds}u^\alpha $$
(22)
and
$$y\frac{d^2y}{ds^2}\left(\frac{dy}{ds}\right)^2+1=0.$$
(23)
Equation (10) — or, equivalently, equation (22) — implies that $`m=\lambda y`$, where $`\lambda `$ is a constant. Moreover, equation (23) can be solved by studying the variation of $`dy/ds`$ with respect to $`y`$. The result is
$$\left(\frac{dy}{ds}\right)^2+\frac{K}{L_0^2}y^2=1,$$
(24)
where $`K=0,\pm 1`$ and $`L_0`$ is a constant length. It follows that
$$\pm y=[\begin{array}{cc}L_0\mathrm{sin}\left(\frac{ss_0}{L_0}\right)\hfill & \mathrm{for}K=+1,\hfill \\ ss_0\hfill & \mathrm{for}K=0,\hfill \\ L_0\mathrm{sinh}\left(\frac{ss_0}{L_0}\right)\hfill & \mathrm{for}K=1.\hfill \end{array}$$
(25)
If we now substitute for $`y`$ from equation (25) in $`dS^2=ds^2dy^2`$, we find that $`dS^2`$ is greater than, equal to, or less than zero for $`K=+1,0,1`$, respectively. In general relativity, we know that the worldline of a massive (massless) particle should have $`ds^2>0`$ ($`ds^2=0`$). Extending this requirement from 4D to 5D, we choose $`dS^2>0`$ ($`dS^2=0`$) for the motion of a massive (massless) particle. This implies that $`K=+1`$ and $`K=0`$ in equation (25) hold for massive and massless particles, respectively. It follows from this interpretation that for $`K=0`$ we have $`ds=0`$, so that $`s=s_0`$ and hence $`y=0`$. Let us note that this interpretation is consistent with equation (25), since for $`ss_00`$ the timelike and spacelike cases $`K=\pm 1`$ reduce to the null case $`K=0`$ regardless of the value of $`L_0`$. Moreover, for $`ss_00,y0`$, and hence $`m=\lambda y0`$, so that lightlike propagation occurs only for a massless particle. Thus null rays propagate only in the 4D spacetime part of the 5D manifold.
The case of a massless particle is important for the treatment of light propagation; therefore, it is necessary to show in detail how the null geodesic case comes about as a limiting case of the motion of a massive particle for $`m0`$. This involves a standard limiting procedure in general relativity that will be adapted here to the situation at hand. Let us note that for a massive particle the momentum $`p^\alpha =mu^\alpha `$ is covariantly constant along its 4D worldline. Here $`m=\lambda y`$ and $`y`$ is given by equation (25). Thus we choose a new variable $`z`$ defined by $`ss_0=\epsilon z`$, where $`\epsilon ,0<\epsilon 1`$, is a constant such that $`\epsilon 0`$ in the massless limit. If we now let $`q=\pm \lambda ^1`$ $`\mathrm{ln}z`$ be an affine parameter along the worldline, then it is simple to show that for $`\epsilon 0,p^\alpha k^\alpha =dx^\alpha /dq`$, which is tangent to a null ray, and that the equation of motion reduces in this limit to the null geodesic equation $`Dk^\alpha /dq=0`$.
For a massive particle we have from the first equation in (25)
$$m=m_{\mathrm{max}}\mathrm{sin}\frac{ss_0}{L_0},\frac{1}{m}\frac{dm}{ds}=\frac{1}{L_0}\mathrm{cot}\frac{ss_0}{L_0},$$
(26)
where $`m_{\mathrm{max}}=\pm \lambda L_0`$. The present upper limit on $`|m^1(dm/ds)|`$ is of order $`10^{12}`$/yr, so that $`L_0`$ must be a sufficiently large cosmic length to render equation (26) compatible with observation .
Let us now explore the implications of equation (19) for the motion of the gyro axis. We find that $`D\sigma ^\mu /ds=T^\mu `$, where the torque is given by
$$T^\mu =\frac{1}{y}\left(\frac{dS}{ds}\right)\mathrm{\Sigma }^4u^\mu $$
(27)
and the variation of $`\mathrm{\Sigma }^4`$ is governed by
$$\frac{d\mathrm{\Sigma }^4}{ds}+\frac{1}{y}\left(\frac{ds}{dS}\right)\sigma _\alpha u^\alpha =0.$$
(28)
These equations are consistent with the fact that $`\mathrm{\Sigma }_AU^A`$ and $`\mathrm{\Sigma }_A\mathrm{\Sigma }^A`$ are constants along the path. For instance,
$$\left(\frac{ds}{dS}\right)^2\sigma ^2+\left(\mathrm{\Sigma }^4\right)^2=\mathrm{\Sigma }_0^2,$$
(29)
where $`\mathrm{\Sigma }_0=(\mathrm{\Sigma }_A\mathrm{\Sigma }^A)^{1/2}`$ is the constant magnitude of the 5D spin vector. One can now compute the variation of $`\sigma _\alpha u^\alpha `$ and $`\sigma ^2`$ along the path using equations (15) and (16). In particular, we find that $`\sigma _\alpha u^\alpha =C_0+C_1\mathrm{cos}\theta `$ and
$$\mathrm{\Sigma }^4=\xi (C_0\mathrm{cot}\theta +C_1\mathrm{csc}\theta ),$$
(30)
where $`\theta =(ss_0)/L_0`$ and $`\xi ^2=1`$ (i.e. $`\xi `$ is either $`+1`$ or $`1`$). Here $`C_0`$ and $`C_1`$ are constants of integration. To simplify matters, let us assume that $`C_0=\mathrm{\Sigma }_AU^A`$ vanishes along the path, generalizing the standard 4D constraint for a “pointlike” gyroscope. Moreover, we set $`C_1=0`$; then, $`T^\mu =0`$ and $`\sigma _\alpha u^\alpha =0`$. It follows that one can choose initial conditions such that $`\sigma ^\mu `$ is nonrotating along the path, since $`\sigma _\mu A^\mu =0`$ follows from our assumptions. However, the magnitude of spin would still vary in accordance with equation (29), i.e. $`\sigma =\lambda ^{}y`$, where $`\lambda ^{}0`$ is a constant. Hence the “pole-dipole” particle’s mass and spin both vary along its trajectory in such a way that $`\sigma /m`$ remains constant.
5 Discussion
A preliminary analysis of the basic spacetime measurements in higher-dimensional gravity theory reveals that the inertial mass and spin of an ideal classical “pointlike” gyroscope may vary along its worldline. These results could be significant in Kaluza-Klein cosmology as well as the search for extra dimensions.
Our discussion has been primarily concerned with classical general relativity; however, the results are expected to be of more general validity. In fact, the intrusion of the extra dimensions in the spacetime domain implies that it is not possible in general to reduce the spacetime metric to the Minkowski form at an arbitrary event in spacetime. This would indicate a breakdown of some of the basic concepts of standard relativistic physics; for instance, the inertial mass and the magnitude of spin are invariant constants that characterize the irreducible unitary representations of the Poincaré group, but could now become variables. We have explored this circumstance in this paper within the eikonal approximation of the classical theory of gravitation; nevertheless, our treatment could be extended to other physical quantities such as the variation of the phase of a wave.
H. L. acknowledges the support of the National Natural Science Foundation of China (grant no. 19975007).
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C. Romero, R. Tavakol and R. Zalatdinov, Gen. Rel. Grav. 28, 365 (1996); J. E. Lidsey, C. Romero, R. Tavakol and S. Rippl, Class. Quantum Grav. 14, 865 (1997).
J.E. Campbell, A Course of Differential Geometry (Clarendon Press, Oxford, 1926).
B. Mashhoon, H. Liu and P. S. Wesson, Phys. Lett. B 331, 305 (1994); H. Liu and B. Mashhoon, Ann. Physik 4, 565 (1995); B. Mashhoon, H. Liu and P.S. Wesson, in Proc. Seventh Marcel Grossman Meeting on General Relativity (Stanford, 1994), edited by R. T. Jantzen and G. Mac Keiser (World Scientific, Singapore, 1996), p. 333; P. S. Wesson, B. Mashhoon and Hongya Liu, Mod. Phys. Lett. A 12, 2309 (1997).
J. M. Overduin and F. I. Cooperstock, Phys. Rev. D 58, 043506 (1998).
I. I. Shapiro, in Quantum Gravity and Beyond — Essays in Honor of Louis Witten, edited by F. Mansouri and J. J. Scanio (World Scientific, Singapore, 1993), p. 180; J. D. Bekenstein, Phys. Rev. D 15, 1458 (1977).
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# Imaging the Phase of an Evolving Bose-Einstein Condensate Wavefunction
## Abstract
We demonstrate a spatially resolved autocorrelation measurement with a Bose-Einstein condensate (BEC) and measure the evolution of the spatial profile of its quantum mechanical phase. Upon release of the BEC from the magnetic trap, its phase develops a form that we measure to be quadratic in the spatial coordinate. Our experiments also reveal the effects of the repulsive interaction between two overlapping BEC wavepackets and we measure the small momentum they impart to each other.
A trapped Bose-Einstein condensate has unique value as a source for atom lasers and matter-wave interferometry because its atoms occupy the same quantum state, with uniform spatial phase. However, when released from the trapping potential, a BEC with repulsive atom-atom interactions expands, developing a non-uniform phase profile. Understanding this phase evolution will be important for applications of coherent matter waves. We have developed a new interferometric technique using spatially resolved autocorrelation to measure the functional form and time evolution of the phase of a BEC wavepacket expanding under the influence of its mean field repulsion.
In 1997, the coherence of weakly interacting BECs was demonstrated by releasing two spatially separated condensates and observing their interference . Subsequent experiments have further investigated condensate coherence properties. One used velocity-resolved Bragg diffraction to probe the momentum spectrum of trapped and released BECs. A complementary experiment that used matter-wave interferometry can be interpreted as a measurement of the spatial correlation function, whose Fourier transform is the momentum spectrum. These experiments showed that a trapped condensate has a uniform phase, and a released condensate develops a non-uniform phase profile. (Recently the influence of non-zero temperature on coherence properties was also investigated ). The experiments reported in this Letter combine spatial resolution and interferometry to measure the functional form of the time-dependent phase profile of a released condensate. We also make the first measurement of the velocity imparted to two equal BEC wavepackets from their mutual mean-field repulsion .
We perform our experiments with a condensate of $`1.8(4)\times 10^6`$ sodium atoms in the $`3S_{1/2}`$, $`F=1`$, $`m_F=1`$ state. The sample has no discernable non-condensed (i.e. thermal) component. The condensate is prepared following the method of Ref. and is held in a magnetic trap with trapping frequencies $`\omega _x=\sqrt{2}\omega _y=2\omega _z=2\pi \times `$27 Hz. Using a scattering length of $`a=2.8`$ nm, the calculated Thomas-Fermi diameters are 47 $`\mu `$m, 66 $`\mu `$m, and 94 $`\mu `$m, respectively.
We release the BEC from the magnetic trap and it expands, driven mostly by the mean-field repulsion of the atoms. This expansion implies the development of a nonuniform spatial phase profile (recall that the velocity field is proportional to the gradient of the quantum phase). After an expansion time $`T_0`$, we probe the phase profile with matter-wave Bragg interferometry . Our interferometer splits the BEC into two wavepackets and recombines them with a chosen overlap, producing interference fringes, which we measure with absorption imaging . From the dependence of the fringe spacing on the overlap, we extract the phase profile of the wavepackets.
Our atom interferometer consists of three optically-induced Bragg-diffraction pulses applied successively in time (Fig. 1). Each pulse consists of two counter-propagating laser beams whose frequencies differ by 100 kHz. They are detuned by about $`2`$ GHz from atomic resonance ($`\lambda =2\pi /k=589`$ nm) so that spontaneous emission is negligible. The first pulse has a duration of 6 $`\mu `$s and intensity sufficient to provide a $`\pi /2`$ pulse, which coherently splits the BEC into two wavepackets, $`\psi _A`$ and $`\psi _B`$. The wavepackets have about the same number of atoms and only differ in their momenta: $`p=0`$ and $`p=2\mathrm{}k`$. At a time $`T_1=1`$ ms after the first Bragg pulse, the two wavepackets are completely separated and a second Bragg pulse (a $`\pi `$ pulse) of 12 $`\mu `$s duration transfers $`\psi _B`$ to a state with $`p0`$ and $`\psi _A`$ to $`p2\mathrm{}k`$ . After a variable time $`T_2`$ the wavepackets partially overlap again and we apply a third pulse, of 6 $`\mu `$s duration (a $`\pi /2`$ pulse). This last pulse splits each wavepacket into the two momentum states. The interference of the overlapping wavepackets in each of the two momentum states allows the determination of the local phase difference between them. By changing the time $`T_2`$ we vary $`\delta x=x_Ax_B`$, the separation of $`\psi _A`$ and $`\psi _B`$ at the time of the final Bragg pulse. The set of data at different $`\delta x`$ constitutes a new type of spatial autocorrelation measurement that is similar to the “FROG” technique used to measure the complete field of ultrafast laser pulses. From these measurements we obtain the phase profile of the wavepackets in the $`x`$ direction.
Figure 2a-e shows one interferometer output port for different $`\delta x`$ (different $`T_2`$) after an expansion time $`T_0`$ = 4 ms. In general, we observe straight, evenly spaced fringes (although for small $`T_0`$ and $`T_2`$ the fringes may be somewhat curved). There is a value of $`\delta x=x_00`$ where we observe no fringes (Fig. 2c) and the fringe spacing decreases as $`|\delta xx_0|`$ increases. Figure 2f, a cut through Fig. 2d, shows the high-contrast fringes . Our data analysis uses the average fringe period $`d`$, obtained from plots like Fig. 2f.
The fringes come from two different effects: the interference of two wavepackets with quadratic phase profile, and a relative velocity between the wavepackets’ centers. The data can be understood by calculating the fringe spacing along $`x`$ at output port 1 . We assume that the phase $`\varphi `$ of the wavefunction $`f\mathrm{e}^{i\varphi }`$ can be written as $`\varphi =\frac{\alpha }{2}x^2+\beta x`$. The equal spacing of the fringes implies, as predicted in the Thomas-Fermi limit , that $`\varphi `$ has no significant higher-order terms . The curvature coefficient $`\alpha `$ describes the mean-field expansion of the wavepackets and $`\beta `$ describes a relative repulsion velocity. The velocity arises because the wavepackets experience a repulsive push as they first separate and again as they recombine. The density at port 1 (see Fig. 1) just after the final interferometer pulse is the interference pattern $`|\psi _{A1}+\psi _{B1}|^2`$ of the wavepackets $`\psi _{A1}`$ and $`\psi _{B1}`$:
$$|f(x\delta x)e^{i(\frac{\alpha }{2}(x\delta x)^2\beta (x\delta x))}+f(x)e^{i(\frac{\alpha }{2}x^2+\beta x)}|^2,$$
(1)
where we assume that the amplitudes and curvatures of the wavepackets are equal and their velocities have equal magnitude and opposite direction. The cross term of (1) is
$$2f(x\delta x)f(x)\mathrm{cos}\left[\left(\alpha \delta x+\frac{M\delta v}{\mathrm{}}\right)x+C\right],$$
(2)
where $`M`$ is the sodium mass, $`M\delta v/\mathrm{}2\beta `$, and $`C`$ is independent of $`x`$ . $`\delta v=v_Bv_A`$ is the relative repulsion velocity between the wavepackets $`\psi _{A1}`$ and $`\psi _{B1}`$. Expression (2) predicts fringes with spatial frequency,
$$\kappa =\alpha \delta x+\frac{M\delta v}{\mathrm{}},$$
(3)
where $`|\kappa |=2\pi /d`$. When there are no fringes, $`\kappa `$ = 0 and the wavepacket separation $`\delta x=x_0M\delta v/\alpha \mathrm{}`$.
Figure 3 plots the measured $`\kappa `$ vs. $`\delta x`$ for $`T_0`$ = 1 and 4 ms. The data are well fit by a straight line as expected from Eq. (3) in the approximation that $`\alpha `$ and $`\delta v`$ are independent of $`\delta x`$. The slopes of the lines are the phase curvatures $`\alpha `$, and the $`\kappa `$ intercepts give the relative velocities $`\delta v`$.
We checked the validity of the data analysis procedure by analyzing data simulated with a 1-D Gross-Pitaevskii (GP) treatment. Despite variations of $`\delta v`$ and $`\alpha `$ with $`\delta x`$ (due to their continued evolution during the variable time $`T_2`$), we find that $`\kappa `$ is still linear in $`\delta x`$. The slopes and intercepts in general are averages over the range of $`\delta x`$ used in the experiment.
The interference fringes used to determine $`\alpha `$ and $`\delta v`$ are created at the time of the final interferometer pulse. Because the two outputs overlap at that moment, we wait a time $`T_3`$ for them to separate before imaging. During this time, the wavepackets continue to expand. The 1-D simulations show that the fringe spacings and the wavepackets expand in the same proportion. We correct $`\kappa `$ (by typically 15 $`\%`$) for this, using the calculated expansion from a 3-D solution of the GP equation described below.
The different slopes and intercepts of the two lines in Fig. 3 show that the curvature $`\alpha `$ and relative velocity $`\delta v`$ of the wavepackets depend on the release time $`T_0`$ before the first interferometer pulse. Figure 4 plots the dependence of $`\alpha `$ and $`\delta v`$ on various release times $`T_0`$. The condensate initially has a uniform phase so that immediately after its release from the trap $`\alpha =0`$. We nevertheless measure a nonzero $`\alpha `$ for $`T_0`$ = 0 ms because the BEC expands during $`T_1`$ and $`T_2`$. As a function of time, $`\alpha `$ behaves as $`\dot{D}`$/$`D`$ where $`D`$ is the wavepacket diameter and $`\dot{D}`$ is its rate of change . At early times when the mean-field energy is being converted to kinetic energy, $`\dot{D}`$ increases rapidly, increasing $`\alpha `$. At late times, after the mean-field energy has been converted, $`D`$ increases while $`\dot{D}`$ is nearly constant, decreasing $`\alpha `$.
We predict the time evolution of $`\alpha `$ using the Lagrangian Variational Method (LVM) . The LVM uses trial wavefunctions with time dependent parameters to provide approximate solutions of the 3-D time-dependent GP equation. In the model, the effect of the interferometer pulses is to replace the original wavepacket with a superposition of wavepackets having different momenta; e.g., the action of our first interferometer pulse is $`\psi _0\left(\psi _0+e^{i2kx}\psi _0\right)/\sqrt{2}`$. We use Gaussian trial wavefunctions in the LVM and, for simplicity, neglect the interaction between the wavepackets, to calculate the phase curvature $`\alpha `$ at the time of the last interferometer pulse. This result, with $`T_1=T_2`$, is the solid line of Fig. 4a.
We use energy conservation to calculate the relative repulsion velocity $`\delta v`$ between $`\psi _{A1}`$ and $`\psi _{B1}`$ because we neglect wavepacket interactions in the LVM. In the Thomas-Fermi approximation, we can calculate the amount of energy available for repulsion when $`T_0`$ = 0. A trapped condensate has $`\frac{5}{7}\mu `$ average total energy per particle, where $`\mu `$ is the chemical potential . After release from the trap, it has $`\frac{2}{7}\mu `$ average mean-field energy per particle. Applying a $`\pi /2`$ Bragg pulse to the BEC causes a density corrugation, which increases the mean-field energy to $`\frac{3}{7}\mu `$ per particle. In the approximation that the wavepackets do not deform as they separate and recombine, one can show that 1/3 of the total mean-field energy goes into expansion of the wavepackets, and 2/3 is available for kinetic energy of center-of-mass motion. Therefore $`\frac{2}{7}\mu `$ of mean-field energy per particle is available for repulsion. The corresponding repulsion velocity is only about $`10^2`$ of a photon recoil velocity. The repulsion energy and $`\delta v`$ decrease for larger $`T_0`$ because both are inversely proportional to the condensate volume, which we calculate with the LVM. The two curves shown in Fig. 4b are the calculated $`\delta v`$ when $`\delta x=0`$ (solid curve) and $`\delta v`$ averaged over the different $`\delta x`$ used in the experiment (dashed curve). The 1-D GP simulations suggest that for small $`T_0`$, the results of the experiment should be closer to the solid curve; and for large $`T_0`$, closer to the dashed curve. The data is consistent with this trend.
In a related set of experiments we performed interferometry in the trap. This differs from the experiments on a released BEC because there is no expansion before the first interferometer pulse and the magnetic trap changes the relative velocity of the wavepackets between the interferometer pulses (Fig. 5a). To better reveal the velocity differences, we choose $`T_1`$ = $`T_2=T`$ to suppress fringes arising from the phase curvature. As with the released BEC measurements, we observe equally spaced fringes at the output of the interferometer, although the fringes are almost entirely due to a relative velocity $`v`$ between the wavepackets $`\psi _{A1}`$ and $`\psi _{B1}`$ at the time of the third interferometer pulse. We obtain $`v`$ from the fringe periodicity after a small correction for residual phase curvature .
Two effects contribute to $`v`$: the mutual repulsion between the wavepackets $`\psi _A`$ and $`\psi _B`$ and the different action of the trapping potential on the two wavepackets in the interferometer. The latter effect occurs because after the first Bragg pulse, $`\psi _A`$ remains at the minimum of the magnetic potential while $`\psi _B`$ is displaced. Wavepacket $`\psi _B`$ therefore spends more time away from the center of the trap and experiences more acceleration than $`\psi _A`$.
Following the last Bragg pulse, $`\psi _{A1}`$ and $`\psi _{B1}`$ have a velocity difference which for our parameters can be approximated by $`v\frac{2\mathrm{}k}{M}\mathrm{sin}^2(\omega _xT)+\delta v`$ . Figure 5b plots $`v`$ versus $`T`$, and the curve is a fit to the above expression. We obtain the trap frequency $`\omega _x/2\pi =26.7(15)`$ Hz, in excellent agreement with an independent measurement. We also obtain the relative velocity from the mean-field repulsion $`\delta v=0.49(12)`$ mm/s, which we expect to be somewhat larger than for the released measurements because the wavepackets contract, producing a larger mean field.
In conclusion, we demonstrate an autocorrelating matter-wave interferometer and use it to study the evolution of a BEC phase profile by analyzing spatial images of interference patterns. We study how the phase curvature of the condensate develops in time and measure the repulsion velocity between two BEC wavepackets. Our interferometric method should be useful for characterizing other interesting condensate phase profiles. For example, it can be applied to detect excitations of a BEC with characteristic phase patterns, such as vortices and solitons . The method should be useful for further studies of the interaction of coherent wavepackets and to study the coherence of atom lasers.
We thank T. Busch, D. Feder, and L. Collins for helpful discussions. This work was supported in part by the US Office of Naval Research and NASA. J.D. acknowledges support from the Alexander von Humboldt foundation. M.E. and C.W.C. acknowledge partial support from NSF grant numbers 9802547 and 9803377.
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# Resistivity of a Metal between the Boltzmann Transport Regime and the Anderson Transition
## Abstract
We study the transport properties of a finite three dimensional disordered conductor, for both weak and strong scattering on impurities, employing the real-space Green function technique and related Landauer-type formula. The dirty metal is described by a nearest neighbor tight-binding Hamiltonian with a single s-orbital per site and random on-site potential (Anderson model). We compute exactly the zero-temperature conductance of a finite-size sample placed between two semi-infinite disorder-free leads. The resistivity is found from the coefficient of linear scaling of the disorder-averaged resistance with sample length. This “quantum” resistivity is compared to the semiclassical Boltzmann expression computed in both Born approximation and multiple scattering approximation.
Ever since Anderson’s seminal paper, a prime model for the theories of the disorder induced metal-insulator, or localization-delocalization (LD), transition in non-interacting electron systems has been the tight-binding Hamiltonian (TBH) on a hypercubic lattice
$$\widehat{H}=\underset{𝐦}{}\epsilon _𝐦|𝐦𝐦|+t\underset{𝐦,𝐧}{}|𝐦𝐧|,$$
(1)
with nearest neighbor hopping matrix element $`t`$ between s-orbitals $`𝐫|𝐦=\psi (𝐫𝐦)`$ on adjacent atoms located at sites $`𝐦`$ of the lattice. The disorder is simulated by taking random on-site potential such that $`\epsilon _𝐦`$ is uniformly distributed in the interval \[-W/2,W/2\]. This is commonly called the “Anderson model”. There are many numerical studies of the LD transition, which occurs in three-dimensions (3D) for a half-filled band at the critical disorder strength $`W_c16.5t`$. Experiments on real metals with strong scattering or strong correlations often yield resistivities which are hard to analyze. Theory gives guidance in two extreme regimes: (a) the semiclassical case where quasiparticles with definite $`𝐤`$ vector justify a Boltzmann approach and “weak localization” (WL) correction, and (b) a scaling regime near the LD transition to “strong localization”. Lacking a complete theory it is often assumed that the two limits join smoothly with nothing between. Experiments, however, are very often in neither extreme limit. The middle is wide and needs more attention.
Here we give a 3D numerical analysis focused not on the transition itself but instead on the resistivity for $`1<W/t<W_c/t`$; specifically we ask how rapidly does the resistivity $`\rho (W)`$ deviate from the values predicted by the usual Boltzmann theory valid when $`Wt`$. It has long been assumed that “Ioffe-Regel condition” $`\mathrm{}1/k_Fa`$ ($`\mathrm{}`$ being the mean free path, and $`a`$ being the lattice constant) gives the criterion for sufficient disorder to drive the metal into an Anderson insulator. Figure 1 shows that this is wrong. By $`W/t4`$, where $`\mathrm{}`$ is close to $`a`$, there is little sign of a divergence away from the semiclassical extrapolation, and the LD transition is postponed to much larger values of $`W/t`$.
A cleaner discussion is possible using Kubo theory, which does not define $`\mathrm{}`$, but allows a definition of the diffusivity $`D_i`$ of an eigenstate $`|i`$, as shown below in Eq. (3). In the semiclassical regime, $`D_iD_k=v_k\mathrm{}_k/3`$. The diffusivity $`D_k`$ diminishes as $`(W/t)^2`$ in Boltzmann theory. As $`\mathrm{}/a`$ approaches a minimum value ($`1`$), $`D_i`$ decreases toward $`D_{\text{min}}=ta^2/\mathrm{}`$, which can be regarded as a minimum metallic diffusivity below which localization sets in. But there is a wide range of $`W/t`$ over which $`D_iD_{\text{min}}`$ and yet the Boltzmann scaling $`D(t/W)^2`$ is approximately right. In this regime single particle eigenstates $`|i`$ are neither ballistically propagating nor are they localized. There is a third category: intrinsically diffusive. A wave packet built from such states has zero range of ballistic motion but an infinite range of diffusive propagation. Such states are found not only in a narrow crossover regime but over a wide range of parameters physically accessible in real materials and mathematically accessible in models like the Anderson model. In this regime, there is not a simple scaling parameter nor a universal behavior. But the behavior is quite insensitive to a changes in Fermi energy $`E_F`$ or $`k_BT`$, and scales smoothly with $`W/t`$.
The traditional tool for computation of $`\rho `$ has been the Kubo formula, originally derived for a system in the thermodynamic limit. In a basis of exact single particle electron state $`|i`$ of energy $`ϵ_i`$, this can be written as
$$\sigma =\frac{1}{\rho }=\frac{e^2}{\mathrm{\Omega }}\underset{i}{}\left(\frac{f}{ϵ_i}\right)D_i=e^2N(E_F)\overline{D},$$
(2)
where $`\mathrm{\Omega }`$ is the sample volume, $`f`$ is the equilibrium Fermi-Dirac distribution, $`N(E_F)`$ the density of states at $`E_F`$, $`\overline{D}`$ the mean diffusivity, and state diffusivity is given by
$$D_i=\pi \mathrm{}\underset{j}{}|i|\widehat{v}_x|j|^2\delta (ϵ_iϵ_j)$$
(3)
where $`\widehat{𝐯}`$ is the velocity operator. These formulas, while correct, are hard to use numerically. Thanks to the recent advances in mesoscopic physics, it is now apparent that the Landauer scattering approach (or equivalent “mesoscopic” Kubo reformulation for the finite-size systems ) provides superior numerical efficiency when computing the transport properties of finite disordered conductors. It relates the conductance of a sample to its quantum-mechanical transmission properties. This formalism emphasizes the importance of taking into account the interfaces between the sample and the rest of the circuit. Transport in the sample is phase-coherent (i.e. effectively occurring at zero temperature); the dissipation and thus thermalization of electrons (necessary for the establishment of steady state) takes place in other parts of the circuit.
Our principal result for the (quantum) resistivity of Anderson model, using Landauer-type approach, is shown on Fig. 1 for two different Fermi energies $`E_F=0`$ (half-filled band) and $`E_F=2.4t`$ (approximately 70% filled band but falling somewhat as $`W`$, and thus the band-width, increases). The linearized Boltzmann equation $`e𝐄𝐯_𝐤f/ϵ_𝐤=(dF_𝐤/dt)_{\text{scatt}}`$ serves as a reference theory. Here $`ϵ_𝐤`$ is the energy band for $`W=0`$, namely $`ϵ_𝐤=2t\mathrm{cos}k_\alpha `$, $`\mathrm{}v_{k\alpha }`$ is $`ϵ/k_\alpha `$, and $`F_𝐤`$ is the non-equilibrium distribution. The collision integral is
$$\left(\frac{dF}{dt}\right)_{\text{scatt}}=\frac{2\pi }{\mathrm{}}\underset{𝐤^{}}{}|V_{\mathrm{𝐤𝐤}^{}}|^2(F_𝐤F_𝐤^{})\delta (ϵ_𝐤ϵ_𝐤^{}).$$
(4)
The mean squared matrix element of the random potential $`|V_{\mathrm{𝐤𝐤}^{}}|^2`$, in Born approximation, is $`\overline{\epsilon _𝐦^2}=W^2/12`$, where $`\overline{(\mathrm{})}`$ denotes average over probability distribution $`P(\epsilon _𝐦)=(1/W)\theta (W/2|\epsilon _𝐦|)`$. This equation assumes that quasiparticles propagate with mean free path $`\mathrm{}>a`$ between isolated collision events. The equation is exactly solvable, yielding (for $`k_BTt`$) $`1/\rho _\text{B}=e^2\tau (n/m)_{\text{eff}}`$, with $`(n/m)_{\text{eff}}=v_{kx}^2\delta (ϵ_kE_F)/\mathrm{\Omega }`$, and $`\mathrm{}/\tau =2\pi N(E_F)W^2/12`$. We have evaluated $`(n/m)_{\text{eff}}`$ and $`N(E_F)`$ numerically. To within factors of order one, the Boltzmann-Born answer for the semiclassical resistivity is $`\rho _\text{B}=(\pi \mathrm{}a/e^2)(W/4t)^2`$. When $`W=3t`$ and $`a=3\text{Å}`$, $`\rho _\text{B}`$ is $`125\mu \mathrm{\Omega }`$cm, typical of dirty transition metal alloys, and close to the largest resistivity normally seen in dirty “good” metals. Figure 1 plots $`\rho /\rho _\text{B}`$ versus $`(W/t)^2`$. Even for $`W=10t`$ there is less than a factor of 2 deviation from the (unwarranted) extrapolation of the Boltzmann theory into the regime $`W>t`$. Boltzmann theory can be “improved” by including multiple scattering from single impurities, that is, replacing the impurity potential by the $`𝐓`$-matrix $`𝐓_𝐦(z)=\epsilon _𝐦/(1\epsilon _𝐦g(z))`$ where $`g(z)=(1/N_s)(zϵ_𝐤)^1`$ is the free particle Green function ($`N_s`$ is the number of lattice sites). To next order the mean square $`𝐓`$-matrix is
$$\overline{|𝐓_𝐦(z)|^2}=\frac{W^2}{12}\left(1+\frac{3W^2}{20t^2}(gg^{}+gg+g^{}g^{})+\mathrm{}\right),$$
(5)
where the first term is the Born approximation and the coefficient of the correction ($`𝒪(W^4)`$) changes sign from negative to positive as $`E_F`$ moves from 0 to $`2.4t`$. As shown on Fig. 1, the resistivity does not behave like $`\overline{|𝐓_𝐦(z)|^2}`$; multiple scattering with interference from pairs of impurities is at least equally important, and the “exact” $`\rho (W)`$ is less sensitive to details like $`E_F`$ than is the $`𝐓`$-matrix approximation. The rest of the paper presents the method of calculation and describes a bit of mesoscopic physics of very dirty metals.
The central linear transport quantity in the mesoscopic view, as well as in the scaling theory of localization, is conductance $`G`$ rather than conductivity $`\sigma (L)=L^{2d}G(L)`$ (the bulk conductivity is an intensive material constant defined only in the thermodynamic limit, $`\sigma =lim_L\mathrm{}L^{2d}G(L)`$). We use a Landauer-type formula to get the exact quantum conductance $`G`$ of finite samples with disorder configurations chosen by a random number generator. Finite-size samples permit exact solutions for any strength of disorder. Similar to other recent works, the bulk resistivity is extracted from the disorder-averaged resistance $`R`$ by finding the linear (Ohmic) scaling of $`R`$ versus the length of the sample $`L`$ at fixed cross section $`A`$ (Fig. 2). Two kinds of errors may arise: (a) The transition from the Ohmic regime to the localized regime occurs for length of the sample $`L\xi `$ which happens when $`G=𝒪(2e^2/h)`$. If $`L`$ is made large enough, $`G`$ will always diminish to this magnitude. Therefore, we avoid using the sample sizes with too small $`G`$. (b) Finite-size boundary conditions and non-specular reflection cause density of states and scattering properties of the sample to be slightly altered as compared to the true bulk. We expect these effects to be small for our samples where $`\mathrm{}`$ is smaller than the transverse size $`\sqrt{A}`$.
A two probe measuring configuration is used for computation. The sample is placed between two disorder-free ($`\epsilon _𝐦=0`$) semi-infinite leads connected to macroscopic reservoirs which inject thermalized electrons at electrochemical potential $`\mu _L`$ (from the left) or $`\mu _R`$ (from the right) into the system. The electrochemical potential difference $`eV=\mu _L\mu _R`$ is measured between the reservoirs. The leads have the same cross section as the sample. The hopping parameter in the lead and the one which couples the lead to the sample are equal to the hopping parameter in the sample. Thus, extra scattering (and resistance) at the sample-lead interface is avoided but transport at Fermi energies $`|E_F|`$ greater than the clean-metal band edge $`|E_b|=6t`$ cannot be studied. Hard wall boundary conditions are used in the $`\widehat{y}`$ and $`\widehat{z}`$ directions. The sample is modeled on a cubic lattice with $`N\times N_y\times N_z`$ sites, where $`N_y=N_z=15`$ and lengths $`L=Na`$ are taken from the set $`N\{5,10,15,20\}`$.
The linear conductance is calculated using an expression obtained from the Keldysh technique
$$G=\frac{4e^2}{\pi \mathrm{}}\text{Tr}\left(\text{Im}\widehat{\mathrm{\Sigma }}_L\widehat{G}_{1N}^r\text{Im}\widehat{\mathrm{\Sigma }}_R\widehat{G}_{N1}^a\right).$$
(6)
Here $`\text{Im}\widehat{\mathrm{\Sigma }}_{L,R}=(\widehat{\mathrm{\Sigma }}_{L,R}^r\widehat{\mathrm{\Sigma }}_{L,R}^a)/2i`$ are self-energy matrices ($`r`$-retarded, $`a`$-advanced) which describe the coupling of the sample to the leads, and $`\widehat{G}_{1N}^r`$, $`\widehat{G}_{N1}^a`$ are Green function matrices connecting the layer $`1`$ and $`N`$ of the sample: $`\widehat{G}^{r,a}=(E\widehat{H}\widehat{\mathrm{\Sigma }}^{r,a})^1`$ ($`\widehat{G}^a=[\widehat{G}^r]^{}`$), with $`\widehat{\mathrm{\Sigma }}^r=\widehat{\mathrm{\Sigma }}_L^r+\widehat{\mathrm{\Sigma }}_R^r`$ ($`\widehat{\mathrm{\Sigma }}^a=[\widehat{\mathrm{\Sigma }}^r]^{})`$. The self-energy matrices introduced by the leads are non-zero only on the end layers of the sample adjacent to the leads. They are given by $`\widehat{\mathrm{\Sigma }}_{L,R}^r(𝐧,𝐦)=t^2\widehat{g}_{L,R}^r(𝐧_S,𝐦_S)`$ with $`\widehat{g}_{L,R}^r(𝐧_S,𝐦_S)`$ being the surface Green function of the bare semi-infinite lead between the sites $`𝐧_S`$ and $`𝐦_S`$ in the end atomic layer of the lead (adjacent to the corresponding sites $`𝐧`$ and $`𝐦`$ inside the conductor). Positive definiteness of the operators $`2\text{Im}\widehat{\mathrm{\Sigma }}_{L,R}`$ makes it possible to find their square root and recast the expression under the trace of Eq. (6) as a Hermitian operator. The expression (6) then looks like the Landauer formula involving the transmission matrix $`𝐭`$
$`G`$ $`=`$ $`{\displaystyle \frac{e^2}{\pi \mathrm{}}}\text{Tr}(\mathrm{𝐭𝐭}^{})={\displaystyle \frac{e^2}{\pi \mathrm{}}}{\displaystyle \underset{n=1}{\overset{N_yN_z}{}}}T_n,`$ (7)
$`𝐭`$ $`=`$ $`2\sqrt{\text{Im}\widehat{\mathrm{\Sigma }}_L}\widehat{G}_{1N}^r\sqrt{\text{Im}\widehat{\mathrm{\Sigma }}_R},`$ (8)
or transmission eigenvalues $`T_n`$ when the trace is evaluated in a basis which diagonalizes $`\mathrm{𝐭𝐭}^{}`$.
For the case of two probe geometry the average transmission in the semiclassical transport regime ($`a<\mathrm{}L\xi `$) is given by $`T=\mathrm{}_0/(\mathrm{}_0+L)`$, with $`\mathrm{}_0`$ being of the order of $`\mathrm{}`$.
Thus, the semiclassical limit of the Landauer formula for conductance $`G=(e^2/\pi \mathrm{})MT`$ (measured between points deep inside the reservoirs) in the case of not too strong scattering should have the form
$$G^1=R_C+\rho \frac{L}{A}.$$
(9)
It describes the (classical) series addition of two resistors. The “contact” resistance $`R_C=\pi \mathrm{}/e^2M`$ is non-zero, even in the case of ballistic transport when the second term containing the resistivity $`\rho =(\pi \mathrm{}/e^2)A/\mathrm{}_0M`$ vanishes. Here $`Mk_F^2A`$ is the number of propagating transverse modes at $`E_F`$, also referred to as “channels”. A ballistic conductor with a finite cross section can carry only finite currents (the voltage drop occurs at the lead-reservoir interface). Using this simple analysis for guidance, we plot average resistances (taken over $`N_{\text{conf}}=200`$ realization of disorder) versus $`L`$ in Fig. 2, and fit with the linear function
$$R=C_1+C_2L,C_2=\rho /A.$$
(10)
The resistivity $`\rho `$ on Fig. 1 is obtained from the fitted value of $`C_2`$. For very small values of $`W`$ the constant $`C_1`$ is approximately equal to $`R_C=\pi \mathrm{}/e^2M`$ (where $`M=147`$ is the number of open channels in the band center). To our surprise, $`C_1`$ diminishes steadily with increasing $`W`$, and even turns negative around $`W7t`$.
The quantum conductance $`G`$ fluctuates from sample to sample exhibiting universal conductance fluctuations (UCF) $`\mathrm{\Delta }G=\sqrt{\text{Var}G}e^2/\pi \mathrm{}`$ in the semiclassical transport regime $`Ge^2/\pi \mathrm{}`$. The inset on Fig. 2 shows the distribution of resistance $`P_L(R)`$ for our numerically generated impurity ensemble. The error bars, used as weights in the fit (10), are computed as $`\delta R=\sqrt{\text{Var}R/N_{\text{conf}}}`$. We find that $`\mathrm{\Delta }G`$ is indeed independent of the size $`L`$ (of cubic samples), but decreases systematically by a factor $`3`$ as $`W`$ increases to the critical value $`W_c`$ (Fig. 3). On the other hand, $`\mathrm{\Delta }R`$, being similar to $`\mathrm{\Delta }G/G^2`$, depends on sample size. As $`W`$ approaches $`W_c`$, $`G`$ gets smaller until (for our finite samples) $`\mathrm{\Delta }G/G1`$. At this point the distribution of resistances $`R=1/G`$ becomes very broad and $`R`$ begins to rise above $`1/G`$. For $`L=15`$ this happens when $`W12t`$. At large $`W`$ the conductance of long samples ($`N=20`$) becomes close to $`e^2/\pi \mathrm{}`$ and deviations from Ohmic scaling are expected. Therefore, we do not use these points in the fitting procedure when $`W10t`$ (keeping the conductance of the fitted samples $`G>2e^2/\pi \mathrm{}`$).
We do not have a complete explanation for the deviation of $`C_1`$ (10) from the quantum point contact resistance $`R_C`$. In the semiclassical regime $`Ge^2/\pi \mathrm{}`$ there are corrections to the Ohmic scaling $`GL^{d2}`$. The Diffuson-Cooperon diagrammatic perturbation theory gives a (negative) WL correction
$$\sigma (L)=\sigma +\frac{e^2}{\pi ^2\mathrm{}2\sqrt{2}}\frac{1}{L}\frac{e^2}{\pi ^3\mathrm{}}\frac{1}{\mathrm{}_0^{}},$$
(11)
where $`\mathrm{}_0^{}`$ is a length of order $`\mathrm{}`$ (its precise value does not lead to observable consequences in the experiments studying WL, as long as it is unaffected by the temperature and the magnetic field). The positive $`1/L`$ term in Eq. 11 provides a possible picture for our finding that $`C_1`$in Eq. (10) goes negative as $`W`$ increases. However, this picture is an extrapolation from the semiclassical into the “middle” regime of intrinsically diffusive states, and therefore should be given little weight. The negative values of $`C_1`$ is better regarded as a new numerical result from the mesoscopic dirty metal theory.
It is interesting to note that in many $`d`$-band intermetallic compounds, $`\rho `$ “saturates” at a constant value rather than following the semiclassical extrapolation, that is, increasing linearly with $`T`$ at high $`T`$. High $`T_c`$ materials and doped C<sub>60</sub> metals, on the other hand, do not saturate. Within Boltzmann theory, the static disorder measured by $`(W/t)^2`$ plays the same role as thermal disorder or squared lattice displacement $`k_BT`$. Our numerical results thus can be described as “failing to saturate.” Similar failure was seen in high $`T`$ Monte Carlo studies by Gunnarsson and Han.
We thank I. L. Aleiner for interesting discussions and challenging questions. Suggestions provided by J. A. Vergés are acknowledged. This work was supported in part by NSF grant no. DMR 9725037.
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# REFERENCES
MATROID THEORY AND CHERN-SIMONS
J. A. Nieto<sup>1</sup><sup>1</sup>1nieto@uas.uasnet.mxand M. C. Marín
Facultad de Ciencias Físico-Matemáticas de la Universidad Autónoma
de Sinaloa, 80010 Culiacán Sinaloa, México
Abstract
It is shown that matroid theory may provide a natural mathematical framework for a duality symmetries not only for quantum Yang-Mills physics, but also for M-theory. Our discussion is focused in an action consisting purely of the Chern-Simons term, but in principle the main ideas can be applied beyond such an action. In our treatment the theorem due to Thistlethwaite, which gives a relationship between the Tutte polynomial for graphs and Jones polynomial for alternating knots and links, plays a central role. Before addressing this question we briefly mention some important aspects of matroid theory and we point out a connection between the Fano matroid and D=11 supergravity. Our approach also seems to be related to loop solutions of quantum gravity based in Ashtekar formalism.
Pacs numbers: 04.60.-m, 04.65.+e, 11.15.-q, 11.30.Ly
May, 2000
1.- INTRODUCTION
In the last few years, duality has been a source of great interest to study nonperturbative, as well as perturbative, dynamics of superstrings and supersymmetric Yang-Mills . In fact, duality is the key physical concept that relates the five known superstring theories in 9+1 dimensions (i.e. nine space and one time), Type I, Type IIA, Type IIB, Heterotic SO(32) and Heterotic E$`{}_{8}{}^{}\times `$ E<sub>8</sub>, which may now be understood as different manifestations of one underlying unique theory called M-theory -. However, dualities are still a mystery and up to now a general understanding how these dualities arises is missing. Nevertheless, just as the equivalence principle is a basic principle in general relativity, the recent importance of dualities in gauge field theories and string theories strongly suggest a duality principle as a basic principle in M-theory. In this sense, what it is needed is a mathematical framework to support such a duality principle.
M-theory is defined as a 10+1 dimensional theory arising as the strong-coupling limit of type IIA string theory. Essentially, M-theory is a non-pertubative theory and in addition to the five superstring theories it describes supermembranes , 5-branes and D=11 supergravity . Although the complete M-theory is unknown there are two main proposed routes to construct it. One is the N=(2,1) superstring theory and the other M(atrix)-theory . Martinec has suggested that these two scenarios may, in fact, be closely related. This scenario has been extended to include dualities involving compactifications on timelike circles as well as spacelike circles ones. In particular, it has been shown that T-duality on a timelike circle takes type IIA theory into a type IIB theory and type IIB theory into a type IIA theory and that the strong-coupling limit of type IIA$`^\text{ }`$is a theory in 9+2 dimensional theory, denoted by M.
More recently, Khoury and Verlinde have shed some new light on the old idea of open/closed string duality . This duality is of special interest because emphasizes the idea that closed string dynamics (gravity) is dual to open string dynamics (gauge theory). Two previous examples on this direction are matrix theory , where gravity arises as an effect of open string quantum fluctuations and Maldacena’s conjecture that anti-deSitter supergravity is in some sense dual to supersymmetric gauge theory.
Thus, just as the tensor theory makes mathematical sense of the postulate of relativity “the laws of physics are the same for every observer”, we are pursuing the possibility that the mathematical formalism necessary to make sense of a duality principle in M-theory is matroid theory . This theory is a generalization of matrices and graphs and , in contrast to graphs in which duality can be defined only for planar graphs, it has the remarkable property that duality can be defined for every matroid. Since M(atrix)-theory and N=(2,1) superstrings have had an important success on describing some essential features of M-theory a natural question is to see whether matroid theory is related to these two approaches. As a first step in this direction we may attempt to see if matroid theory is linked somehow to D=11 supergravity which is a common feature of both formalism. In fact, it has been shown that the Fano matroid and its dual are closely related to Englert’s compactification of D=11 supergravity. This result is physically interesting because such a relation allows the connection between the fundamental Fano matroid or its dual and octonions which, at the same time, are one of the alternative division algebras . It is worth mentioning that some time ago Blencowe and Duff raised the question whether the four forces of nature correspond to the four divisions algebras.
In this work, we make further progress on this program. Specifically, we find a route to incorporate matroid theory in quantum Yang-Mills in the context of Chern-Simons action. Our mechanism is based on a theorem due to Thistlethwaite which connect the Jones polynomial for alternating knots with the Tutte polynomial for graphs. Since Witten showed that Jones polynomial can be understood in three dimensional terms through a Chern-Simons formalism it became evident that we have a bridge between graphs and Chern-Simons. In this context duality, which is the main subject in graphic matroids, can be associated to Chern-Simons in a mathematical natural way. This connection may transfer important theorems from matroid theory to fundamental physics. For instance, the theorem due to Whitney that if $`M_1,..,M_p`$ and $`M_1^{^{}},..,M_p^{^{}}`$ are the components of the matroids $`M`$ and $`M^{}`$ respectively, and if $`M_i^{^{}}`$ is the dual of $`M_i`$ ($`i=1,\mathrm{},p`$) then $`M^{}`$ is dual of $`M`$ and conversely, if $`M`$ and $`M^{}`$ are dual matroids then $`M_i^{^{}}`$ is dual of $`M_i`$ may be applied to dual Chern-Simons partition functions. One of the aims of this work is to explain how this can be done.
The plan of this work is as follows. In section 2, we briefly review matroid theory and in section 3 we closely follow the reference to discuss a connection between matroid theory and D=11 supergravity. In section 4, we study the relation between matroid theory and Witten’s partition function for knots. Finally, in section 5, we make some final comments.
2.- A BRIEF REVIEW OF MATROID THEORY
In 1935, while working on abstract properties of linear dependence, Whitney introduced the concept of matroid. In the same year, Birkhoff established the connection between simple matroids (also known as combinatorial geometries ) and geometric lattices. In 1936, Mac Lane gave an interpretation of matroids in terms of proyective geometry. And an important progress to the subject was given in 1958 by Tutte who introduced the concept of homotopy for matroids. At present, there is a large body of information about matroid theory. The reader interested in the subject may consult the excellent books on matroid theory by Welsh , Lawler and Tutte . We also recommend the books of Wilson , Kung and Ribnikov .
An interesting feature of matroid theory is that there are many different but equivalent ways of defining a matroid. In this respect, it seems appropriate to briefly review the Whitney’s discovery of the matroid concept. While working with linear graphs Whitney noticed that for certain matrices duality had a simple geometrical interpretation quite different than in the case of graphs. Further, he observed that any subset of columns of a matrix is either linearly independent or linearly dependent and that the following two theorems must hold:
(a) any subset of an independent set is independent.
(b) if N<sub>p</sub> and N<sub>p+1</sub> are independent sets of p and p+1 columns respectively, then N<sub>p</sub> together with some column of N<sub>p+1</sub> forms an independent set of p+1 columns.
Moreover, he discovered that if these two statements are taking as axioms then there are examples that do not represent any matrix and graph. Thus, he concluded that a system satisfying (a) and (b) should be a new one and therefore deserved a new name: He called to this kind of system a matroid.
The definition of a matroid in terms of independent sets has been refined and is now expressed as follows: A matroid $`M`$ is a pair (E,$`)`$, where E is a non-empty finite set, and $``$ is a non-empty collection of subsets of E (called independent sets) satisfying the following properties:
($``$ i) any subset of an independent set is independent;
($``$ ii) if I and J are independent sets with I$``$ J, then there is an element $`e`$ contained in J but not in I, such that I$`\{e\}`$ is independent.
A base is defined to be any maximal independent set. By repeatedly using the property ($``$ ii) it is straightforward to show that any two bases have the same number of elements. A subset of E is said to be dependent if it is not independent. A minimal dependent set is called a circuit. Contrary to the bases not all circuits of a matroid have the same number of elements.
An alternative definition of a matroid in terms of bases is as follows:
A matroid $`M`$ is a pair (E, $``$), where E is a non-empty finite set and $``$ is a non-empty collection of subsets of E (called bases) satisfying the following properties:
($``$ i) no base properly contains another base;
($``$ ii) if B<sub>1</sub> and B<sub>2</sub> are bases and if $`b`$ is any element of B$`_1,`$ then there is and element $`g`$ of B<sub>2</sub> with the property that (B<sub>1</sub>-{$`b`$})$`\{g\}`$ is also a base.
A matroid can also be defined in terms of circuits:
A matroid $`M`$ is a pair (E, $`𝒞)`$, where E is a non-empty finite set, and $`𝒞`$ is a collection of a non-empty subsets of E (called circuits) satisfying the following properties.
($`𝒞`$ i) no circuit properly contains another circuit;
($`𝒞`$ ii) if $`𝒞_1`$ and $`𝒞_2`$ are two distinct circuits each containing and element $`c`$, then there exists a circuit in $`𝒞_1`$ $``$ $`𝒞_2`$ which does not contain $`c`$.
If we start with any of the three definitions the other two follows as a theorems. For example, it is possible to prove that ($``$ ) implies ($``$) and ($`𝒞`$). In other words, these three definitions are equivalent. There are other definitions also equivalent to these three, but for the purpose of this work it is not necessary to consider them.
Notice that even from the initial structure of a matroid theory we find relations such as independent-dependent and base-circuit which suggests duality. The dual of $`M`$, denoted by $`M^{},`$ is defined as a pair (E, $`^{}`$), where $`^{}`$ is a non-empty collection of subsets of E formed with the complements of the bases of M. An immediate consequence of this definition is that every matroid has a dual and this dual is unique. It also follows that the double-dual $`M^{}`$ is equal to $`M`$. Moreover, if A is a subset of E, then the size of the largest independent set contained in A is called the rank of A and is denoted by $`\rho `$(A). If $`M=M_1+M_2`$ and $`\rho `$($`M`$) = $`\rho `$($`M_1`$) +$`\rho `$($`M_2`$) we shall say that $`M`$ is separable. Any maximal non-separable part of $`M`$ is a component of $`M`$. An important theorem due to Whitney is that if $`M_1,..,M_p`$ and $`M_1^{^{}},..,M_p^{^{}}`$ are the components of the matroids $`M`$ and $`M^{}`$ respectively, and if $`M_i^{^{}}`$ is the dual of $`M_i`$ (i = 1,…,p). Then $`M^{}`$ is dual of $`M`$. Conversely, let $`M`$ and $`M^{}`$ be dual matroids, and let $`M_1,..,M_p`$ be components of $`M`$. Let $`M_1^{^{}},..,M_p^{^{}}`$ be the corresponding submatroids of $`M^{}`$. Then $`M_1^{^{}},..,M_p^{^{}}`$ are the components of $`M^{}`$, and $`M_i^{^{}}`$ is dual of $`M_i.`$
3.- MATROID THEORY AND SUPERGRAVITY
Among the most important matroids we find the binary and regular matroids. A matroid is binary if it is representable over the integers modulo two. Let us clarify this definition. An important problem in matroid theory is to see which matroids can be mapped in some set of vectors in a vector space over a given field. When such a map exists we are speaking of a coordinatization (or representation) of the matroid over the field. Let GF(q) denote a finite field of order q. Thus, we can express the definition of a binary matroid as follows: A matroid which has a coordinatization over GF(2) is called binary. Furthermore, a matroid which has a coordinatization over every field is called regular. It turns out that regular matroids play a fundamental role in matroid theory, among other things, because they play a similar role that planar graphs in graph theory . It is known that a graph is planar if and only if it contains no subgraph homeomorphic to K<sub>5</sub> or K<sub>3,3</sub>. The analogue of this theorem for matroids was provided by Tutte . In fact, Tutte showed that a matroid is regular if and only if is binary and includes no Fano matroid or the dual of this. In order to understand this theorem it is necessary to define the Fano matroid. We shall show that the Fano matroid may be connected with octonions which, in turn, are related to the Englert’s compactification of D=11 supergravity.
A Fano matroid F is the matroid defined on the set E={1,2,3,4,5.6.7} whose bases are all those subsets of E with three elements except f$`{}_{1}{}^{}=`${1,2,4}, f$`{}_{2}{}^{}=`${2,3,5}, f$`{}_{3}{}^{}=`$ {3,4,6}, f$`{}_{4}{}^{}=`${4,5,7}, f$`{}_{5}{}^{}=`${5,6,1}, f$`{}_{6}{}^{}=`${6,7,2} and f$`{}_{7}{}^{}=`${7,1,3}. The circuits of the Fano matroid are precisely these subsets and its complements. It follows that these circuits define the dual F of the Fano matroid.
Let us write the set E in the form $``$=$`\{e_1,e_2,e_{3,}e_{4,}e_5,e_6,e_7\}`$. Thus, the subsets used to define the Fano matroid now become $`f_1=\{e_1,e_2,e_4\}`$, $`f_2\{e_2,e_3,e_5\}`$, $`f_3\{e_3,e_4,e_6\}`$, $`f_4\{e_4,e_5,e_7\}`$, $`f_5\{e_5,e_6,e_1\}`$, $`f_6\{e_6,e_7,e_2\}`$ and $`f_7`$ $`\{e_7,e_1,e_3\}`$. The central idea is to identify the quantities $`e_i,`$ where $`i=1,\mathrm{}`$,$`7,`$ with the octonionic imaginary units. Specifically, we write an octonion $`q`$ in the form $`q=q_0e_0+q_1e_1+q_2e_2+q_3e_3+q_4e_4+q_5e_5+q_6e_6+q_7e_7,`$ where $`q_0`$ and $`q_i`$ are real numbers. Here, $`e_0`$ denotes the identity. The product of two octonions can be obtained with the rule:
$$e_ie_j=\delta _{ij}+\psi _{ij}^ke_k,$$
(1)
where $`\delta _{ij}`$ is the Kronecker delta and $`\psi _{ijk}=`$ $`\psi _{ij}^l\delta _{lk}`$ is the fully antisymmetric structure constants, with $`i,j,k=1,\mathrm{},7`$. By taking the $`\psi _{ijk}`$ equals 1 for one of the seven combinations $`f_i`$ we may derive all the values of $`\psi _{ijk}`$.
The octonion (Cayley) algebra is not associative, but alternative. This means that the basic associator of any three imaginary units is
$$(e_i,e_j,e_k)=(e_ie_j)e_ke_i(e_je_k)=\phi _{ijkm}e_m,$$
(2)
where $`\phi _{ijkl}`$ is a fully antisymmetric four index tensor. It turns out that $`\phi _{ijkl}`$ and $`\psi _{ijk}`$ are related by the expression
$$\phi _{ijkl}=(1/3!)ϵ_{ijklmnr}\psi _{mnr},$$
(3)
where $`ϵ_{ijklmnr}`$ is the completely antisymmetric Levi-Civita tensor, with $`ϵ_{12\mathrm{}7}=1`$. It is interesting to note that given the numerical values f<sub>i</sub> for the indices of $`\psi _{mnr}`$ and using (3) we get the other seven subsets of E with four elements of the dual Fano matroid F$`^{}.`$ For instance, if we take f<sub>1</sub> then we have $`\psi _{124}`$ and (3) gives $`\phi _{3567}`$ which leads to the circuit subset $`\{3,5,6,7\}`$.
We would like now to relate the above structure to the Englert’s octonionic solution of eleven dimensional supergravity. First, let us introduce the metric
$$g_{ab}=\delta _{ij}h_a^ih_b^j,$$
(4)
where $`h_a^i`$ = $`h_a^i(x^c)`$ is a sieben-bein, with $`a,b,c=1,\mathrm{},7`$. Here, $`x^c`$ are a coordinates patch of the geometrical seven sphere S<sup>7</sup>. The quantities $`\psi _{ijk}`$ can now be related to the S<sup>7</sup> torsion in the form
$$T_{abc}=R_0^1\psi _{ijk}h_a^ih_b^jh_c^k,$$
(5)
where $`R_0`$ is the S<sup>7</sup> radius. While the quantities $`\phi _{ijkl}`$ can be identified with the four index gauge field $`F_{abcd}`$ through the formula
$$F_{abcd}=R_0^1\phi _{ijkl}h_a^ih_b^jh_c^kh_d^l.$$
(6)
Furthermore, it is possible to prove that the Englert’s 7-dimensional covariant equations are solve with the identification $`F_{abcd}=\lambda T_{[abc},_{d]},`$ where $`\lambda `$ is a constant. Therefore, $`\lambda T_{abc}=A_{abc}`$ is the fully antisymmetric gauge field which is a fundamental object in 2-brane theory .
It is important to mention that in the Englert’s solution of D= 11 supergravity the torsion satisfies the Cartan-Schouten equations
$$T_{acd}T_{bcd}=6R_0^2g_{ab},$$
(7)
$$T_{ead}T_{dbf}T_{fce}=3R_0^2T_{abc.}$$
(8)
But as Gursey and Tze noted, these equations are mere septad-dressed, i.e. covariant forms of the algebraic identities
$$\psi _{ikl}\psi _{jkl}=6\delta _{ij},$$
(9)
$$\psi _{lim}\psi _{mjn}\psi _{nkl}=3\psi _{ijk},$$
(10)
respectively. It is worth mentioning that Englert solution realizes the riemannian curvature-less but torsion-full Cartan-geometries of absolute parallelism on S<sup>7</sup>.
So, we have shown that the Fano matroid is closely related to octonions which at the same time are an essential part of the Englert’s solution of absolute parallelism on S<sup>7</sup> of D=11 supergravity. The Fano matroid and its dual are the only minimal binary irregular matroids. We know from Hurwitz theorem (see reference ) that octonions are one of the alternative division algebras (the others are the real numbers, the complex numbers and the quaternions). While among the only parallelizable spheres we find S<sup>7</sup> (the other are the spheres S<sup>1</sup> and S<sup>3</sup> ). This distinctive and fundamental role played by the Fano matroid, octonions and S<sup>7</sup> in such different areas in mathematics as combinatorial geometry, algebra and topology respectively lead us to believe that the relationship between these three concepts must have a deep significance not only in mathematics, but also in physics. Of course, it is known that the parallelizability of S<sup>1</sup>, S<sup>3</sup> and S<sup>7</sup> has to do with the existence of the complex numbers, the quaternions and the octonions respectively (see reference ). It is also known that using an algebraic topology called K-theory we find that the only dimensions for division algebras structures on Euclidean spaces are 1, 2, 4, and 8. We can add to these remarkable results another fundamental concept in matroid theory; the Fano matroid.
4.- MATROID THEORY AND CHERN-SIMONS
Before going into details, it turns out to be convenient to slightly modify the notation of the previous section. In this section, we shall assume that the Greek indices $`\alpha ,\beta ,\mathrm{},etc`$ run from 0 to 3, the indices $`i,j,\mathrm{},etc`$ run from 0 to 2 and finally the indices $`a,b,\mathrm{},etc`$ take values in the rank of a compact Lie Group G. Further, we shall denote a compact oriented four manifold as $`M^4.`$
Consider the second Chern class action
$$S=\frac{k}{16\pi }_{M^4}ϵ^{\mu \nu \alpha \beta }F_{\mu \nu }^aF_{\alpha \beta }^b\text{ }g_{ab},$$
(11)
with the curvature given by
$$F_{\alpha \beta }^a=_\alpha A_\beta ^a_\beta A_\alpha ^a+C_{bc}^aA_\alpha ^bA_\beta ^c.$$
(12)
Here $`g_{ab}`$ is the Killling-Cartan metric and $`C_{bc}^a`$ are the completely antisymmetric structure constants associated to the compact simple Lie group G. The action (11) is a total derivative and leads to the Chern-Simons action
$$S_{CS}=\frac{k}{4\pi }_{M^3}\{ϵ^{ijk}(A_i^a(_jA_k^b_kA_j^b)g_{ab}+\frac{2}{3}C_{abc}A_i^aA_j^bA_k^c)\},$$
(13)
where $`M^3=M^4`$ is a compact oriented three dimensional manifold. In a differential forms notation (13) can be rewritten as follows:
$$S_{CS}=\frac{k}{2\pi }_{M^3}Tr(AdA+\frac{2}{3}AAA),$$
(14)
where $`A=A_i^aT_adx^i`$, with $`T_a`$ the generators of the Lie algebra of G.
Given a link $`L`$ with $`r`$ components and irreducible representation $`\rho _r`$ of G, one for each component of the link, Witten defines the partition function
$$Z(L,k)=D𝒜\text{exp(}iS_{cs})_{r=1}^nW(L_r\text{,}\rho _r),$$
(15)
where $`W(C_i,\rho _i)`$ is the Wilson line
$$W(L_r,\rho _r)=Tr_{\rho _r}P\mathrm{exp}(\underset{L_r}{}A_i^aT_a𝑑x^i).$$
(16)
Here the symbol $`P`$ means the path-ordering along the knots $`L_r.`$
If we choose $`M^3=S^3,`$ $`G=SU(2)`$ and $`\rho _r`$=C<sup>2</sup> for all link components then the Witten’s partition function (15) reproduces the Jones polynomial
$$Z(L,k)=V_L(t),$$
(17)
where
$$t=e^{\frac{2\pi i}{k}}.$$
(18)
Here $`V_L(t)`$ denotes the Jones polynomial satisfying the skein relation:
$$t^1V_{L_+}tV_L_{}=(\sqrt{t}\frac{1}{\sqrt{t}})V_{L_0},$$
(19)
where $`L_+,L_{}`$ and $`L_0`$ are the standard notation for overcrossing, undercrossing and zero crossing.
Now, lets pause about the relation between the knots and Chern-Simons term and let us discuss the Tutte polynomial. To each graph $`𝒢`$, we associate a polynomial $`T_𝒢(x,x^1)`$ with the property that if $`𝒢`$ is composed solely of isthmus and loops then $`T_𝒢(x,x^1)=x^Ix^l,`$ where $`I`$ is the number of isthmuses and $`l`$ is the number of loops. The polynomial $`T_𝒢`$ satisfies the skein relation
$$T_𝒢=T_𝒢^{^{}}+T_{𝒢^{^{\prime \prime }}},$$
(20)
where $`𝒢^{}`$ and $`𝒢^{\prime \prime }`$ are obtained by delating and contracting respectively an edge that is neither a loop nor an isthmus of $`𝒢`$.
There is a theorem due to Thistlethwaite which assures that if $`L`$ is an alternating link and $`𝒢(L)`$ the corresponding planar graph then the Jones polynomial $`V_L(t)`$ is equal to the Tutte polynomial $`T_𝒢(t,t^1)`$ up to a sign and factor power of $`t.`$ Specifically, we have
$$V_L(t)=(t^{\frac{3}{4}})^{w(L)}t^{\frac{(rn)}{4}}T_𝒢(t,t^1)$$
(21)
where $`w(L)`$ is the writhe and $`r`$ and $`n`$ are the rank and the nullity of $`𝒢`$ respectively. Here $`V_L(t)`$ is the Jones polynomial of alternating link $`L.`$
On the other hand, a theorem due to Tutte allows to compute $`T_𝒢(t,t^1)`$ from the maximal trees of $`𝒢`$. In fact, Tutte proved that if $``$ denotes the maximal trees in a graph $`𝒢`$, $`i(B)`$ denotes the number of internally active edges in $`𝒢`$ and $`e(B)`$ the number the externally active edges in $`𝒢`$ (with respect to a given maximal tree $`Bϵ`$) then the Tutte polynomial is given by the formula
$$T_𝒢(t,t^1)=x^{i(B)}x^{e(B)}$$
(22)
where the sum is over all elements of $``$.
First, note that $``$ is the collection of bases of $`𝒢`$. If we now remember our definition of matroid $`M`$ in terms of bases discussed in section 2 we note the Tutte polynomial $`T_𝒢(t,t^1)`$ computed according to (22) uses the concept of a graphic matroid $`M(𝒢)`$ defined as the pair (E, $``$), where E is the set of edges of $`𝒢`$. In fact, the elements of $``$ satisfy the two properties
($``$ i) no base properly contains another base;
($``$ ii) if $`B_1`$ and $`B_2`$ are bases and if $`b`$ is any element of $`B_1,`$ then there is and element $`g`$ of $`B_2`$ with the property that ($`B_1`$-{$`b`$})$`\{g\}`$ is also a base.
which identifies a $`M(𝒢)`$ as a matroid. With this remarkable connection between the Tutte polynomial and a matroid we have found in fact a connection between the partition function $`Z(L,k)`$ given in (15) and matroid theory. This is because according to (21) the Tutte polynomial $`T_𝒢(t,t^1)`$ are related to the Jones polynomial $`V_L(t)`$ which at the same time according to (17) are related to the partition function $`Z(L,k)`$. Specifically, for $`M^3=S^3,`$ $`G=SU(2)`$, $`\rho _r`$=C<sup>2</sup> for all alternating link components of $`L`$, we have the relation
$$Z(L,k)=V_L(t)=(t^{\frac{3}{4}})^{w(L)}t^{\frac{(rn)}{4}}T_𝒢(t,t^1).$$
(23)
Thus, the matroid (E, $``$) used to compute $`T_𝒢(t,t^1)`$ can be associated not only to $`V_L(t),`$ but also to $`Z(L,k).`$
Now that we have at hand this slightly but important connection between matroid theory and Chern-Simons theory we arre able to transfer information from matroid theory to Chern-Simons and conversely from Chern-Simons to matroid theory. Let us discuss two examples for the former possibility.
First of all, it is known that in matroid theory the concept of duality is of fundamental importance. For example, there is a remarkable theorem that assures that every matroid has a dual. So, the question arises about what are the implications of this theorem in Chern-Simons formalism. In order to address this question let us first make a change of notation $`T_𝒢(t,t^1)T_{M(𝒢)}(t)`$ and $`Z(L,k)Z_{M(𝒢)}(k).`$ The idea of this notation is to emphasize the connection between matroid theory, Tutte polynomial and Chern- Simons partition function. Consider the planar dual graph $`𝒢^{}`$ of $`𝒢`$. In matroid theory we have $`M(𝒢^{})`$ =$`M^{}(𝒢)`$. Therefore, the duality property of the Tutte polynomial
$$T_𝒢(t,t^1)=T_𝒢^{}(t^1,t)$$
(24)
can be expressed as
$$T_{M(𝒢)}(t)=T_{M^{}(𝒢)}(t^1)$$
(25)
and consequently from (23) we discover that for the partition function $`Z_{M(𝒢)}(k)`$ we should have the duality property
$$Z_{M(𝒢)}(k)=Z_{M^{}(𝒢)}(k).$$
(26)
This duality symmetry for the partition function $`Z_{M(𝒢)}(k)`$ is not really new, but is already known in the literature as mirror image symmetry (see, for instance , and references quoted there). However, what seems to be new is the way we had derived it.
As a second example let us first mention another theorem due to Withney : If $`M_1,..,M_p`$ and $`M_1^{^{}},..,M_p^{^{}}`$ are the components of the matroids $`M`$ and $`M^{}`$ respectively, and if $`M_i^{^{}}`$ is the dual of $`M_i`$ $`(i=1,\mathrm{},p)`$. Then $`M^{}`$ is dual of $`M`$. Conversely, let $`M`$ and $`M^{}`$ be dual matroids, and let $`M_1,..,M_p`$ be components of $`M`$. Let $`M_1^{^{}},..,M_p^{^{}}`$ be the corresponding submatroids of $`M^{}`$. Then $`M_1^{^{}},..,M_p^{^{}}`$ are the components of $`M^{}`$, and $`M_i^{^{}}`$ is dual of $`M_i.`$ Thus, according to (26) we find that
$$Z_{M_i(𝒢_i)}(k)=Z_{M_i^{^{}}(𝒢_i)}(k)$$
(27)
if and only if
$$Z_{M(𝒢)}(k)=Z_{M^{}(𝒢)}(k),$$
(28)
where $`𝒢_i`$ are the components of $`𝒢.`$
5.- COMMENTS
Motivated by a possible duality principle in M-theory we have started to bring information from matroid theory to fundamental physics. We now have two good examples which indicate that this task makes sense. In the first example, we have found enough evidence for a connection between the Fano matroid and supergravity in D=11. While in the second example, we have found a relation between the graphic matroid and the Witten’s partition function for Chern-Simons. This relation is of special importance because leads us to a duality symmetry in the partition function $`Z_{M(𝒢)}(k)`$. In fact, if there is a duality principle in M-theory we should expect to a have a duality symmetry in the corresponding partition function associated to M-theory.
In this work, we have concentrated in the original connection between Chern-Simons action and knots theory. But it is well known that Chern-Simons formalism and knots connection has a number of extensions . It will be interesting to study such extensions from the point of view of matroid theory. It is also known Chern-Simons formalism is closely related to conformal field theory and this in turn is closely related to string theory. So, it seems that the present work may eventually leads to a connection between matroid theory and string theory. In order to achieve this goal we need to study the relation between matroids and Chern-Simons using signed graphs . This is because general knots and links (not only alternating) are one to one correspondence with signed planar graphs. This in turn are straightforward connected with Kauffmann polynomials which at the same time leads to the Jones polynomials. But, signed graphs leads to signed matroids. So, one of our future goals will be to find a connection between signed matroids and string theory. Moreover, matrix Chern-Simons theory has a straightforward relation with Matrix-model and non-commutative geometry . So, a natural extension of the present work will be to study the relation between matroid theory and matrix Chern-Simons theory.
An important duality in M-theory is the strong/weak coupling $`S`$-duality which provides with one of the most important techniques to study non-perturbative aspects of gauge field theory and string theory. For further work it may also be important to find the relation between the duality symmetry for $`Z_{M(𝒢)}(k)`$ given in (27) and $`S`$-duality.
Besides of the possible connection between M(atroid)-theory and M-theory there is another interesting physical application of the present work. This has to do with loop solutions of quantum gravity based in Ashtekar formalism. It is known that the Witten’s partition function provides a solution of the Ashtekar constriants . So, the duality symmetries (27) also applies to such solutions. In other words, it seems that we have also found a connection between matroid theory and loop solutions of quantum canonical gravity.
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# 1 Introduction
## 1 Introduction
Coherency and interference are the basic physical phenomena, which lead to new effects in quantum optics. Interference can be destructive or, vice versa, constructive, causing mutual suppression or amplification of simultaneous processes. In quantum optics the interference effects can reduce interaction of radiation with absorbing atoms at lower energy levels without substantial variations in the interaction with emitting atoms at overlying levels. In turn, this leads to a difference in absorption and emission spectra. The intratomic coherence conditions many fundamental effects in high-resolution nonlinear spectroscopy, the light amplification without population inversion and resonant refraction increase in the absence of absorption, coherent population trapping, increase of the resonant nonlinear-optical radiation at a simultaneous primary radiation absorption decrease, and the laser induction of structures of the autoionization type in spectral continua. The resonant nonlinear interference effects, theoretically and experimentally studied since the time of first masers and lasers creation \[1-6\], currently again attracted a great attention \[7-10\]. They are promising for new laser sources in the VUV and X-ray bands, laser accelerators of atomic particles, microscopes with increased resolution, supersensitive magnetometers, etc. (see, e.g. . A lot of scientific meetings are dedicated to this problem (see, e.g. ).
A publication flux are devoted to the related effect of electromagnetically induced transparency as applied to improve characteristics of laser light conversion to short-wave bands (see. e.g. ). (It appears that in many respects this is the phenomenon comprehensively studied in Russia in 1960-70 s \[12-36, 40\]). A quite complete survey of results acquired by western authors can be found in \[7-10\] and references therein. Therefore, below a basic attention would be given to certain less-known results of Russian authors.
A basic contribution into studies of the nonlinear interference processes in absorption (emission) spectra at the interaction of atom-molecular systems with electromagnetic radiation and the effect of dynamic splitting of spectral lines in strong electromagnetic fields was made by the Yerevan, Moscow, Nizhny Novgorod, Novosibirsk, St. Petersburg, and Minsk schools. These results are generalized in monographs \[4, 15-18, 22, 23, 30, 33, 36\]. In the possible inversionless amplification was analyzed for three-level systems at the discrete optical transitions. Corresponding experimental studies were carried out in . (For two-level optical systems this effect was predicted in and experimentally studied in the radio band in ). The effect of self-transparency induced by a strong field at an adjacent transition was theoretically and experimentally studied in detail, e.g. in . The coherent population trapping was first observed in . Later the study in nonlinear interference phenomena at the discrete transitions was extended to the transitions to a continuous spectrum , autoionization-type resonances were predicted in and experimentally revealed at the atomic transitions to continuum. (Later this effect was called as the laser-induced continuum structure (LICS)). They also predicted the possible nonlinear response increase at the laser short-wave generation with a simultaneous decrease in its absorption and improvement in phase matching . In it was shown that, similar to discrete transitions, the inversionless amplification is possible also at the transitions to autoionization and antoionization-type states. The cited papers initiated the coherence effect study, first at the transitions to continuum, and then at the discrete optical ones.
The coherence and interference phenomena are the basis for inversionless amplification, coherent population trapping, and electromagnetically induced transparency both at the discrete transitions and those to continuous spectrum. As it was already indicated, these effects offer unconventional solution for actual problems of quantum electronics. However, the peculiarities of optical transitions and real experimental designs can qualitatively change an expected manifestation of these processes. These problems remain to be a subject of great attention. Hence, it is required to develop theoretical approaches considering the most important accompanying processes and involving numerical analysis, if necessary. Now the least understood phenomena are the effects of nonuniform broadening and level degeneration, relaxation and motion of population on the coherence degradation. It appears that sometimes, vice versa, the relaxation promotes the intratomic coherence. A fairly small number of papers is dedicated to the effect of above processes on the resonant nonlinear-optical frequency mixing.
## 2 Resonant nonlinear-optical interference
### 2.1 Destructive and constructive interference in classical and quantum optics
The interference is one of fundamental physical phenomena. Oscillations of various nature depending on a phase relationship can interfere constructively or destructively. Varying oscillation phases and amplitudes, the resulting process can be amplified or suppressed. The quantum interference can proceed, when there is coherent superposition of real states. Moreover, the degenerate (in frequency) interfering intra-atomic oscillations can be conditioned by different correlating quantum transitions contributing into the same process. These are, e.g., one-and two-photon contributions into the optical process related to emission or absorption at a specified frequency. The process can result from the coherent superposition of a neighboring real energy level and a quasi-level (virtual state) created by a strong auxiliary field . Such a superposition is realized even more simply than in the case of real doublet state.
The interference is more general concept, than notions of one-, two-, and multi-step and multi-photon processes. The latter were introduced and classified by their frequency-correlation properties in the framework of perturbation theory. However, these properties are significantly varied as the field intensity rises, in particular at resonant interaction . As a result, the qualitative effects become possible in nonlinear spectroscopy of the Doppler-broadened transitions, such as the induced compensation of residual Doppler broadening in two-photon absorption or Raman scattering under conditions of a difference in photon frequencies . Even when many elementary processes contribute into the optical one at a given frequency and their classification into stepped and multi-photon processes is difficult, experimental data can be often explained and predicted using the concept on interfering components of nonlinear polarization. Amplitudes and phases of these components are varied by controlling the corresponding field intensities and the detunings from one- and two-photon resonances.
### 2.2 Equation for the density matrix: the effects of intermediate level population and relaxation
In the general case of an open configuration of energy levels, when a lower level is not ground, various relaxation processes’ rates are different, and all the levels can be populated, the density matrix approach is most convenient to analyze resonant nonlinear-optical processes. Simple formulas for spectral properties of responses at a weak probing field frequency in the presence of strong one at an adjacent transition are uniformly deduced for $`V`$, $`\mathrm{\Lambda }`$, and cascade schemes . Let us show this by the example of transition diagram displayed in Fig. 1.
The scheme is assumed open, i.e. level $`l`$ is not ground one. For simplicity, we consider the fields $`E_2`$ and $`E_4`$ at frequencies $`\omega _2`$ and $`\omega _4`$, are probing, i.e., they do not disturb the level population. The fields $`E_1`$ and $`E_3`$ at frequencies $`\omega _1\omega _{gl}`$ and $`\omega _3\omega _{mn}`$, are strong. Below, in Sec. 3, we eliminate this limitation while considering interaction of two strong fields. Let us acquire the conditions of inversionless amplification at the transitions $`gn`$ and $`ml`$ so that to consider both the configurations $`V`$ and $`\mathrm{\Lambda }`$. The probing field frequencies can be either higher or lower than those of strong fields. In the interaction representation the density matrix components and corresponding equations take the form
$$\rho _{lg}=r_1\mathrm{exp}(i\mathrm{\Omega }_1t),\rho _{nm}=r_3\mathrm{exp}(i\mathrm{\Omega }_3t),\rho _{ng}=r_2\mathrm{exp}(i\mathrm{\Omega }_1t)+\stackrel{~}{r}_2\mathrm{exp}[i(\mathrm{\Omega }_1+\mathrm{\Omega }_3\mathrm{\Omega }_4)t],$$
$$\rho _{lm}=r_4\mathrm{exp}(i\mathrm{\Omega }_4t)+\stackrel{~}{r}_4\mathrm{exp}[i(\mathrm{\Omega }_1\mathrm{\Omega }_2+\mathrm{\Omega }_3)t],\rho _{ln}=r_{12}\mathrm{exp}[i(\mathrm{\Omega }_1\mathrm{\Omega }_2)t]+r_{43}\mathrm{exp}[i(\mathrm{\Omega }_4\mathrm{\Omega }_3)t],\rho _{ii}=r_i,$$
$$P_2r_2=iG_2\mathrm{\Delta }r_2iG_3r_{32}^{}+ir_{12}^{}G_1,d_2\stackrel{~}{r}_2=iG_3r_{41}^{}+ir_{43}^{}G_1,P_4r_4=i[G_4\mathrm{\Delta }r_4G_1r_{41}+r_{43}G_3],$$
$$d_4\stackrel{~}{r}_4=iG_1r_{32}+ir_{12}G_3,P_{41}r_{41}=iG_1^{}r_4+ir_1^{}G_4,P_{43}r_{43}=iG_4r_4^{}+ir_4G_3^{},$$
$$P_{32}r_{32}=iG_2^{}r_3+ir_2^{}G_3,P_{12}r_{12}=iG_1r_2^{}+ir_1G_2^{},\mathrm{\Gamma }_mr_m=2\mathrm{R}\mathrm{e}\{iG_3^{}r_3\}+q_m,$$
$$\mathrm{\Gamma }_nr_n=2\mathrm{R}\mathrm{e}\{iG_3^{}r_3\}+\gamma _{gn}r_g+\gamma _{mn}r_m+q_n,\mathrm{\Gamma }_gr_g=2\mathrm{R}\mathrm{e}\{iG_1^{}r_1\}+q_g,$$
$$\mathrm{\Gamma }_lr_l=2\mathrm{R}\mathrm{e}\{iG_1^{}r_1\}+\gamma _{gl}r_l+\gamma _{ml}r_m+q_l,$$
where $`\mathrm{\Delta }r_1=r_lr_g,\mathrm{\Delta }r_2=r_nr_g,\mathrm{\Delta }r_3=r_nr_m,\mathrm{\Delta }r_4=r_lr_m;`$ $`\mathrm{\Omega }_1=\omega _1\omega _{lg},\mathrm{\Omega }_3=\omega _3\omega _{mn},\mathrm{\Omega }_2=\omega _2\omega _{gn},\mathrm{\Omega }_4=\omega _4\omega _{ml};`$ $`G_1=E_1d_{lg}/2\mathrm{},G_2=E_2d_{gn}/2\mathrm{},G_3=E_3d_{nm}/2\mathrm{},G_4=E_4d_{ml}/2\mathrm{};`$ $`P_1=\mathrm{\Gamma }_{lg}+i\mathrm{\Omega }_1,P_2=\mathrm{\Gamma }_{ng}+i\mathrm{\Omega }_2,P_3=\mathrm{\Gamma }_{nm}+i\mathrm{\Omega }_3,P_4=\mathrm{\Gamma }_{lm}+i\mathrm{\Omega }_4,P_{12}=\mathrm{\Gamma }_{ln}+i(\mathrm{\Omega }_1\mathrm{\Omega }_2),P_{43}=\mathrm{\Gamma }_{ln}+i(\mathrm{\Omega }_4\mathrm{\Omega }_3),P_{32}=\mathrm{\Gamma }_{gm}+i(\mathrm{\Omega }_3\mathrm{\Omega }_2),P_{41}=\mathrm{\Gamma }_{gm}+i(\mathrm{\Omega }_4\mathrm{\Omega }_1),d_2=\mathrm{\Gamma }_{ng}+i(\mathrm{\Omega }_1+\mathrm{\Omega }_3\mathrm{\Omega }_4),d_4=\mathrm{\Gamma }_{lm}+i(\mathrm{\Omega }_1\mathrm{\Omega }_2+\mathrm{\Omega }_3).`$ Here $`\mathrm{\Omega }_i`$ are the frequency detunings from resonances, $`G_i`$ are the Rabi frequencies, $`\mathrm{\Delta }r_i`$, are the population differences depending on intensity, $`\mathrm{\Gamma }_{ij}`$ are the uniform halfwidths of transitions, $`\mathrm{\Gamma }_i^1`$ are the lifetimes, $`\gamma _{ij}`$ is the rate of relaxation from the $`i`$th to $`j`$th levels, and $`q_i`$ are the rates of excitation by additional noncoherent pumping. The off-diagonal density matrix amplitudes $`r_i`$, define the coefficients of absorption (amplification) and refraction indices and $`\stackrel{~}{r}_i`$, define the nonlinear polarization of four-wave mixing. For the cascade configurations, the equations and their solutions are deduced by a simple substitution of detuning signs or by a complex conjugation of corresponding co-factors.
### 2.3 Laser-induced intra-atomic coherence and classification of resonant nonlinear-optical effects
Solution to the set of coupled equations for the density matrix components is given by
$$r_{1,3}=iG_{1,3}\mathrm{\Delta }r_1/P_1,r_{2,4}=iG_{2,4}R_{2,4}/P_{2,4},$$
$$R_2=\frac{\mathrm{\Delta }r_2(1+g_7+v_7)v_3(1+v_7g_8)\mathrm{\Delta }r_3g_3(1+g_7v_8)\mathrm{\Delta }r_1}{(1+g_2+v_2)+[g_7+g_2(g_7v_8)+v_7+v_2(v_7g_8)]},$$
$`(1)`$
$$R_4=\frac{\mathrm{\Delta }r_4(1+v_5+g_5)g_1(1+g_5v_6)\mathrm{\Delta }r_1v_1(1+v_5g_6)\mathrm{\Delta }r_3}{(1+g_4+v_4)+[v_5+v_4(v_5g_6)+g_5+g_4(g_5v_6)]},$$
$`(2)`$
$$\mathrm{\Delta }r_1=\frac{(1+\text{æ}_3)\mathrm{\Delta }n_1+b_1\text{æ}_3\mathrm{\Delta }n_3}{(1+\text{æ}_1)(1+\text{æ}_3)a_1\text{æ}_1b_1\text{æ}_3},\mathrm{\Delta }r_3=\frac{(1+\text{æ}_1)\mathrm{\Delta }n_3+a_1\text{æ}_1\mathrm{\Delta }n_1}{(1+\text{æ}_1)(1+\text{æ}_3)a_1\text{æ}_1b_1\text{æ}_3},$$
$$\mathrm{\Delta }r_2=\mathrm{\Delta }n_2b_2\text{æ}_3\mathrm{\Delta }r_3a_2\text{æ}_1\mathrm{\Delta }r_1,\mathrm{\Delta }r_4=\mathrm{\Delta }n_4a_3\text{æ}_1\mathrm{\Delta }r_1b_3\text{æ}_3\mathrm{\Delta }r_3;$$
$$r_m=n_m+(1b_2)\text{æ}_3\mathrm{\Delta }r_3,r_g=n_g+(1a_3)\text{æ}_1\mathrm{\Delta }r_1,r_n=n_nb_2\text{æ}_3\mathrm{\Delta }r_3+a_1\text{æ}_1\mathrm{\Delta }r_1,$$
$`(3)`$
$$r_l=n_lb_1\text{æ}_3\mathrm{\Delta }r_3+a_3\text{æ}_1\mathrm{\Delta }r_1,\mathrm{\Delta }r_i(E_1=0,E_3=0)=\mathrm{\Delta }n_i;$$
$$g_1=\frac{|G_1|^2}{P_{41}P_1^{}},g_2=\frac{|G_1|^2}{P_{12}^{}P_2},g_3=\frac{|G_1|^2}{P_{12}^{}P_1^{}},g_4=\frac{|G_1|^2}{P_{41}P_4},g_5=\frac{|G_1|^2}{P_{43}d_2^{}},g_6=\frac{|G_1|^2}{P_{41}d_2^{}},g_7=\frac{|G_1|^2}{P_{32}^{}d_4^{}},g_8=\frac{|G_1|^2}{P_{12}^{}d_4^{}},$$
$$v_1=\frac{|G_3|^2}{P_{43}P_3^{}},v_2=\frac{|G_3|^2}{P_{32}^{}P_2},v_3=\frac{|G_3|^2}{P_{32}^{}P_3^{}},v_4=\frac{|G_3|^2}{P_{43}P_4},v_5=\frac{|G_3|^2}{P_{41}d_2^{}},v_6=\frac{|G_3|^2}{P_{43}d_2^{}},v_7=\frac{|G_3|^2}{P_{12}^{}d_4^{}},v_8=\frac{|G_3|^2}{P_{32}^{}d_4^{}};$$
$$\text{æ}_1=\text{æ}_1^0\frac{\mathrm{\Gamma }_{lg}^2}{|P_1|^2},\text{æ}_1^0=\frac{2(\mathrm{\Gamma }_l+\mathrm{\Gamma }_g\gamma _{gl})}{\mathrm{\Gamma }_l\mathrm{\Gamma }_g\mathrm{\Gamma }_{lg}}|G_1|^2,\text{æ}_3=\text{æ}_3^0\frac{\mathrm{\Gamma }_{mn}^2}{|P_3|^2},\text{æ}_3^0=\frac{2(\mathrm{\Gamma }_m+\mathrm{\Gamma }_n\gamma _{mn})}{\mathrm{\Gamma }_m\mathrm{\Gamma }_n\mathrm{\Gamma }_{mn}}|G_3|^2;$$
$$a_1=\frac{\gamma _{gn}a_2}{\mathrm{\Gamma }_n\gamma _{gn}}=\frac{\gamma _{gn}\mathrm{\Gamma }_la_3}{\mathrm{\Gamma }_n(\mathrm{\Gamma }_g\gamma _{gl})}=\frac{\gamma _{gn}\mathrm{\Gamma }_l}{\mathrm{\Gamma }_n(\mathrm{\Gamma }_l+\mathrm{\Gamma }_g\gamma _{gl})},$$
$$b_1=\frac{\gamma _{ml}\mathrm{\Gamma }_nb_2}{\mathrm{\Gamma }_l(\mathrm{\Gamma }_m\gamma _{mn})}=\frac{\gamma _{ml}b_3}{\mathrm{\Gamma }_l(\mathrm{\Gamma }_l\gamma _{ml})}=\frac{\gamma _{ml}\mathrm{\Gamma }_n}{\mathrm{\Gamma }_l(\mathrm{\Gamma }_m+\mathrm{\Gamma }_n\gamma _{mn})}.$$
At $`G_3=0`$ equations (1) and (2) convert in solutions for $`\mathrm{\Lambda }`$ and $`V`$ schemes
$$r_2=i\frac{G_2}{P_2}\frac{\mathrm{\Delta }r_2g_3\mathrm{\Delta }r_1}{1+g_2},r_4=i\frac{G_4}{P_4}\frac{\mathrm{\Delta }r_4g_1\mathrm{\Delta }r_1}{1+g_4}.$$
$`(4)`$
According to , it is convenient to classify the effects of the strong radiation resonant to an adjacent transition, as (i) population saturation (formulas (3)), (ii) dynamic splitting of the resonance for a probing field (or splitting of a common level, i.e., the ac Stark effect, denominators in formulas (4), and (iii) nonlinear interference effects (NIEF) (the terms in the numerators of (4)). The two last effects are conditioned by quantum coherence.
## 3 Difference in pure emission and absorption spectra due to <br>the nonlinear interference effects: inversionless amplification, <br>resonantly amplified refraction in the absence of absorption, <br>and laser-induced transparency
A light emitted of absorbed, e.g., at the frequency $`\omega _2`$ whose power is proportional to $`Re(iG_2^{}r_2)`$, can be considered as a difference between pure emission (a term proportional to $`r_g`$) and pure absorption (other terms in formulas (1) and (4)). The two constituents are positive, but differently depend on detuning due to the NIE. Thus, a sign alternation arises in spectral line contour, resulting in the inversionless amplification. This was emphasized in (see also ). Optimum conditions for the inversionless amplification in a uniformly broadened three-level system were analyzed in \[28-30\] in detail. The refraction index at frequency $`\omega _2`$ is defined as $`Im(iG_2^{}r_2)`$ and, generally, the laser-induced minimum (including zero) absorption can coincide with the resonant refraction index maximum . As is emphasized in \[22-30, 36\], the splitting effect and the NIE as a whole, varying the spectral line shape and causing the difference in pure emission (spontaneous or induced) and absorption spectra, does not vary its integral intensity, which is defined only by saturation effects,
$$𝑑\mathrm{\Omega }_2\mathrm{Re}(ir_2/G_2)=\mathrm{\Delta }r_2,𝑑\mathrm{\Omega }_4\mathrm{Re}(ir_4/G_4)=\mathrm{\Delta }r_4.$$
$`(5)`$
Thus, namely NIE lead to the coherent population trapping, electromagnetically induced transparency and the inversionless amplification, e.g., at the transition $`gn`$ (or $`ml`$), when the second terms in nominators of (4) become equal or begin to exceed $`\mathrm{\Delta }r_2`$ (or $`\mathrm{\Delta }r_4`$). It is seen from the density matrix equation that the considered effects are finally defined by the coherence at transitions $`gm`$ and $`ln`$ ($`r_{32}`$ and $`r_{12}`$), induced jointly by probing and strong fields.
### 3.1 Inversionless amplification of the probing wave
Now we enlarge on the problem, what are the elementary processes usually defined by the perturbation theory, which contribute into absorption (amplification) in the analyzed cases. For instance, let us consider the absorption index $`\alpha (\mathrm{\Omega }_4)`$ at frequency $`\omega _4>\omega _1`$ (Fig. 1) and $`E_3=0`$ normalized to its maximum $`\alpha ^0(0)`$, in the absence of all strong fields. From formulas (4) we find
$$\frac{\alpha (\mathrm{\Omega }_4)}{\alpha ^0(0)}=\mathrm{Re}\left\{\frac{\mathrm{\Gamma }_4[\mathrm{\Delta }r_4g_1\mathrm{\Delta }r_1]}{P_4\mathrm{\Delta }n_4(1+g_4)}\right\}.$$
$`(6)`$
Further we consider the two following cases. (i) Great yields from one-photon resonances
$$|\mathrm{\Omega }_1||\mathrm{\Omega }_4|\mathrm{\Gamma }_1,\mathrm{\Gamma }_4,|g_4|1,|g_1|1,P_4i\mathrm{\Omega }_4,P_1i\mathrm{\Omega }_1i\mathrm{\Omega }_4.$$
Formula (6) takes on the form
$$\frac{\alpha (\mathrm{\Omega }_4)}{\alpha ^0(0)}\frac{\mathrm{\Gamma }_4^2\mathrm{\Delta }r_4}{\mathrm{\Omega }_4^2\mathrm{\Delta }n_4}\mathrm{Re}\left\{\frac{\mathrm{\Gamma }_4(\mathrm{\Delta }r_4g_4+\mathrm{\Delta }r_1g_1)}{i\mathrm{\Omega }_4\mathrm{\Delta }n_4}\right\}\frac{\mathrm{\Gamma }_4^2\mathrm{\Delta }r_4}{\mathrm{\Omega }_4^2\mathrm{\Delta }n_4}\frac{\mathrm{\Gamma }_4\mathrm{\Gamma }_{14}}{\mathrm{\Gamma }_{14}^2+(\mathrm{\Omega }_4\mathrm{\Omega }_1)^2}\frac{|G_1|^2(\mathrm{\Delta }r_1\mathrm{\Delta }r_4)}{\mathrm{\Omega }_4^2\mathrm{\Delta }n_4}$$
$$=\frac{\mathrm{\Gamma }_{lm}^2(r_lr_m)}{(n_ln_m)\mathrm{\Omega }_4^2}\frac{\mathrm{\Gamma }_{gm}\mathrm{\Gamma }_{lm}}{\mathrm{\Gamma }_{gm}^2+(\mathrm{\Omega }_4\mathrm{\Omega }_1)^2}\frac{|G_1|^2(r_mr_g)}{\mathrm{\Omega }_4^2(n_ln_m)}.$$
$`(7)`$
The two last co-factors in (7) describe the Raman scattering and arise from the nominator (NIE) and denominator in (6). It ensues from (7) that the population inversion of initial and final unperturbed states ($`r_m=n_m>r_g`$) is required for the probing field amplification.
(ii) Resonance $`\mathrm{\Omega }_1=\mathrm{\Omega }_4=0`$.
The amplification and transparency conditions have the form
$$g_1\mathrm{\Delta }r_1\mathrm{\Delta }r_4,\frac{|G_1|^2}{\mathrm{\Gamma }_{lg}\mathrm{\Gamma }_{gm}}(r_lr_g)r_lr_m.$$
$`(8)`$
As it follows from (8), the amplification, due to NIE, does not require population inversion between initial and final states. The lower is the relaxation rate at a two-photon transition as compared to the coherence relaxation at coupled one-photon transitions, the more favorable are conditions for the inversionless amplification. An optimum strong field intensity is defined by the common level splitting into two quasilevels, which reduces the interference and, hence, the amplification at the $`ml`$ transition center. The population difference saturation at the strong field transition also reduces the system coherence. There is an optimum relationship between the initial population differences $`\mathrm{\Delta }n_4`$ and $`\mathrm{\Delta }n_3`$, created by an additional noncoherent excitation. Optimum conditions for the inversionless amplification and transparency for opened and closed systems are analyzed in \[28-30\] in more detail.
### 3.2 Three-level system in strong fields: inversionless amplification for the strong wave
Above expressions can be easily generalized to the case of inversionless amplification of the strong fields which can drive a quantum system. This case is of interest in connection with creation of ”laser without population inversion”. For certainty, let us consider the interaction of two strong fields $`E_3`$ and $`E_4`$ (Fig. 1). Taking the strong field $`E_4`$ effects in density matrix equations into account, the set of equations can be reduced to an algebraic. The solution has the form
$$r_4=i\frac{G_4}{P_4}\frac{(1+u_2^{})\mathrm{\Delta }r_4v_1\mathrm{\Delta }r_3}{1+v_4+u_2^{}},r_3=i\frac{G_3}{P_3}\frac{(1+v_4^{})\mathrm{\Delta }r_2u_3\mathrm{\Delta }r_4}{1+v_4^{}+u_2},$$
$`(9)`$
where
$$u_2=|G_4|^2/P_3P_{43}^{},u_3=|G_4|^2/P_4^{}P_{43}^{},$$
other notations are the same.
It is seen comparing (4) and (9) that, apart from the population difference saturation, a growth in the amplified wave intensity makes more difficult to achieve the conditions for inversionless amplification and self-transparency at the line center (factors $`(1+u_2^{})`$ and $`(1+v_4^{})`$ in the nominators), as well as reduces the gain due to an additional resonance splitting (the additional term in denominators). An extended analysis of this problem, accounting for the saturated populations, will be published elsewhere.
### 3.3 Inversionless amplification and resonantly amplified refraction in the absence of absorption in sodium vapor: a simple experiment
Currently a small number of experiments contrasts to a flux of theoretical publications. The most experiments concern with coherent excitation of a doublet or a set of neighboring sublevels in the short-pulse mode, as well as with accompanying interference effects. In \[41-43\] there was proposed a design, comprehensive theoretical grounds, and estimations for a possible experiment on simultaneous observation of inversionless amplification and absorptionless resonant refraction. This was a scheme of interfering two-quantum transitions induced by an auxiliary field in the uniformly broadened three-level system with collisions. Such an experiment is of interest due to minimized accompanying processes. Meanwhile, this simplest model enters more complex experimental schemes.
Let us consider again the energy level diagram displayed in Fig. 1. We assume the level $`n`$ to be a ground one and the field $`E_1`$ and $`E_4`$ to be absent. Thus, we separate the $`V`$-shaped three-level configuration $`gnm`$. The strong field $`E_3`$ at frequency $`\mathrm{\Omega }_3`$ couples the levels $`m`$ and $`n`$. The weak field at frequency $`\omega _2`$ probes the transition $`gn`$. Using (3) and (4) we derive the absorption $`\alpha _2`$ and refraction $`n_2`$ indices at frequency $`\omega _2`$, (see also \[28-30\]),
$$\alpha _2(\mathrm{\Omega }_2)=\alpha _2^0(0)\mathrm{Im}f(\mathrm{\Omega }_2,|E_3|^2),$$
$`(10)`$
$$\mathrm{\Delta }n(\mathrm{\Omega }_2)=n(\mathrm{\Omega }_2)n(\mathrm{\Omega }_2)^{(nr)}=\delta n_2^0\mathrm{Re}f(\mathrm{\Omega }_2,|E_3|^2,$$
$`(11)`$
$$f(\mathrm{\Omega }_2,|E|^2)=i\frac{\mathrm{\Gamma }_{gn}}{\mathrm{\Delta }n_{ng}}\frac{\mathrm{\Gamma }_{gm}+i(\mathrm{\Omega }_2\mathrm{\Omega }_3)(r_nr_g)ir_{mn}G_3}{[\mathrm{\Gamma }_{gn}+i\mathrm{\Omega }_2][\mathrm{\Gamma }_{gm}+i(\mathrm{\Omega }_2\mathrm{\Omega }_3)]+|G_3|^2},$$
$`(12)`$
where $`\alpha _2^0(0)`$ is the absorption (or amplification, depending on the population difference $`\mathrm{\Delta }n_2`$ sign) at the spectral line center if the strong field $`E_3`$ is turned off, $`\delta n_2^0`$ is the maximum contribution of transition $`ng`$ to the refraction index at $`E_3=0`$, and $`n(\mathrm{\Omega }_2)^{nr}`$ is the linear contribution of all other nonresonant levels.
As was noted above, the NIE leads to inversionless amplification and is created there by the coherence at the transition $`gm`$. The coherence is induced by the strong field (factor $`r_{mn}`$) in combination with the probing field. The greater is $`|G_3r_{nm}/\mathrm{\Gamma }_{gm}|`$ as compared to $`r_nr_g`$, the more pronounced is the effect,
$$r_{mn}=\frac{iG_3(r_nr_m)}{\mathrm{\Gamma }i\mathrm{\Omega }_3},$$
$`(13)`$
hereafter $`\mathrm{\Gamma }\mathrm{\Gamma }_3`$. At $`\mathrm{\Omega }_3=0`$ the absorption (amplification) maximum corresponds to $`\mathrm{\Omega }_2=0`$, hence,
$$f(0)=\frac{r_nr_g(r_nr_m)|G_3|^2/\mathrm{\Gamma }\mathrm{\Gamma }_{gm}}{(1+|G_3|^2/\mathrm{\Gamma }_{gm}\mathrm{\Gamma }_{gn})\mathrm{\Delta }n_2}.$$
$`(14)`$
Thus, even at $`(r_nr_g)>0`$ and $`(r_nr_m)>0`$, a negative absorption, i.e. amplification, could take place, if
$$|r_nr_m||G_3|^2/\mathrm{\Gamma }\mathrm{\Gamma }_{gm}>|r_nr_g|.$$
$`(15)`$
The lower is the coherence relaxation rate $`\mathrm{\Gamma }_{gm}`$ at two-photon transition $`gm`$ as compared to the coherence relaxation at coupled one-photon transitions, the more favorable are conditions for inversionless amplification. At $`|G_3|^2\mathrm{\Gamma }_{gm}\mathrm{\Gamma }_{gn}`$ a splitting of the level $`n`$ into two quasi-levels significantly reduces interference and, hence, amplification at the transition $`gn`$ center. There is also an optimum relationship between saturated population differences at the interacting transitions. It depends on the strong field intensity and the relation between initial population differences $`\mathrm{\Delta }n_2=n_nn_g`$ and $`\mathrm{\Delta }n_3=n_nn_m`$ \[28-30\], created by an additional noncoherent radiation. To vary this relation in a wide range, we proposed in \[41-43\] to use alkali atoms placed into a high-pressure buffer gas. The strong field couples $`P_{3/2}`$ and the ground $`S`$ \- level. A fast collisional exchange furnishes population transfer from $`P_{3/2}`$ to the lower level $`P_{1/2}`$. For simplicity, it can be believed that the population distribution over the fine structure levels is Boltzmann’s one due to collisions. Thus, it becomes possible to vary the population difference at the probing transition in a wide range varying the strong field intensity and buffer gas pressure. Due to the saturation of $`P_{3/2}S`$ transition, even the population inversion at the $`P_{1/2}S`$ becomes possible (similar to a ruby laser). Hence, $`P_{1/2}S`$ can be chosen as a probing transition. The population inversion was experimentally observed by a similar scheme in the mixture of sodium and helium vapors .
Collisions play a double part, i.e., on the one hand they considerably worsen coherence, on the other hand the population transfer due to collisions furnishes simple control and optimization of the population differences at the coupled transitions. Moreover, a wide collisional broadening allows one to neglect the nonuniformity of interaction with atoms due to the Doppler effect, hyperfine splitting, and some other processes. This makes an experiment be governed by the simplest theoretical model. By estimations and numerical examples for sodium atoms, now we show that inversionless amplification and absorptionless resonant refraction can be significant under proposed experimental conditions, in spite of collisions reducing the coherence.
We will discuss the concrete $`D1`$ and $`D_2`$ transitions in sodium. According to , we write the kinetic equations of level populations
$$(\mathrm{\Gamma }_m+\nu _{mg})r_m\nu _{gm}r_gP=0,$$
$$P=\frac{2|G_3|^2\mathrm{\Gamma }(r_nr_m)}{\mathrm{\Gamma }^2+i\mathrm{\Omega }_3^2},$$
$`(16)`$
$$(\mathrm{\Gamma }_g+\nu _{gm})r_g\nu _{mg}r_m=0,\mathrm{\Gamma }_mr_m+\mathrm{\Gamma }_gr_gP=0,$$
$`(17)`$
$$r_m+r_n+r_g=N,$$
$`(18)`$
where $`\nu _{gm}`$ and $`\nu _{mg}`$, are the frequencies of collisions transferring populations, $`\mathrm{\Gamma }_g^1`$ and $`\mathrm{\Gamma }_m^1`$ are the lifetimes of relevant levels. From (16)-(18) we find
$$r_nr_m=\frac{N}{1+\text{æ}},r_nr_g=\frac{N}{1+\text{æ}}\left[\text{æ}\frac{\nu _{mg}(\nu _{gm}+\mathrm{\Gamma })}{\nu _{mg}+2(\nu _{gm}+\mathrm{\Gamma })}1\right],$$
$`(19)`$
where
$$\text{æ}=\frac{2|G_3|^2\mathrm{\Gamma }}{\mathrm{\Gamma }^2+\mathrm{\Omega }_3^2}\frac{\nu _{mg}+2(\nu _{gm}+\mathrm{\Gamma })}{\mathrm{\Gamma }_g\nu _{mg}+\mathrm{\Gamma }_m(\nu _{gm}+\mathrm{\Gamma })}.$$
$`(20)`$
Let us assume that $`\mathrm{\Gamma }_g\mathrm{\Gamma }_m`$ and a buffer gas pressure is high that $`(\nu _{mg}\nu _{gm})\mathrm{\Gamma }_{gm}`$. Taking under given conditions $`\nu _{gm}=\nu _{mg}\mathrm{exp}(\mathrm{\Delta }E/k_bT)`$, where $`\mathrm{\Delta }E=E_mE_g`$, is the fine splitting energy, $`k_b`$ is the Boltzmann constant, and T is the temperature, we get
$$r_nr_g=\frac{N}{1+\text{æ}}\left[\text{æ}\frac{1\mathrm{exp}(\mathrm{\Delta }E/k_BT)}{1+2\mathrm{exp}(\mathrm{\Delta }E/k_BT)}1\right],$$
$`(21)`$
$$\text{æ}=\frac{2|G_3|^2\mathrm{\Gamma }}{\mathrm{\Gamma }_m(\mathrm{\Gamma }^2+\mathrm{\Omega }_3^2)}\frac{1+2\mathrm{exp}(\mathrm{\Delta }E/k_BT)}{1+\mathrm{exp}(\mathrm{\Delta }E/k_BT)}.$$
$`(22)`$
For sodium $`\mathrm{\Delta }E=17.2`$ cm$`^1`$ and, at T = 550 K, estimations yield $`\mathrm{\Delta }E/k_bT=4.310^2`$, $`\text{æ}3|G_3|^2/\mathrm{\Gamma }\mathrm{\Gamma }_m9\lambda ^3I/64\pi ^3ϵ_0\mathrm{}c\mathrm{\Gamma }`$,
$$r_nr_g=\frac{N}{1+\text{æ}}\left[1.310^2\text{æ}1\right].$$
$`(23)`$
where $`\lambda `$ and $`I`$ are the strong field wavelength and energy flux density, $`ϵ_0`$ is the dielectric constant of vacuum. From (13), (14), (20), and (21), it is seen that principally attainable inversionless amplification rises as $`\mathrm{\Delta }E`$ grows (e.g., in $`K`$ and $`Rb`$).
The inelastic collision cross section in sodium and helium for the transition $`3P_{3/2}3P_{1/2}`$ is $`\sigma _{mg}410^{15}`$ cm<sup>2</sup>. At T=550 K and helium atmospheric pressure, estimations yields $`\nu _{mg}=N_{He}\overline{v}\sigma _{mg}7.510^9`$ s<sup>-1</sup>. Since $`\mathrm{\Gamma }_g\mathrm{\Gamma }_m6.210^7`$ s<sup>-1</sup>, the validity conditions for approximation (21) are satisfied. Using data for the collisional broadening D of sodium and helium lines, we estimate the collisional halfwidth as $`\mathrm{\Gamma }510^{10}`$ s<sup>-1</sup>, which exceeds the Doppler’s width of this transition, $`\mathrm{\Delta }\omega _D/2=4.710^9`$ s<sup>-1</sup> ($`\mathrm{\Delta }\nu _D/2=0.75`$ GHz). For our conditions, we have
$$\frac{|G_3|^2}{\mathrm{\Gamma }\mathrm{\Gamma }_{gm}}\frac{|G_3|^2}{\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm}}\frac{\text{æ}\mathrm{\Gamma }_m}{3\mathrm{\Gamma }_{gm}},$$
$`\text{æ}510^9I`$, where $`I`$ is expressed in $`W\text{cm}^2`$.
For the radiation power of 0.1 W focused into the footprint $`A=10^5`$ cm<sup>2</sup> (the confocal parameter is $`b1`$ cm), we find $`|G_3|3.6`$ GHz, $`\text{æ}510^2`$, and $`|G_3|^2/\mathrm{\Gamma }\mathrm{\Gamma }_{gm}0.1`$. These values are optimal to vary the population difference around zero at the probing transition $`r_nr_g`$. Above estimations for the intensity 1-10 kW$``$cm<sup>2</sup> required to change noticeably the line shape agree well to the experiment with a change in the population difference ratio at the coupled transitions (44\]. Inversionless amplification at $`r_nr_g=0`$ is estimated as $`\alpha _2(0)/\alpha _2^0(0)\mathrm{\Gamma }_m/3\mathrm{\Gamma }_{gm}`$. Accepting $`\mathrm{\Gamma }_{gm}\nu _{mg}`$, we find a value about 0.3% of the absorption in the absence of strong field. It is seen that this value is very sensitive to the coherence decay rate at the transition $`gm`$.
Absorption-amplification spectral line shapes, frequency interval positions and halfwidths can be controlled, as analyzed in . The line shape is very sensitive to the strong field intensity and frequency detunings from resonance. The amplification halfwidth increases and maximum decrease as the strong field intensity grows above a certain value. The population difference saturation at the strong field transition and the common energy level splitting reduce the inversionlesa amplification, so that it should be optimized by an appropriate choice of the strong field intensity and detuning. In our case the inelastic collision frequency is the important optimization parameter. The refraction index (dispersion) is described by $`\text{Re}f(\mathrm{\Omega }_2,|E_3|^2)`$. In the framework of proposed experiment the absorption-amplification coefficient and refraction index shapes can be controlled so that the refraction maximum falls within the spectral interval of vanishing absorption.
Thus, the considered model of three-level system with the interference controlled by collisions makes it possible the amplification without population inversion and the resonantly amplified refraction at vanishing absorption. Exact formulas are presented to analyze optimum experimental conditions. Collisions destroying the coherence reduce the effects, as compared to atomic beams. However, this decrease is comparable to the effect of Doppler broadening in metal vapors. Advantages of the proposed experimental design are a simplicity and the possibility to control populations at the coupled transitions and to avoid interfering effects. This makes experiment adequate to a simple theoretical model. An experiment on inversionless amplification in the continuous mode in potassium vapor using collisional population of the upper level of probing transition was carried oat in , however the four-level configuration contribution was significant there. The latter was conditioned by the transitions between hyperfine splitting sublevels and the incomplete overlapping of optical transitions due to insufficient collisional broadening.
## 4 Coherence and frequency mixing: multiple resonances at the condition of induced transparency
A nonlinear-optical response sharply increases as interacting wave frequencies approach one - and multiphoton resonances. This reduces required intensities of initial fields down to the values corresponding to cw lasers \[47-51\]. However, due to resonant absorption of primary and generated waves, there arise limitations from above onto the atomic density. Quantum coherence alternatively manifests itself in various optical processes. In particular, as was shown by an example of bound-free transitions , the absorption decrease can be not accompanied by an effective nonlinear susceptibility decrease on frequency mixing and varies the refraction index in another way.
Recently, an interest is growing to control matter - optical properties via the quantum coherence effects, especially promising for shortwave generation . Therewith, an accent is on the wave conversion at frequency mixing under the condition of resonance with an absorbing transition between discrete levels and only by the generated field. In a scheme of totally resonant multiphoton interaction was proposed, in which the quantum transition coherence and interference suppress absorption of both primary and generated field. Therewith the atomic nonlinear susceptibility is not subject to a significant destructive interference and rises by many orders of magnitude due to simultaneous multiphoton and one-photon resonances. Now available great atomic concentration rises additionally the nonlinear-optical response of medium and yields new spectral dependencies conditioned by local field effects (to be considered further).
In this section we consider several qualitative effects conditioned by intratomic coherence and possible totally resonant four-wave interaction under low absorption of both generated and initial fields. We show also the possibility of significant efficiency of generation. The results are easily extended to nonlinear-optical processes of higher order. Now we turn to the energy level diagram shown in Fig.2a .
The strong fields $`E_3`$ and $`E_2`$ at frequencies $`\omega _3`$ and $`\omega _2`$ couple nonpopuiated levels 3, 2 and 2, 1, respectively. The fields $`E_1`$ and $`E_s`$, generated at frequencies $`\omega _1\omega _{10}`$ and $`\omega _s=\omega _1+\omega _2+\omega _3`$ are assumed to be weak and not changing the level populations. The latter fields are considered only in the lowest order of perturbation theory. The absorption coefficient and refraction index at frequencies $`\omega _1`$ and $`\omega _s`$, as well as the nonlinear polarization generating the wave at frequency $`\omega _s`$ are defined by real and imaginary parts of effective linear
$$\chi _1(\omega _1;\omega _1)=(\chi _1^0/P_{01})f_1,\chi _s(\omega _s;\omega _s)=(\chi _s^0/P_{03})f_s,$$
$`(24)`$
and nonlinear
$$\chi ^{NL}(\omega _s;\omega _1+\omega _2+\omega _3)=(\chi _0^{NL}/P_{01}P_{02}P_{03})f$$
$`(25)`$
susceptibilities, which, in turn, are proportional to the pre-exponentiat factors $`r_i`$, and $`\stackrel{~}{r}_i`$ of the corresponding components of nondiagonal density matrix elements (see the similar equations of Sec. 2). Here $`\chi _1^0`$, $`\chi _s^0`$ and $`\chi _s^{NL}`$ – are the resonant susceptibilities at negligibl $`G_2`$ and $`G_3`$. The factors $`f_1`$, $`f_2`$, and $`f`$ describe the strong field effects. Simple calculations by the density matrix procedure similar to yield
$$f_1=\left\{1+g_2/P_{01}P_{02}[1+(g_3/P_{02}D_{03})]\right\}^1,$$
$`(26)`$
$$f_s=\left\{1+g_3/P_{03}D_{02}[1+(g_2/D_{02}D_{01})]\right\}^1,$$
$`(27)`$
$$f=f_1[1+g_3/D_{03}P_{02}]^1=\left[1+(g_2/D_{02}D_{01})+(g_3/D_{03}P_{02})\right]^1,$$
$`(28)`$
where
$$P_{01}=1+ix_1,P_{02}=1+ix_{02},P_{03}=1+ix_s,$$
$$D_{01}=1+iy_1,D_{02}=1+iy_{02},D_{03}=1+iy_s,$$
$$x_1=\frac{\omega _1\omega _{10}}{\mathrm{\Gamma }_{10}}=0,x_{02}=\frac{\omega _1+\omega _2\omega _{20}}{\mathrm{\Gamma }_{20}}=0,x_s=\frac{\omega _s\omega _{30}}{\mathrm{\Gamma }_{30}}=0,$$
$$y_1=\frac{\omega _s\omega _3\omega _2\omega _{10}}{\mathrm{\Gamma }_{10}}=0,y_{02}=\frac{\omega _s\omega _3\omega _{20}}{\mathrm{\Gamma }_{20}}=0,y_s=\frac{\omega _1+\omega _2+\omega _3\omega _{30}}{\mathrm{\Gamma }_{30}}=0,$$
$$g_2=G_2^2/\mathrm{\Gamma }_{10}\mathrm{\Gamma }_{20},g_3=G_3^2/\mathrm{\Gamma }_{30}\mathrm{\Gamma }_{20},$$
and $`\mathrm{\Gamma }_{ij}`$ are the uniform halfwidths of corresponding transitions. To analyze the cases. when $`E_s`$, is not an independent probing field, we should put $`\omega _s=\omega _1+\omega _2+\omega _3`$ and $`D_{0i}=P_{0i}`$.
Differences of the factors $`f_1`$, $`f_2`$, and $`f`$ from unity and between each other, as well as their frequency dependency, are conditioned by two different coherence initiation channels $`\rho _{02}`$ (two combinations of strong and weak fields, $`E_1`$, $`E_2`$ and $`E_s`$, $`E_3`$) and their evolution as the field intensities rise. The absorption coefficients are defined by imaginary parts of the relevant susceptibilities, relative to which the considered processes manifest themselves as resonance splitting and absorption minima. The generated wave power $`Pg_2g_3|\chi ^{NL}|^2`$ is defined by not only an imaginary part, but also by the real part of $`\chi ^{NL}`$. (A relative phase of the generated wave depends on their ratio.) Therefore, the quantum coherence effects can be used to match the most important generation conditions, i.e., significant decrease of absorption for all the interacting fields without a noticeable decrease in atomic nonlinear-optical response. Furthermore, these effects can be used to increase additionally the generation efficiency improving the wave phase velocity matching. The laser-induced spectral structures in the susceptibility real parts (additional dispersion conditioned by the coherence) enable such a matching by a slight detuning from the resonance for $`\omega _1`$ or $`\omega _s`$.
An approach to the two resonances increases $`|\chi ^{NL}|^2`$ by the factor $`x_i^2`$, equal to $`10^6`$ and more. The triple resonance gives a gain of the order of $`10^{18}`$. Due to the induced transparency, the atomic density N and, hence, the power $`PN^2`$ can be additionally increased by a few orders of magnitude.
These results are easily generalized to the case of higher order mixing. For instance, if the 3-2 and (or) 2-1 transitions are multiphoton, the generalization is achieved by substituting the one-photon Rabi frequencies $`G_2`$ and $`G_3`$ with the corresponding matrix elements for multiphoton transitions.
In a particular case, when the transparency is induced only at the frequency $`\omega _s`$, by the wave $`E_4`$ not directly participating in conversion (Fig. 2c), the relevant expressions are also deduced by a simple subscript substitution. It is necessary to put $`g_2=0`$ and change subscripts in the formulas for $`g_3`$ and the corresponding resonance denominator. This case is similar to considered in for the transition to continuum (Fig. 26) (see also ).
Thus, the intratomic coherence effects in strong electromagnetic fields enable one to control absorption, refraction, and nonlinear-optical generation spectral properties. In particular, the above choice of transition diagram and interacting wave intensities make it possible to gain a medium nonlinear-optical response increased by many orders, meanwhile eliminating the initial and generated field absorption.
## 5 Atomic coherence: effects of a local field and the spectral line nonuniform broadening
### 5.1 Local field
The local field acting on an atom substantially differs from the external field both in value and phase, as the atomic density grows. This varies shapes of the spectral lines conditioned by the quantum transitions’ interference . The effect is revealed at particle concentration of the order of $`10^{17}`$ cm<sup>3</sup>. Let us show by a simple example that qualitatively new spectral dependencies can arise in above problems due to the local field formation and drastic variation in the field of additional strong laser radiation at an adjacent transition.
According to a conventional (but approximate) concept (see, e.g. ), the local $`E_L`$ and external $`E`$ fields in isotropic media are related by the simple formula
$$E_L=E+\frac{P}{3ϵ_0}.$$
$`(29)`$
The medium polarization $`P`$ to a linear approximation can be presented by $`P=ϵ_0N\alpha E_L`$, where $`N`$ is the atomic concentration, $`\alpha `$ is the microscopic (atomic) polarizability, and $`ϵ_0`$ is the dielectric constant of vacuum.
One of the consequences from (29) is the Clausius-Mossotti equation relating polarizability $`\alpha `$ to the dielectric constant $`ϵ_0`$ of material
$$ϵ=1+LN\alpha ,$$
$`(30)`$
where the local field factor $`L=(ϵ+2)/3=(1\alpha N/3)^1`$ shows how the local field differs from the external one. The former plays an important part in linear and nonlinear optical phenomena (see, e.g. \[56-59\]. In spite of an approximate nature of formula (29), the authors of showed that it well describes linear and nonlinear responses of the dense atomic gas.
Let us consider the interaction of two optical fields with three-level systems of the cascade or $`\mathrm{\Lambda }`$-schemes. Therewith, one field is strong and (for simplicity) interacts with the transition between nonpopulated or equipopulated levels. According to the classification of resonant nonlinear processes (see Sec. 2), only the resonance splitting effect would be observed for the probing field in this case. For certainty we consider the $`nml`$ $`\mathrm{\Lambda }`$-scheme (Fig. 1), where the levels $`m`$ and $`n`$ are nonpopulated, but the state $`l`$ is ground. The strong field of frequency $`\omega _3`$ and amplitude $`E_3`$ acts at the $`nm`$ transition, the weak (probing), wave of frequency $`\omega _4`$ and amplitude $`E_4`$ acts at the $`lm`$ transition. The set of equations for this problem differs used in Sec. 3 only by the substitution of $`E_4`$ by $`E_{4L}`$. Such an approach is widely used in the local field theory \[56-59\].
For the macroscopic complex polarization $`P_4(\omega _4)=d_{ml}\rho _{lm}N`$ using formulas of Sec. 3 we find
$$P_4(\omega _4)=ϵ_0\chi _4(\omega _4)E_4,\chi _4(\omega _4)=\chi _4^0f(\omega _4),$$
$`(31)`$
$$\chi _4^0=i\frac{N|d_{ml}|^2}{ϵ_0\mathrm{\Gamma }_{lm}},f(\omega _4)=\frac{\mathrm{\Gamma }_{lm}P_{43}}{(P_4i\delta _{4L})P_{43}+|G_4|^2},$$
$`(32)`$
where $`\chi _4`$ is the macroscopic susceptibility at the probing field frequency $`\omega _4`$ in the presence of strong one at the frequency $`\omega _3`$, $`\chi _4^0`$ is the susceptibility in the absence of the latter, and $`f(\omega _4)`$ is the form-factor. The parameter
$$\delta _{4L}=|d_{ml}|^2\frac{N}{3ϵ_0\mathrm{}}$$
$`(33)`$
appears as the transition $`lm`$ frequency shift conditioned by the concentration rise (local field). It is substantial that in so doing the two-photon and strong field transition frequencies are not varied. As a result, the local field effect is not reduced to redefining the detuning of the weak field resonance, but qualitatively changes the whole spectral line shape, if this factor becomes comparable to the resonance width. At $`d_{ml}`$=1 Db and $`N=10^{23}`$ m<sup>-3</sup> we estimate this shift as $`\delta _{4L}=810^{11}`$ s<sup>-1</sup>, which can be comparable to characteristic shock widths of resonances. For instance, when the shock width is defined by resonance exchange (self-broadening) and significantly exceeds a natural one, $`\mathrm{\Gamma }_{ml}`$ has the form (see, e.g. (58,59\] and references therein)
$$\mathrm{\Gamma }_{ml}|d_{ml}|^2\frac{N}{6ϵ_0\mathrm{}},$$
whence it follows that the ratio $`C_4=\delta _{4L}/\mathrm{\Gamma }_{ml}`$ can approach two in this case.
Peculiarities of the local field, exhibited in spectral nonlinear interference dependencies of absorption, are analyzed in . The red shift, entering all the resonant denominators, and hence, not leading to a simple redefinition of the detuning from one-photon resonance, changes qualitatively the spectral dependencies as the absorbing particle concentration rises. Thus, the local field strongly varies spectral dependencies of the probing field absorption in the presence of a strong field at the adjacent transition.
The factor $`L_4=E_{4L}/E_4`$ characterizing the difference between local and external fields by value and phase can be presented in the form
$$L_4=1+iG_4f(\omega _4).$$
$`(34)`$
Whence it follows that this difference also increases as the parameter $`C^4`$ rises. Varying the strong field intensity and frequency, the local field spectral dependence can be also significantly varied.
The resonant exchange also shifts the frequency of the groond-to-excited state transition. This shift is proportional to atomic concentration, but is usually two- or threefold shorter than the broadening and can be often neglected.
The acquired results are genual and applicable to other interaction schemes. In the case of cascade transitions under the conditions of zero or equal level populations at the strong field transition, the results are found by a simple redefinition of corresponding quantities.
Since the local field acting on an atom can substantially differ by value and phase from the external field as atomic density rises, this drastically changes spectral dependencies not only for absorption and refraction, but also for generating nonlinear poiarizations . Similar, it can be shown using the Lorentz-Lorentz approximation that the local field induces red shifts in the resonant denominators at the allowed transitions for effective nonlinear susceptibilities (Sec. 4). Thus, a supplementary opportunity appears to increase efficiency sharply by controlling concentration, initial field intensities, and detuninga from resonances at the conditions of multiple resonances, induced transparency, and phase matching. The considered local field effects should be taken into account when designing and interpreting experiments.
### 5.2 Constructive and destructive interference as a consequence of atomic velocity distribution: nonuniform broadening of spectral lines
As it was already mentioned, the laser-induced coherence contribution into spectra can be constructive or destructive depending on a detuning sign. Therefore, inversionless amplification and resonance splitting in gases can significantly differ from those for stationary atoms at the Doppfer broadening prevailing above uniform one. Nevertheless, as was shown in , spectral profiles with alternating signs are induced sometimes even in this case. In the relatively weak laser fields the NIE manifest itself in narrow spectral structure variations within the wide Doppler’s contour. This structure is anisotropic, i.e., depends on the angle between the strong and probing field wave vectors, as well as on collisions changing atom velocities. Destructivity or constructivity effect of the Maxwell velocity distribution depends on the position of probing wave frequency compared to the strong-field frequency. The probing field transition line is being deformed as a whole as the strong field transition uniform broadening or the wave intensity rise. Formulas of averaging over velocities for a number of cases are presented in .
Since the dependence of responses on detuning signs often has a sign-alternating character at coherent interaction, some interfering components can vanish when averaged over velocities. We present an example showing the effect of nonuniform broadening on coherent interaction. Let us consider the four-wave mixing in the lowest order of perturbation theory (Fig. 1). From the solutions to equations for $`\overline{r}_4`$ (Sec. 2) we deduce the nonlinear susceptibility at frequency $`\omega _4=\omega _1\omega _2+\omega _3`$
$$\chi (\omega _4=\omega _1\omega _2+\omega _3)=\frac{iK}{\mathrm{\Gamma }_{ml}+i(\mathrm{\Omega }_1^{^{}}\mathrm{\Omega }_2^{^{}}+\mathrm{\Omega }_3^{^{}})}\times $$
$`(35)`$
$$\left\{\frac{1}{\mathrm{\Gamma }_{gm}+i(\mathrm{\Omega }_3^{}\mathrm{\Omega }_2^{})}\left[\frac{n_gn_n}{\mathrm{\Gamma }_{ng}i\mathrm{\Omega }_2^{}}+\frac{n_mn_n}{\mathrm{\Gamma }_{mn}+i\mathrm{\Omega }_3^{}}\right]+\frac{1}{\mathrm{\Gamma }_{ln}+i(\mathrm{\Omega }_1^{}\mathrm{\Omega }_2^{})}\left[\frac{n_gn_n}{\mathrm{\Gamma }_{ng}i\mathrm{\Omega }_2^{}}+\frac{n_gn_l}{\mathrm{\Gamma }_{lg}+i\mathrm{\Omega }_1^{}}\right]\right\}$$
where $`K`$ is the constant, $`\mathrm{\Omega }_1^{}=\omega _1\omega _{gl}𝐤_\mathrm{𝟏}𝐯=𝛀_\mathrm{𝟏}𝐤_\mathrm{𝟏}𝐯`$ is the detuning from resonance tailing the Doppler’s shift into account, other $`\mathrm{\Omega }_i^{}`$ are the similar detunings depending on velocity, and $`n_i`$, are the populations of corresponding levels, also depending on velocity.
As is seen from (35), all the terms, besides those proportional to $`n_gn_n`$ as functions of velocity, have poles at the same complex semiplane. Therefore, if the Doppler’s shifts corresponding to the heat velocity $`u`$ much exceed uniform halfwidths of transitions, then only the polarization components proportional to $`n_gn_n`$ are nonzero after averaging over velocities with the Maxwell distribution. The averaging result has the form
$$\chi _v=\frac{iK\pi ^{1/2}\mathrm{exp}\left\{(\mathrm{\Omega }_2/k_2u)^2\right\}(N_gN_n)}{k_2u[\stackrel{~}{\mathrm{\Gamma }}_1+i(\mathrm{\Omega }_1k_1\mathrm{\Omega }_2/k_2)][\stackrel{~}{\mathrm{\Gamma }}_3+i(\mathrm{\Omega }_3k_3\mathrm{\Omega }_2/k_2)]},$$
$`(36)`$
where $`N_g`$, and $`N_n`$ are the integral over velocities unperturbed level populations, $`\mathrm{\Omega }_4=\mathrm{\Omega }_1\mathrm{\Omega }_2+\mathrm{\Omega }_3`$, $`\mathrm{\Omega }_1=\omega _1\omega _{gl}`$, $`\mathrm{\Omega }_2=\omega _2\omega _{gn}`$, $`\mathrm{\Omega }_3=\omega _3\omega _{mn}`$, $`k_i=\omega _i/c`$ and
$$\stackrel{~}{\mathrm{\Gamma }}_1=\mathrm{\Gamma }_{nl}+(k_1/k_21)\mathrm{\Gamma }_{ng},\stackrel{~}{\mathrm{\Gamma }}_3=\mathrm{\Gamma }_{gm}+(k_3/k_21)\mathrm{\Gamma }_{ng}.$$
$`(37)`$
Substituting $`\mathrm{\Omega }_2^{}`$ by $`\mathrm{\Omega }_2^{}`$ in the sum $`\omega _4=\omega _1\omega _2+\omega _3`$ in the resonant cascade level configuration, we find that all the poles appears at the same complex plane. This significantly decreases the averaged susceptibility as compared to the frequency subtraction scheme.
The effect of velocity distribution on the coherent four-wave mixing can be described by the following way. For stationary atoms and a totally resonant process, the nonlinear polarization is proportional to the factor $`\mathrm{\Gamma }^3`$. For a gas and the frequency subtraction at an available intermediate level population, this factor is substituted by $`1/ku\mathrm{\Gamma }^2`$, i.e., decreased by $`ku/\mathrm{\Gamma }`$. A more detailed analysis shows that for frequency summation or at $`N_nN_g=0`$, the coherence and nonlinear polarization suppressed by the interference of various velocity contributions yields this factor equal to $`(ku)^3`$. In other words, the susceptibility decreases by the factor of $`(ku/\mathrm{\Gamma })^3`$ as compared to stationary atoms and by the factor of $`(ku/\mathrm{\Gamma })^2`$ as to optimum conditions for frequency subtraction in gases. Hence. the difference scheme was chosen for continuous four-wave generation in the field of helium-neon laser in as distinct from .
When only one resonance ($`\omega _1\omega _2\omega _{ln}`$) of the Raman-type scattering presents, we have
$$\chi ^{(3)}\frac{1}{\mathrm{\Omega }_1\mathrm{\Omega }_4}\mathrm{exp}\left[\left(\frac{\mathrm{\Omega }_1\mathrm{\Omega }_2}{(k_1k_2)u}\right)^2\right].$$
In the absence, of resonancea (for both frequency summation and subtraction), we find
$$\chi ^{(3)}1/\mathrm{\Omega }_1\mathrm{\Omega }_4(\mathrm{\Omega }_1\mathrm{\Omega }_2).$$
Since the NIE and mixing processes are defined by a common one source, i.e., the coherence induced at a forbidden transition, dose effects occur also at averaging over velocities.
Thus, the transition nonuniform broadening can significantly change the influence of various level population on the coherent processes, so that some level contributions prevail. Varying the level populations, strong field intensities and detunings, the amplification-absorption, refraction, and nonlinear polarization spectra can be controlled, thus considerably increasing che generation yield.
## 6 Nonlinear interference effects at bound-free transitions: laser-induced autoionization-type resonances in the continuum
Nonlinear interference effects similar to those occurring at the discrete transitions (including inversionless amplification and induced transparency) manifest themselves also in continuous spectra, e.g., at the transition into an ioniza-tion continuum. The corresponding theory was generalized in \[30-32, 36, 62\]. Similar phenomena for crystal brands were considered in . The laser-induced and autoionization-type resonances theoretically analyzed in were experimentally observed in , and then the nonlinear interference processes became a subject of intense studies in the context of laser-induced continuum structures (LICS), inversionless amplification and electromagnetically induced transparency, first, at bound-free transitions (see, e.g. ) and then also at discrete transitions (see \[7-11, 40, 52, 53\]).
Let us show the potentials to control simultaneously two LICS and to split discrete resonances by strong electromagnetic fields to reduce absorption, correct phase matching, and to improve nonlinear-optical short-wave generation techniques (Fig. 2b). The wave at frequency $`\omega _1`$ is weak, but those at frequencies $`\omega _3`$, and $`\omega `$ are strong. We also take into account possible strong nonresonant transitions to the discrete levels $`k`$. The detunings $`|\omega _1\omega _{gm}|`$, $`|\omega _1+\omega _2\omega _{gn}|`$ and $`|\omega _1\omega _3\omega _{ln}|`$ are assumed to be much smaller than all other. Density matrix calculations similar to yield the following expressions for nonlinear susceptibility $`\chi ^{(3)}(\omega _\mu =\omega _1+\omega _2+\omega _3)`$, the probing wave absorption coefficients $`\alpha (\omega _1)`$ and $`\alpha (\omega _\mu )`$ at the corresponding frequencies,
$$\frac{\chi ^{(3)}(\omega _\mu =\omega _1+\omega _2+\omega _3)}{\chi _{0\mu }^{(3)}}=\frac{K}{D_{gm}X},$$
$`(38)`$
$$\frac{\alpha (\omega _1)}{\alpha _{01}}=\mathrm{Re}\left\{\frac{1g_{mn}/(D_{gm}X)}{D_{gm}}\right\},$$
$`(39)`$
$$\frac{\alpha (\omega _\mu )}{\alpha _{0\mu }}=1k_3\beta _l+\frac{k_3\beta _l(y_l+q_{gl})^2}{1+y_l^2}\mathrm{Re}\left\{k_4g_{nn}A^2\frac{(1iq_{gn})^2}{Y}\right\},$$
$`(40)`$
where $`\chi _{0\mu }^{(3)}`$, $`\alpha _{01})`$ and $`\alpha _{0\mu })`$ are the relevant resonant values for negligible intensities of all the fields,
$$K=1\frac{k_1\beta _l[(1iq_{nl})(1iq_{lg})]}{(1iq_{ng})(1+ix_l)},$$
$`(41)`$
$$A=1\frac{k_1\beta _l[(1iq_{ln})(1iq_{gl})]}{(1iq_{gn})(1+iy_l)},$$
$`(42)`$
$$X=(1+g_{nn})\left[1+ix_n+\frac{q_{mn}}{D_{gm}(1+q_{nn})}k_2\beta _l\beta _n\frac{(1iq_{nl})^2}{1+ix_l}\right],$$
$`(43)`$
$$Y=(1+g_{nn})\left[1+iy_n+\frac{q_{mn}}{p_{gm}(1+q_{nn})}k_2\beta _l\beta _n\frac{(1iq_{nl})^2}{1+iy_l}\right],$$
$`(44)`$
$$D_{gm}=1+\frac{i(\omega _1\omega _{gm})}{\mathrm{\Gamma }_{gm}},p_{gm}=1+\frac{i(\omega _\mu \omega _3\omega _2\omega _{gm})}{\mathrm{\Gamma }_{gm}},$$
$`(45)`$
$$x_l=\frac{\omega _1+\omega _2+\omega _3\omega \omega _{gl}\delta _{ll}}{\mathrm{\Gamma }_{gl}+\gamma _{ll}},x_n=\frac{\omega _1+\omega _2\omega _{gn}\delta _{nn}}{\mathrm{\Gamma }_{gn}+\gamma _{nn}},$$
$`(46)`$
$$y_l=\frac{\omega _\mu \omega \omega _{gl}\delta _{ll}}{\mathrm{\Gamma }_{gl}+\gamma _{ll}},y_n=\frac{\omega _\mu \omega _3\omega _{gn}\delta _{nn}}{\mathrm{\Gamma }_{gn}+\gamma _{nn}},$$
$`(47)`$
$$k_1=\frac{\gamma _{gl}\gamma _{ln}}{\gamma _{gn}\gamma _{nn}},k_2=\frac{\gamma _{nl}\gamma _{ln}}{\gamma _{ll}\gamma _{nn}},k_3=\frac{\gamma _{gl}\gamma _{lg}}{\gamma _{gg}\gamma _{ll}},k_4=\frac{\gamma _{gn}\gamma _{ng}}{\gamma _{gg}\gamma _{nn}},$$
$`(48)`$
$$g_{mn}=\frac{|G_{mn}|^2}{\mathrm{\Gamma }_{gm}\mathrm{\Gamma }_{gn}},\beta _l=\frac{g_{ll}}{1+g_{ll}},\beta _n=\frac{g_{nn}}{1+g_{nn}},$$
$`(49)`$
$$g_{ii}=\frac{\gamma _{ii}}{\mathrm{\Gamma }_{gi}},g_{ij}=\frac{\delta _{ij}}{\gamma _{ij}},$$
$`(50)`$
$$\gamma _{ij}=\pi \mathrm{}G_{iϵ}G_{ϵj}|_{ϵ=\mathrm{}\omega \mu }+\mathrm{Re}\left\{\underset{k}{}\frac{G_{ik}G_{kj}}{p_{gk}}\right\},\delta _{ij}=\mathrm{}P\frac{dϵG_{iϵ}G_{ϵj}}{\mathrm{}\omega _\mu ϵ}+\mathrm{Im}\left\{\underset{k}{}\frac{G_{ik}G_{kj}}{p_{gk}}\right\},$$
$`(51)`$
$`P`$ in (51) designates the sign of the integral principal value. The factors $`k_i`$, take the values $`0k_i1`$ depending on a degree of degeneracy for continuum states (unity for the nondegenerate states).
Formulas (38) and (40) generalize the expressions from to the case of several strong fields. Together with (39) these formulas show the possibility to reduce absorption both for the primary and the generated fields. The absorption falls exponentially as the medium length rises. Since the absorption coefficients, as functions of frequency, and the ratios of squared module of nonlinear susceptibility to these coefficients, do not coincide at certain conditions. These ratios define the generation power in absorbing media. This power quadratically rises with growing atomic concentration and the initial fields intensity under the condition of optimized medium perturbation. Comparing (38) and (40) to corresponding formulas from and Sec. 4, we see new interference spectral structures caused by the joint action of strong fields $`E_3`$and $`E`$ (terms proportional to $`\beta _n`$ and $`g_n`$), arising in nonlinear polarization and absorption (refraction). These nonlinear resonances give additional opportunities for absorption nonlinear spectroscopy and enhancement of efficiency for nonlinear-optical conversion to the short-wave spectral range.
## 7 Nonlinear interference and relaxation
Relaxation processes can extraordinarily exhibit the coherence effects. For instance, the spontaneous relaxation coupling two transitions with close frequencies can promote inversionless amplification even in the absence of strong fields. Let us consider an example , when certain relaxation mechanisms and external dc fields suppressing the destructive interference make the resonant nonlinear-optical interaction allowed. The experiment was carried out with the $`HeNe`$ laser beam ($`\lambda =1.52`$ $`\mu `$m) resonant to the transition $`2s_22p_4`$ of Ne atoms. Upper and lower levels contain three ($`J_1=1)`$ and one ($`J_0=0`$) Zeeman sublevels, respectively. The initial beam contained two linear ortogonaily polarized components $`E_1`$ and $`E_2`$ with the frequency shift $`\mathrm{\Delta }=\omega _2\omega _1`$ much lower than the natural transition width. The wave intensity at frequency $`\omega _1`$ was significantly higher than at $`\omega _2`$. The four-wave generation of $`E_s`$, arose at the frequency $`\omega _s=2\omega _1\omega _2=\omega _22\mathrm{\Delta }`$ with the same polarization as in $`E_2`$. The generation power sharply rose with growing collision frequency and external $`dc`$ magnetic field. This effect can be explained in the following way.
It is convenient to expand each wave and nonlinear polarization $`P^{(NL)}(\omega _s)`$ into two circular components $`P_+^{(NL)}(\omega _s)`$ and $`P_{}^{(NL)}(\omega _s)`$ . These components contain two terms. One describes the four-wave mixing with the same polarization at two-level subsystems, while another does the radiation with contrary polarizations at three-level Zeeman subsystems (Fig. 3a).
At such a choice of polarizations, these two contributions interfere destructively and totally suppress one another, if relaxation rates for populations and quadripole moments (alignment) are equal for the Zeeman sublevels in the upper electron state. The spontaneous radiation trapping, anisotropic collisions, and/or the external magnetic field violate the amplitude balance of destructively interfering components of nonlinear polarization and induce the four-wave mixing. The magnetic field effect on the two-level configuration is compensated for the Doppler shifts. The second component of nonlinear polarization represents a double $`V`$-configuration (Fig. 3a), which is removed from resonance by the magnetic field.
Now let us consider another example (Fig. 3a). The coherence induced at the transition $`n^{}n`$, defining four-wave mixing $`\omega _s=2\omega _1\omega _2`$, is given by
$$\rho _{n^{}n}^{(2)}\alpha V_{n^{}g}\rho _{gn}^{(1)}+\rho _{n^{}g}^{(1)}V_{gn}\alpha \left[\frac{1}{\mathrm{\Omega }_2+i\mathrm{\Gamma }_{n^{}g}}\frac{1}{\mathrm{\Omega }_1+i\mathrm{\Gamma }_{ng}}\right]\frac{1}{\mathrm{\Omega }+i\mathrm{\Gamma }_{n^{}n}}=$$
$$\frac{1}{(\mathrm{\Omega }_2+i\mathrm{\Gamma }_{n^{}g})(\mathrm{\Omega }_1i\mathrm{\Gamma }_{ng})}\left[1i\frac{\mathrm{\Gamma }_{nn^{}}\mathrm{\Gamma }_{n^{}g}\mathrm{\Gamma }_{ng}}{\mathrm{\Omega }+i\mathrm{\Gamma }_{nn^{}}}\right],$$
$`(52)`$
where $`\mathrm{\Omega }_1=\omega _1\omega _{ng}`$, $`\mathrm{\Omega }_2=\omega _2\omega _{n^{}g}`$, $`\mathrm{\Omega }=\omega _2\omega _1\omega _{n^{}n}`$. A spontaneous relaxation we have $`\mathrm{\Gamma }_{ij}=(\mathrm{\Gamma }_i+\mathrm{\Gamma }_j)/2`$ and the resonance $`\mathrm{\Omega }=0`$ disappears. Collisions violate the relevant frequency equality and induce this resonance, suppressing the destructive interference .
## 8 Conclusion
The purpose of this paper is to show the diverse exhibition of nonlinear interference effects in optics and to survey certain earlier works in this field. Interference phenomena can play a governing part in numerous experiments on resonant nonlinear optics . Some such processes are sketched at Fig. 3 and explained in the figure caption.
The author is cordially grateful to V.A. Ignatchenko, with whom he began to deal with laser physics, to S.G. Rautian, who introduced him into nonlinear laser spectroscopy, and to V.P. Chebotaev and S.A. Akhmanov, who taught to him many things but are not among living, unfortunately.
I greatly acknowledge also to my colleagues Im Thek-de, Yu.I. Heller, V.V. Slabko, and V.G. Arkhipkin for their collaboration.
The author is also grateful to B. Wellegehausen, H. Welling, and G. zu Putlitz and the German Research Society; L.J.P. Hermans and J.P. Woerdman and the Netherlands Foundation for Basic Research of Matter; L. Moi from Siena University; A.Dalgamo and G. Victor from the Harvard-Smithonian Center for Astrophysics for the opportunity to work with their groups and to discuss the optical problems at seminars.
The work was partially supported by the International Science Foundation, the Russian Foundation for Basic Research, and the Krasnoyarsk Region Science Foundation.
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# A real-space, real-time method for the dielectric function
## I Introduction
Real-space methods have proven their utility in calculating the linear response of finite systems in time-dependent density functional theory. However, there has been the perception that real-space methods are unsuitable for infinite periodic systems. The problem is that the long range polarization currents are important but are dynamically independent of the local state of the electrons within the unit cell. Stated differently, the polarization gives rise to a surface charge at the surface of any finite sample, but the resulting electric field is independent of the charge density within any cell in the interior. The necessity to introduce the polarization as an independent degree of freedom has been well recognized in the literature of the density functional theory.
We will show here that in fact it is straightforward to treat infinite systems in the real-time formulation of TDDFT, simply by adding as one additional dynamic variable the surface charge. Formally, this is conveniently done adding a gauge field in within a Lagrangian. A gauge formalism has also been very recently applied by Kootskra et al , however using a frequency representation rather than solving real-time equations. In our formulation, we derive the dynamic equations from the Lagrangian:
$$L=_\mathrm{\Omega }d^3r(\frac{\underset{i}{}|\stackrel{}{}\varphi _i/ieA\widehat{z}\varphi _i|^2}{2m}+\frac{1}{8\pi }\stackrel{}{}V(r)\stackrel{}{}V(r)+en(r)V(r)+en_{ion}(r)V(r)+$$
(1)
$$𝒱_{xc}[n(r)]+𝒱_{ion}[\rho (r,r^{})])\frac{\mathrm{\Omega }}{8\pi }(\frac{dA}{dt})^2i_\mathrm{\Omega }d^3r_i\varphi _i^{}\frac{\varphi _i}{t}.$$
Here the $`\varphi _i`$ are the Bloch wave functions of the electrons, normalized so that $`n(r)=_i|\varphi _i(r)|^2`$ is the electron density. The volume of the unit cell is $`\mathrm{\Omega }`$. The electromagnetic interaction is separated into a Coulomb field $`V(\stackrel{}{r})`$ that satisfies periodic boundary conditions in the unit cell and a vector gauge field $`\widehat{z}A(t)`$. Note that the gauge field is uniform, without any dependence on $`r`$. The electric field is then given by
$$\stackrel{}{}=\stackrel{}{}V\widehat{z}\frac{dA}{dt}.$$
In these formal equations, we use units with $`\mathrm{}=c=1`$.
The other pieces of the first integral are the usual terms in the Kohn-Sham energy functional. The term $`en(r)V(r)`$ gives the direct Coulomb interaction of the electrons, except for the surface charging. The ionic interaction is separated<sup>§</sup><sup>§</sup>§The separation is somewhat arbitrary, but is useful because the periodicity of $`V`$ then takes the ionic lattice into account automatically. into a long-range part that can be associated with an ionic charge density $`n_{ion}(r)`$ and a short-range part $`𝒱_{ion}`$. The latter depends on the orbital angular momentum of the electrons in typical ab initio pseudopotentials. It therefore depends on the full one-electron density matrix $`\rho (r,r^{})=_i\varphi _i^{}(r)\varphi _i(r^{})`$. We have emphasized this point because nonlocal interactions do not respect gauge invariance. The invariance is of course restored if the density matrix is gauged. If one expresses the nonlocality in terms of the operator $``$, one makes the replacement $`ieA\widehat{z}`$. Finally, the $`𝒱_{xc}`$ is the usual exchange-correlation energy density of density functional theory.
Requiring the Lagrangian action to be stationary gives equations of motion for $`\varphi _i`$ and $`A`$ and the Poisson equation for $`V`$. The dynamic equation for $`\varphi _i`$ is the time-dependent Kohn-Sham equation,
$$\frac{^2}{2m}\varphi _i\frac{e}{mi}A_z\varphi _i+\frac{e^2}{2m}A^2\varphi _i+(eV+\frac{\delta 𝒱_{ion}}{\delta n}+\frac{\delta 𝒱_{xc}}{\delta n})\varphi _i=i\frac{}{t}\varphi _i$$
(2)
The equation for $`A`$ is
$$\frac{\mathrm{\Omega }}{4\pi }\frac{d^2A}{dt^2}\frac{e}{m}\underset{i}{}\varphi _i|_z/i|\varphi _i+\frac{e^2}{m}AN_e+\frac{\delta }{\delta A}_\mathrm{\Omega }𝒱_{ion}d^3r=0$$
(3)
where $`N_e=_\mathrm{\Omega }d^3rn(r)`$ is the number of electrons per unit cell.
## II Linear Response, Sum Rule and Simple Models
The calculation of the dielectric function using the above real-time dynamic equations is very similar to the corresponding calculation of dynamic polarizability of finite systems. We first solve the static equations (with $`A`$=0) to get the ground state electron orbitals and the periodic Coulomb potential $`V`$. The system is then perturbed by making a sudden change in $`A`$, $`A(t=0_+)=A_0`$. This corresponds to applying a short-duration electric field at $`t=0`$, $`(t)=A_0\delta (t)`$. The dynamic equations are then applied to evolve the variables in time, finding the time evolution of the polarization electric field $`(t)=dA(t)/dt`$. The dielectric function $`ϵ(\omega )`$ is just the ratio of the Fourier components of the external and the total fields; it is given by
$$\frac{1}{ϵ(\omega )}1=\frac{1}{A_0}_{0_+}^{\mathrm{}}e^{i\omega t\eta t}\frac{dA(t)}{dt}𝑑t.$$
(4)
Here $`\eta `$ is a small quantity to establish the imaginary part of the response. In principle the resulting theory automatically respects the Kramers-Kronig relation.
The linear energy-weighted sum rule is easily derived in this formalism. The sum rule may be expressed as
$$_0^{\mathrm{}}\omega \mathrm{Im}ϵ^1(\omega )𝑑\omega =\frac{2\pi ^2e^2N_e}{m\mathrm{\Omega }}.$$
(5)
To calculate the sum rule with our Lagrangian, we write the integral using eq. (4)
$$_0^{\mathrm{}}\omega \mathrm{Im}ϵ^1(\omega )𝑑\omega =\frac{1}{A_0}_0^{\mathrm{}}𝑑t\frac{dA}{dt}\mathrm{Im}_0^{\mathrm{}}𝑑\omega \omega e^{i\omega t\eta t}=\frac{\pi }{2A_0}\left(\frac{d^2A}{dt^2}\right)_{t=0_+}.$$
(6)
The second derivative in the last expression can be easily found from eq. (3). At $`t=0_+`$, the wave functions have not yet had time to change, $`A(0_+)=A_0`$ and $`_z=0`$. Then if the last term in eq. (3) can be neglected,
$$\frac{d^2A}{dt^2}=\frac{4\pi e^2N_eA_0}{m\mathrm{\Omega }}$$
(7)
and eq. (5) follows immediately. Thus the time-dependent treatment satisfies the sum rules automatically to the extent permitted by the last term. That term is only nonzero for nonlocal pseudopotentials, and in fact it may improve the accuracy of the theory by incorporating effects of the core electrons on the dynamic properties.
Let us now see how the gauge field treatment works in a simple analytically solvable model, namely the electron gas. As mentioned before, when the field $`A_0`$ is applied, there is no immediate response to the operator $``$, since the wave function does not change instantaneously. However, in the Fermi gas, the single-particle states are eigenstates of momentum so the response remains $`=0`$ for all time. Putting this in eq. (3), and dropping the pseudopotential term, the equation for $`A`$ becomes simple harmonic motion, with the solution
$$A(t)=A_0\mathrm{cos}\omega _{pl}t$$
(8)
where $`\omega _{pl}`$ is the plasmon frequency,
$$\omega _{pl}^2=\frac{4\pi e^2N_e}{m\mathrm{\Omega }}=\frac{4\pi e^2n}{m}.$$
(9)
The dielectric function may now be calculated from the time integral eq. (4). One obtains the familiar electron gas result,
$$ϵ(\omega )=1\frac{\omega _{pl}^2}{\omega ^2}.$$
(10)
One sees that the derivation here is much simpler than the usual one using the Coulomb gauge. There one formulates the response in a particle-hole representation, and takes the external field to be of the form $`e^{iqr}`$ with $`q`$ finite. The dielectric function is then found by taking the $`q0`$ limit.
We can make another simple model for the opposite extreme of a tightly bound electron in the unit cell. Assume that the ion potential $`eV_{ion}(r)+\delta 𝒱_{ion}/\delta n`$ can be approximated by a harmonic oscillator potential in the region over which the electron wave function is appreciable. According to Kohn’s theorem, the response is just the same as for an isolated electron in the same ionic potential. This comes out of eq. (2-3) in the following way. The initial impulse $`A_0`$ starts the electron moving, and as a result both $`V(r)`$ and $`\delta 𝒱_{xc}/\delta n(r)`$ become time-dependent. Together with the changing $`A`$, the electron in the unit cell drags its self-induced field with it, and the accelerations associated with these three terms in eq. (2) exactly cancel. The remaining ionic terms then produce simple harmonic motion for $`\varphi `$.
## III Numerical details
The computational algorithm we employ is identical to the ones we used for clusters and molecules, which was based on a method introduced in nuclear physics. The Kohn-Sham operator is represented on a real space grid as in ref. . There are a number of technical details associated with the periodicity and with the gauge potential that did not arise for the finite-system calculations. In the new code, the potential $`V(r)`$ is calculated by Fourier transformation of the Poisson equation rather than a relaxation method. This gives automatically the required periodicity to $`V(r)`$. The wave functions $`\varphi _i`$ represent Bloch states of the periodic lattice, and they are constructed with the corresponding periodic boundary conditions labeled by the Bloch momenta $`k`$. The periodic boundary condition on the Bloch wave function $`\varphi _k(r+a)=\mathrm{exp}(iak)\varphi _k(r)`$ is easily implemented in the relaxation method used to find eigenstates. In practice, many Bloch states are needed to obtain smooth dielectric functions. However, constructing the states takes much less time than for the same number of electrons in a finite system, because the Bloch states in a given band are automatically orthogonal.
We use here the same energy density functional that we used previously for finite systems. Only the valence electrons are included explicitly; core electrons are treated by a pseudopotential. The exchange-correlation energy of the electrons is calculated in the local density approximation following the prescription of ref. .
The presence of a vector gauge potential requires a modification in the pseudopotential calculation, as indicated in the introduction. In particular, the $`A`$-dependence of the $`𝒱_{ion}`$ term in eq. (2) must be consistent with the last term in eq. (3) in order to have the algorithm conserve energy. We implement the $`A`$-dependence of $`𝒱_{ion}`$ simply by gauging the density matrix directly,
$$𝒱_{ion,A}(\rho (r,r^{}))=𝒱_{ion}(e^{iA(zz^{})}\rho (r,r^{}))$$
As in the finite systems calculations, energy is conserved to a very high accuracy with the algorithm , provided the time step is less than the inverse energy span of the Kohn-Sham operator.
## IV Lithium
In this section we demonstrate the feasibility of the method with lithium as an example of a simple metal. As other alkali metals, lithium has a Fermi surface which is nearly spherical. However, unlike sodium and potassium, the effective mass of the electrons at the Fermi surface is significantly enhanced over the free-Fermi gas value ($`m^{}1.6m_e`$).
The Kohn-Sham operator is represented in coordinate space with a uniform spatial mesh. The lattice spacing of the bcc unit cell of Li metal is $`a=3.49`$ Å, and we take a mesh spacing of $`\mathrm{\Delta }x=0.58`$ to subdivide the cell into a $`6^3`$ lattice of mesh points. We use a time step of $`\mathrm{\Delta }t=0.01`$ eV<sup>-1</sup> which is sufficient to conserve energy to 10<sup>-4</sup> eV over the time integration interval, $`T=18`$ eV<sup>-1</sup>.
We sample the occupied states with a uniform mesh in momentum space. With a finite set of occupied orbitals, the allowed excitation energies will be discrete, and the metallic behavior, $`ϵ(\omega )\mathrm{}`$, is only reached in the limit of an infinitely dense momentum space lattice. However, it is our view that the TDLDA loses validity at long times when other degrees of freedom can be excited. This is the case for the low frequency response of metals, where the imaginary part of the dielectric responses is dominated by phonons and inelastic electron scattering.
In Fig. 1 we show the normalized time-dependent polarization field, $`A_0^1dA/dt`$ over the time interval $`t=[0,18]`$ eV<sup>-1</sup>. The inset shows an expanded view of the initial response in the interval eV<sup>-1</sup>. The dashed line is the comparison with the linear behavior deduced from the sum rule, $`A_0^1dA/dt=4\pi e^2N_et/m\mathrm{\Omega }`$. The agreement shows that the local sum rule in nearly satisfied, despite the fairly large optical effective mass.
In Fig. 2 we show the inverse dielectric function computed from eq. (4) for various meshes in the Brillouin zone. We employ $`\eta =0.2`$ eV to smooth the response in the Fourier transformation. We see that the response becomes smoother, the more finely the Fermi sea is sampled. With a $`32^3`$ lattice of Bloch states, we get results smooth enough to be compared with measurement.
In Fig. 3 we show the real and imaginary parts of the inverse dielectric function in the frequency interval $`020`$ eV/$`\mathrm{}`$. The dashed lines show the empirical function from ref. . There is also another theoretical calculation in the literature, . The agreement is quite good, especially considering that the calculation is ab initio with a energy density functional that much simpler than more recent ones. The main feature in the absorptive response is the plasmon at 7 eV and its width. The peak position is significantly downshift from the naive plasmon frequency, $`\omega =\sqrt{4\pi e^2n/m}8`$ eV. The width is associated with interband transitions and is also well reproduced.
## V Diamond
In this section we compute the dielectric response of a typical elemental insulator, diamond. The diamond lattice is represented in our calculation by the primitive unit cell which contains 8 carbon atoms. The four valence electrons of each carbon are calculated explicitly while the core electrons are only treated implicitly by the pseudopotential. We found in earlier studies of carbon structures that the Kohn-Sham Hamiltonian requires a mesh spacing $`\mathrm{\Delta }x=0.3`$ Å to get orbital energies to an accuracy of 0.1 eV; in the calculation here we take a $`12^3`$ lattice in the unit cell which implies $`\mathrm{\Delta }x=3.56/12\mathrm{\AA }=0.297`$ Å. With a smaller mesh spacing than for lithium, the span the Kohn-Sham operator is increased and the time step $`\mathrm{\Delta }t`$ must be reduced accordingly. We use here $`\mathrm{\Delta }t=0.002`$ eV<sup>-1</sup>.
With a cubical unit cell and 8 carbon atoms, there are $`84/2=16`$ occupied bandsThe bands are actually two-fold degenerate because we have not exploited the symmetry that allows a smaller unit cell with 4 carbons.. In each band we take a lattice of up to $`16^3`$ points to represent the Bloch states.
For an insulator, a reference point of the vector potential $`A(t)`$ should be irrelevant. Since all the points within Brillouin zone are occupied, there is a static solution of eq.(2-3) for each constant vector potential. However, in our real-space, real-time calculation, we found this is violated due to the discrete representation in coordinate and momentum spaces. The total energy does depend on $`A`$, and there appears a spurious low frequency mode in the vector potential $`A(t)`$ associated with the adiabatic evolution of the electronic wave function. Below that frequency, a spurious conduction feature appears in the response.
These features are shown in the plots of the response in Fig. 4. We see that the spurious adiabatic evolution gives an unphysical plasmon at $`1.2`$ eV, which dominates the dielectric response at lower frequencies. Though the amount of strength associated with this spurious plasmon is very small, 0.007 electrons out of the total of 32, it has a qualitative effect on the response at very low frequency. The frequency and strength decrease the finer the spatial mesh, showing that it is an artifact of the discrete mesh representation of the coordinate space.
To infer the dielectric function near $`\omega =0`$, we apply the Kramers-Kronig relation to the imaginary part of the response, but excluding the spurious plasmon peak. This gives the predicted dielectric function shown in Fig. 5. The empirical dielectric function is shown as the dashed line. The agreement is good, as indeed was found solving the TDLDA equations by other methods, but one can also see the effect of the well-known shortcoming of TDLDA, that the predicted band gaps are too small. The theoretical absorption strength become significant starting at about 5 eV excitation, while the empirical absorption begins at around 7 eV. Nevertheless, the dielectric constant comes out in good agreement with the empirical, being within a percent of the empirical value of $`ϵ(0)=5.67`$.
## VI Conclusions
We see that the method not only works in principle, but produces fairly accurate dielectric functions in the cases of a simple metal and a simple insulator. In lithium, the theory describes the metallicity as well as the interband transitions. In diamond, there is a spurious plasmon at low frequency due to the discrete mesh representation in coordinate and momentum spaces. However, it can be easily dealt with and then the dielectric function has an excellent quality except for a small band gap region. We find that two benefits of the real-space, real-time formulation of the TDLDA in finite systems are preserved in our implementation here. The real-space method allows the Kohn-Sham operator and electron-electron interactions to be evaluated efficiently . Computational efficiency is also gained by calculating the response in real time in that all frequencies are calculated at once. Finally, the method requires much less storage than methods using a particle-hole representation of the time-varying wave function.
## VII Acknowledgment
We acknowledge discussions with E.K.U. Gross, R. Martin, J. Rehr, and R. Resta. This work was supported in part by the Department of Energy under Grant DE-FG03-00-ER41132, by the DGES (PB98-0345) and by the JCyL (VA28/99), and by the Ministry of Education, Science and Cultures (Japan), No. 11640372. JI and KY acknowledge the Institute of Solid State Physics, University of Tokyo, and the Research Center for Nuclear Physics, Osaka University for their use of supercomputers. AR acknowledges the hospitalty of the Institute for Nuclear Theory where this work was started and computer time provided by the Centre de Coputació i Comunicacions de Catalunya.
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# An exploratory study of the hard X-ray variability properties of PG quasars with RXTE
## Abstract
We have monitored with the RXTE PCA the variability pattern of the 2–20 keV flux in four PG quasars (QSOs) from the Laor et al. (1994) sample. Six observations of each target at regular intervals of 1 day were performed. The sample comprises objects with extreme values of Balmer line width (and hence soft X-ray steepness) and spans about one order of magnitude in luminosity. The most robust result is that the variability amplitude decreases as energy increases. Several options for a possible ultimate driver of the soft and hard X-ray variability, such as the influx rate of Comptonizing relativistic particles, instabilities in the accretion flow or the number of X-ray “active sites”, are consistent with our results.
The study that is the subject of this paper aims at a systematic investigation of the variability properties of a sizable sample of optically selected QSOs (the first after the seminal paper by Zamorani et al 1984). The sample is the Laor et al. (1997) sample of 23 PG quasars, selected to have $`\mathrm{z}<0.4`$, Galactic $`\mathrm{N}_\mathrm{H}<1.9\times 10^{20}`$ cm<sup>-2</sup> and $`\mathrm{M}_\mathrm{B}<23`$. The timing results of a monitoring campaign of six of them with the ROSAT/HRI are discussed by Fiore et al. 1998 (F98). They discovered large amplitude (factor of 2) and rapid (timescale $``$1 day) variability. “Steep” QSOs (i.e., those whose ROSAT/PSPC energy spectral index $`\alpha _{\mathrm{sx}}`$ is $`\stackrel{>}{}3`$) show systematically larger amplitude variation than “flat” ones. Laor et al. (1994) discovered a strong correlation between $`\alpha _{\mathrm{sx}}`$ and the Full Width Half Maximum (FWHM) of the H$`\beta `$ optical lines. This correlation can be explained if the size of the Broad-Line Region (BLR) is uniquely determined by the luminosity of the active nucleus and the BLR gas is virialized. In this case, $`\mathrm{L}/\mathrm{L}_{\mathrm{Edd}}`$ scales inversely as the square of the bulk velocity. In this framework, objects with steep $`\alpha _{\mathrm{sx}}`$ display larger variability, implying a lower mass black hole and therefore a higher $`\mathrm{L}/\mathrm{L}_{\mathrm{Edd}}`$.
An exploratory program to study the hard X-ray variability properties of a sub-sample of PG QSOs on timescales $`\stackrel{>}{}`$1 day was started with the Rossi X-ray Timing Explorer (RXTE) Observatory. The objects are listed in Tab. 1.
The monitoring campaign consisted of six pointings for each target, with intervals of approximately one day between each pointing.
The 1-day averaged light curves for all sources of the sample are shown in Fig. 1. All light curves in Fig. 1
exhibit some degree of variability, and it is generally more marked in the soft than in the hard energy band. To characterize quantitatively the variability in our sample, we used the so-called average structure function (SF, Di Clemente et al. 1996). In the formulation suggested by F98, it consists of the mean of the logarithmic ratio between each pair of flux measurements:
$$SF=<|2.5\times \mathrm{log}[f(t_j)]/[f(t_i)]|>$$
The errors on the SF are given by the standard uncertainties on the average. For our light curves, sampled six times at regular intervals of approximately 1 day between consecutive observations, we can build a five-point SF. We have calculated the SF in the soft and hard energy bands for all the QSOs in our sample (the thresholds between the bands are defined in Tab. 1). The “soft” SF is systematically higher than the “hard” one on all timescales (see the left panel of Fig. 2), although the difference is
significant at more than the 1-$`\sigma `$ level only for data points corresponding to $`\mathrm{\Delta }\mathrm{t}\mathrm{t}_\mathrm{j}\mathrm{t}_\mathrm{i}3`$ days. The results are far less clear if objects with “narrow” and “broad” optical lines are compared (see the right panel in Fig. 2). However, the suggestion exists that narrow-line QSOs remain more variable in the hard X-ray band (see Tab. 2).
The cause of the lower variability amplitude in harder X-rays can be multifold: a) Compton-reprocessing of the primary nuclear continuum (for which, however, there is no convincing evidence yet; George et al. 2000); b) smearing induced by Comptonization on a variability pattern, driven by instabilities in the accretion disk; c) smearing induced by the larger number of interactions that higher energy photons undergo in the Comptonizing plasma; d) the number of “active sites”, whereby the soft X-rays might be produced by the superposition of a large number of “flares” (originating, e.g. in an accretion disk corona), while the harder photons are produced in a more homogeneous and probably compact innermost region.
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# Considerations about universality in phase-ordering of binary liquids
## 0.1 Introduction
Phase-ordering is observed in many spatial systems, including ones that underwent spinodal decomposition or mechanical mixing. Examples include a wide variety of systems including magnetic systems, binary alloys, binary fluids, but also some reaction diffusion systems and some models of gene evolution.
In a seminal paper Hohenberg and Halperin showed that most of these systems can be categorized to belong to a small number of categories, and each member of a category shows the same universal phase-ordering behavior. In particular they distinguished between systems in which the order parameter is not conserved,e.g. magnetic systems or genetic systems, and systems in which the order parameter is conserved, e.g. binary alloys or binary fluids. Liquid systems, such as binary fluids and magnetic fluids, can phase-order not only by diffusion but also via flow and therefore they are in a universality class of their own.
The phase-ordering process has been conjectured to observe scaling, i.e. the systems are characterized by a single length-scale $`L(t)`$ such that the morphologies of the system scaled but that length-scale are statistically self-similar. The time dependence of this length-scale is then known as the scaling law. The scaling laws can be deduced by dimensional analysis from the equations of motion. For systems without hydrodynamics $`L(t)`$ scales as $`t^{1/2}`$ for systems with non-conserved order parameter and as $`t^{1/3}`$ for systems with a conserved order parameter.
Because binary fluids have more that one mechanism for domain coarsening there is more that one growth law. Off-critical quenches that contain a much smaller volume fraction of one component form droplets and this morphology grows mostly via diffusion leading to a $`L(t)t^{1/3}`$ growth law. Most studies, however, deal with critical quenches with equal volume fractions for both components. Because of this symmetry bi-continuous morphologies are typically formed and both diffusion and hydrodynamic growth mechanisms operate.
There are, however, exceptions to this general rule. Sometimes dynamical asymmetries in the diffusion constant of the two phases or their viscoelastic properties can lead to droplet morphologies. A stationary morphology that consists of circular droplets has no hydrodynamic pathway of coarsening. Once two droplets coalesce they can induce a flow that can induce further coalescence. The central idea of this paper is that one of the main hydrodynamics pathways of coarsening can be suppressed if the morphology consists of droplets leading to a new scaling state. The possibility of the non-uniqueness of the scaling state was first suggested by A. Rutenberg, although at the time no persuasive numerical evidence could be found.
The droplet morphologies that we require are not unusual and can often seen in the early stages of viscoelastic spinodal decomposition as well as in systems with an order-parameter dependent mobility. We will use a viscoelastic lattice Boltzmann method to create the initial droplet morphology and then use a symmetric Newtonian lattice Boltzmann method to evolve the initial droplet morphology to verify that this new scaling state exists.
In viscoelastic phase-separation the droplet morphology consists of low-viscosity droplets suspended in the viscoelastic matrix. The opposite morphology is observed in mechanical mixing where the low viscosity material will form the matrix phase. We show that both droplet morphologies are stable under phase-ordering, in agreement with experimental results.
## 0.2 Numerical method
For the simulations we use the viscoelastic two-component lattice Boltzmann simulation introduced in an earlier paper. Briefly, in lattice Boltzmann simulations densities $`f_i`$ that are associated with velocities $`v_i`$ are streamed on a lattice according to the lattice Boltzmann equation
$`f_i(𝐱+𝐯_i\mathrm{\Delta }t,t+\mathrm{\Delta }t)=`$
$`f(𝐱,t)+\mathrm{\Delta }t{\displaystyle \underset{j}{}}\mathrm{\Lambda }_{ij}\left[f_j^0(𝐱,t)f_j(𝐱,t)\right]`$ (1)
where $`f_i^0`$ is an equilibrium distribution, $`\mathrm{\Lambda }_{ij}`$ is a collision matrix, and $`𝐯_i\mathrm{\Delta }t`$ is a lattice vector. The velocity set for our simulation consists of 17 velocities given by $`\{(0,0)`$, $`(1,0)`$, $`(0,1)`$, $`(1,0)`$, $`(0,1)`$, $`(1,1)`$, $`(1,1)`$, $`(1,1)`$, $`(1,1)`$, $`(1,0)`$, $`(0,1)`$, $`(1,0)`$, $`(0,1)`$, $`(1,1)`$, $`(1,1)`$, $`(1,1)`$, $`(1,1)\}`$. Note that the last 8 velocities are the same as the previous eight velocities. This duplicity allows the simulation to have two independent stresses which represent a viscoelastic and a purely viscous contribution to the total stress tensor. The two contributions are used to produce a Jeffrey’s model for the stress (see eqn. (8)). The algorithm is required to conserve mass and momentum, but not energy. Energy conservation is replaced by a condition of constant temperature. The macroscopic density $`\rho `$ and velocity $`𝐮`$ are defined as
$$\rho =\underset{i}{}f_i\rho 𝐮=\underset{i}{}f_i𝐯_i.$$
(2)
To simulate a two-component mixture we now have to consider the densities of the two fluids $`\rho _1`$ and $`\rho _2`$. The first lattice Boltzmann equation (1) is now an equation for the total density $`\rho =\rho _1+\rho _2`$ and we introduce a second lattice Boltzmann equation for $`g_i`$ to describe the evolution of the density difference $`\varphi =\rho _1\rho _2`$:
$`g_i(𝐱+𝐯_i\mathrm{\Delta }t,t+\mathrm{\Delta }t)=`$
$`g(𝐱,t)+{\displaystyle \frac{\mathrm{\Delta }t}{\tau }}\left[g_i^0(𝐱,t)g_i(𝐱,t)\right]`$ (3)
where we choose a diagonal collision matrix with a single relaxation time $`\tau `$. These densities are only defined on the first nine velocities $`𝐯_i`$. The density difference $`\varphi `$ is given by
$$\varphi =\underset{i}{}g_i.$$
(4)
By choosing appropriate equilibrium distributions and an appropriate collision matrix, we ensure that the following partial differential equations are being simulated up to second order in the derivatives but assuming that the relaxation of the viscoelastic stress $`\sigma `$ is slow ($`\theta 1/\sqrt{ϵ}`$):
$`_t\rho +_𝐱(\rho 𝐮)`$ $`=`$ $`0`$ (5)
$`\rho _t𝐮+\rho 𝐮.𝐮`$ $`=`$ $`_𝐱P+_𝐱(\sigma _v+\sigma )`$ (6)
$`\sigma _v`$ $`=`$ $`\nu _{\mathrm{}}((\rho 𝐮)+((\rho 𝐮))^T.𝐮\delta )`$ (7)
$`+\xi _{\mathrm{}}.𝐮\delta `$
$`\sigma +\theta (\varphi )\sigma _{(1)}`$ $`=`$ $`(\nu _0(\varphi )\nu _{\mathrm{}})((\rho 𝐮)+((\rho 𝐮))^T)`$ (8)
$`_t\varphi +_𝐱(\varphi 𝐮)`$ $`=`$ $`D^2\mu +.((\varphi /\rho ).(P\sigma ))`$ (9)
where $`\delta `$ is the identity matrix, $`\sigma `$ is the viscoelastic stress tensor, $`\sigma _{(1)}=_t\sigma +𝐮.\sigma \sigma .(𝐮)(𝐮)^T\sigma `$ is its upper convected derivative, $`P=0.5\rho +0.007(\varphi \varphi 0.5\varphi .\varphi \delta \varphi ^2\varphi \delta `$ is the pressure tensor, and $`\mu =0.55\varphi /\rho +0.25\mathrm{ln}((\rho +\varphi )/(\rho \varphi ))0,007^2\varphi `$ is the chemical potential. The parameters $`\nu _{\mathrm{}}`$, $`\xi _{\mathrm{}}`$, and $`\theta `$ are determined by the eigenvalues of the collision matrix. The values for the parameters were $`\mathrm{\Delta }t=1`$, $`\tau =1`$, $`D=0.5`$, $`\xi _{\mathrm{}}=0.31`$, and $`\nu _{\mathrm{}}=0.01`$. For the low-viscosity phase we used $`\theta =0.055`$, $`\nu _0=0.013`$ and for the viscoelastic phase $`\theta =39.5`$ and $`\nu _0=1.97`$. For the symmetric simulations of Figure 1 we used $`\theta =0.055`$ and $`\nu _0\nu _{\mathrm{}}=0.075`$.
## 0.3 Simulations
We performed simulations of critical spinodal decomposition of a viscoelastic binary mixture in two dimensions where one component is much more viscoelastic than the other. These simulations lead to the usual morphologies in which the viscoelastic phase is connected and the less viscoelastic phase is dispersed. We performed our simulations on a $`256^2`$ and a $`1024^2`$ lattice and after about 1000 iterations the less viscoelastic phase is completely dispersed, although the domains are still highly deformed. We used this morphology as an initial condition for a simulation where we make both components purely viscous to examine the effects of initial conditions that are not symmetric on symmetric binary fluid mixtures. We choose the viscosity such that a system started with a random initial condition will be in the viscous scaling state. This allows us to distinguish the effect that the morphology created by viscoelastic phase-separation has from the effect of viscoelasticity itself in the late-time phase-ordering process.
In Figure 1 we see a comparison of the phase-ordering from a droplet morphology (a) and a symmetric initial condition (b). The morphologies for both systems are shown after 1000, 2000, 4000, and 8000 iterations. We see that the phase-ordering of the morphology with a dispersed phase leads to an even more pronounced dispersed phase where almost all domains are circular at late times. We see that droplet coalescence occurs frequently with only few domains vanishing due to the evaporation-condensation mechanism underlying Oswald ripening. The droplet coalescence, however, is not frequent enough to change the connectivity of the domains. Instead, we observe that domains become more circular on average, suggesting that even for very long times we do not expect a transition to a bi-continuous morphology. (This is no longer true when we leave viscous regime and enter into the inertial regime. Here a return to the inertial scaling state is observed). In Figure 1(c) we see that the growth law for the droplet phase appears to be $`L(t)t^{1/2}`$ and is smaller than the $`L(t)t^{2/3}`$ seen for the symmetric phase-ordering shown in Figure 1(b). More recent simulations of the droplet morphology, however, establish scaling over several orders of magnitude and find that the actual growth-law is $`Lt`$ . This emphasizes the fact that dynamic scaling analysis that does not cover several decades can be misleading. Unfortunately such studies are very computationally demanding and exist only in few studies.
The existence of this second scaling state, distinct from the bi-continuous scaling state, is important to understanding the late-time regime of viscoelastic phase-separation because, even in the absence of viscoelastic effects, we observe a droplet-morphology evolving from the initial morphology created by viscoelastic phase-separation. This result is also important for practical reasons in processes where late-time morphologies need to be controlled. It is well known that mixing of high-viscosity and low-viscosity components by means of mechanical agitation leads to morphologies where the high-viscosity phase is dispersed in droplets for volume fractions of the low-viscosity component of much less than 50%. This effect is enhanced if the high-viscous phase is viscoelastic. This leads us to consider what the late-time morphology of a phase-ordering system with an early morphology created by mechanical mixing would be.
In order to answer this question we performed a simulation of viscoelastic phase-separation and after a droplet morphology had been formed at 1000 time-steps, we inverted the properties of the two components. We used this state as a model for the morphologies of dispersed droplets of the viscoelastic phase in a matrix of the low viscosity Newtonian phase that is typical of mechanically mixed morphologies. We then continued the simulations and observed the phase-ordering behavior of the new morphology. The results of these simulations are shown in Figure 2. In Figure 2(a) the morphology for a viscoelastic phase-separation of a 50%-50% mixture is shown. We should emphasize that the phase-ordering (see eqn. (9)) does not have a composition-dependent diffusion constant and therefore no domain shrinkage is observed in these simulations. Domain shrinkage can be a very slow process in viscoelastic phase-separation that prolongs the spinodal decomposition process and makes it more difficult to differentiate the early-stage decomposition and the late-time phase-ordering processes. This does not reduce the validity of our results, however, since we are only interested in the late-time behavior when the domain shrinkage is completed.
The morphologies shown in Figure 2(b) are of a simulation where at 1000 iterations, after the dispersion was achieved, the viscoelastic and the low viscosity components were exchanged. We see that the morphology remains a droplet morphology, albeit now of the viscoelastic phase. Comparing Figures 2 (a) and (b) we see that the most important factor in determining the final morphology is the early stage dynamics, and that phase-ordering scaling morphologies of both dispersed viscoelastic and dispersed low viscosity domains exist.
From Figure 2(c) we see that each scaling state appears to have a different growth law. The morphology of dispersed viscoelastic domains grows as $`Lt^{0.4}`$ whereas the morphology of dispersed low viscosity domains grows as $`t^{0.2}`$ after near circular droplets have been formed. An anomalous slow growth in viscoelastic phase-separation has first been observed experimentally by Tanaka who found $`Lt^{0.15}`$ for a system of high molecular-weight polystyrene/ diethyl malonate (4.0 wt.%). It must, however, be emphasized that it is impossible to conclusively determine scaling exponents from such a short run. It is conceivable that a more extensive scaling analysis could lead to different exponents.
These simulations also emphasize the difference between a morphology after spinodal decomposition and after mechanical mixing. The morphology after spinodal decomposition is a dispersed low viscosity phase, whereas the state after mechanical mixing has a dispersed viscoelastic phase. Subsequent phase-ordering does not change the connectivity of these states, in agreement with the conventional wisdom that there is a profound difference in states produced by spinodal decomposition and mechanical mixing.
## 0.4 Conclusions
In this article we have shown that more than one scaling state exists for late-time spinodal decomposition of two-dimensional binary fluids and that local correlations in the initial conditions or the early time behavior of a phase-separating binary mixture can be very important in selecting one of these scaling states. We have also explained the difference between viscoelastic phase-ordering states after spinodal decomposition and mechanical mixing. Our results show that the volume fraction and the physical properties of a mixture do not select a morphology by themselves, but that the morphology of the initial state is of paramount importance. This is why viscoelastic phase-separation can lead to unusual late-time scaling states even when viscoelasticity is no longer important at large length scales.
## Acknowledgments
The authors acknowledge the financial support of DuPont Chemical Company. A.W. acknowledges the support of EPSRC Grant GR/M56234 and would like to thank Craig Carter for the generous permission to use his Origin2000 computer.
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# The effects of interactions and disorder in the two-dimensional chiral metal
## I Introduction
An unusual type of two-dimensional conductor, the chiral metal, may be formed at the surface of a layered, three-dimensional system in the presence of a magnetic field which has a component perpendicular to the layers . It occurs in a system of weakly-coupled layers at magnetic field-strengths for which an isolated layer would have an integer quantised Hall conductance. Under these circumstances, bulk states are localised and the extended, chiral edge states associated with each layer hybridise to form a two-dimensional conductor at the surface of the sample. Properties of this chiral metal may be probed by conductance measurements with current flow normal to the layers.
Early experimental work on semiconductor multilayer samples established that the integer quantum Hall effect of a single layer is indeed robust against weak interlayer tunneling, and suggested that, within a quantum Hall plateau, interlayer conductivity in the bulk vanishes in the low-temperature limit. Subsequent theoretical discussions drew attention to the unusual nature of surface states in multilayer quantum Hall systems, and some aspects of the theory of the chiral metal have since been explored in considerable detail . Recent experiments by Druist et al have isolated the contribution of surface states to vertical transport and demonstrated that mesoscopic conductance fluctuations offer a way to study inelastic scattering in the chiral metal . Conduction by surface states has also been invoked in the interpretation of further experiments on semiconductor multilayer samples , and on the bulk quantum Hall effect in both organic conductors and an inorganic quasi-two dimensional conductor .
In this paper we extend the theoretical treatment of mesoscopic chiral conductors in two respects. First, we develop some aspects of the theory of conductance fluctuations which we hope will be useful for the interpretation of experiments. Second, we study electron-electron interactions in the dirty chiral metal.
The existing understanding of conductance fluctuations in the chiral metal is based on calculations for the variance of sample-to-sample conductance fluctuations in systems that are fully phase-coherent. From the results of those calculations, estimates have been made for the variance of conductance fluctuations measured as a function of magnetic field in samples that are only partially phase-coherent, and also for the range in magnetic field of conductance correlations. These estimates determine the relevant scales but do not yield numerical coefficients or functional forms. Moreover, they depend on a continuum treatment of the system of coupled edge states, and apply if the inelastic scattering rate, $`\tau _{in}^1`$ is smaller than the inter-edge tunneling rate, $`\tau _{}^1`$, while existing experiments appear to be in the opposite regime . We supplement this past work by calculating in full the conductance correlation function $`\delta g(B)\delta g(B+\delta B)`$, both for $`\tau _{in}^1\tau _{}^1`$ and for $`\tau _{in}^1\tau _{}^1`$.
Separately, we consider conductance as a function of Fermi energy. For conventional disordered conductors, changes in Fermi energy provide a way to sample the conductance distribution. In contrast, we find that the conductance of a partially-coherent chiral metal sample, though dependent on the disorder realisation, is unchanged by small variations in Fermi energy. One consequence is that thermal smearing of the electron distribution function does not decrease the amplitude of conductance fluctuations in the chiral metal.
Turning to electron-electron interactions, we establish the form of the dynamically screened interaction and examine some of its consequences. Most importantly, we calculate the inelastic scattering rate due to electron-electron interactions, obtaining $`\tau _{in}^1T^{3/2}`$ for $`\tau _{in}\tau _{}`$. We also examine other interaction effects which are known to be of interest in conventional disordered conductors. Interactions in diffusive metals are responsible for a zero-bias anomaly in the tunneling density of states , but – as noted previously , and as we confirm – this is absent from the chiral metal. Similarly, the leading interaction contribution to the mean conductance has a singular temperature dependence in conventional disordered conductors , but we find that it cancels in the chiral metal.
The remaining sections of this paper are organised as follows. We review briefly past work and collect our most important new results in Sec. II. We describe calculations on conductance fluctuations in Sec. III, and treat electron-electron interactions in Sec. IV. We summarise and discuss recent experiments in Sec. V.
## II Background and main results
It is useful first to introduce some notation. Consider a layered sample which exhibits the bulk quantum Hall effect, with one Landau level below the Fermi energy. A chiral metal is formed at the sample surface, from hybridisation of one edge state per layer. We denote the interlayer spacing by $`a`$, the chiral velocity by $`v`$ and the interlayer diffusion constant by $`D`$. The electron density of states per unit area is $`n=1/hva`$ and the surface conductivity in the transverse direction, as given by the Einstein relation, is $`\sigma =D/va`$, in units of $`e^2/h`$, which we use throughout this paper.
Phase coherent mesoscopic conductors are known to have three distinct mesoscopic regimes, according to relative sample dimensions . For definiteness, suppose that the sample is a cylinder of height $`L`$ and circumference $`C`$, with its axis perpendicular to the layers. Take the number of layers to be $`N`$, so that $`L=Na`$. Two intrinsic length scales arise, with which $`L`$ should be compared. The first of these involves the mean conductance of the sample, which if Ohm’s law applies is $`g=C\sigma /L`$: the quasi one-dimensional localisation length, $`\xi =2C\sigma `$, is the value of $`L`$ at which $`g`$ is of order the quantum unit of conductance. Samples with $`L\xi `$ have states localised in the vertical direction and are quasi one-dimensional insulators. Samples with $`L\xi `$ are metallic and have Ohmic dependence of $`g`$ on sample dimensions. A second length scale influences the amplitude of conductance fluctuations. It is the distance $`L_1`$ that an electron diffuses in the transverse direction during the time taken to circumnavigate the sample: $`L_1=(DC/v)^{1/2}(a\sigma C)^{1/2}`$. Samples with $`L_1L\xi `$ are quasi one-dimensional metals without time-reversal symmetry and, in common with other examples of such systems, have a variance for the conductance of $`(\delta g)^2=1/15`$. Samples with $`LL_1`$ are two-dimensional chiral metals. If contacts are attached to these samples for conductance measurements, electrons will escape by diffusion to the contacts before completing a circuit of the sample in the chiral direction. As a result, multiple scattering of an electron from a given impurity is completely supressed. Such samples present a number of essentially independent parallel strips for conduction in the transverse direction, each having a width in the chiral direction given by the distance an electron propagates before reaching a contact. This width is $`L^2v/DC(L/L_1)^2`$. The conductance of the two-dimensional chiral metal may be thought of as a sum of contributions from each such strip, and its variance has the value $`\delta g^2=L_1^2/3L^2`$.
Our concern in the following is with samples that are not completely phase coherent. They have two different metallic regimes, analogous to those described above for phase coherent conductors. We restrict our attention to one of these, the incoherent two-dimensional metal, in which $`\tau _{in}<C/v`$ so that dephasing occurs before circumnavigation: this is the simpler regime theoretically and seems the one relevant for existing experiments. Cho, Balents and Fisher have estimated the size of conductance fluctuations in the incoherent metal by viewing each phase-coherent region as a classical resistor and the whole sample as a random resistor network. We summarise their argument here for convenience. Let the number of regions into which the sample is divided be $`n_x`$ in the chiral direction, and $`n_z`$ in the transverse direction. Let the conductance of a single region have mean value $`g_{\mathrm{patch}}`$, with fluctuations of magnitude $`\delta g_{\mathrm{patch}}`$; provided $`\delta g_{\mathrm{patch}}g_{\mathrm{patch}}`$, the corresponding resistances are $`R_{\mathrm{patch}}=1/g_{\mathrm{patch}}`$ and $`\delta R_{\mathrm{patch}}=\delta g_{\mathrm{patch}}/g_{\mathrm{patch}}^2`$. A strip of $`n_z`$ such regions forms a single conduction path in the transverse direction, having mean resistance $`R_{\mathrm{strip}}=n_zR_{\mathrm{patch}}`$ with fluctuations $`\delta R_{\mathrm{strip}}^2=n_z\delta R_{\mathrm{patch}}^2`$; the corresponding conductances are $`g_{\mathrm{strip}}=1/R_{\mathrm{strip}}`$ and $`\delta g_{\mathrm{strip}}=\delta R_{\mathrm{strip}}/R_{\mathrm{strip}}^2`$. Since the sample consists of $`n_x`$ such paths in parallel, its conductance has mean value $`g=n_xg_{\mathrm{strip}}=(n_x/n_z)g_{\mathrm{patch}}`$ and fluctuations $`\delta g^2=n_x\delta g_{\mathrm{strip}}^2=(n_x/n_z^3)\delta g_{\mathrm{patch}}^2`$. The implications of this final result depend on the size of a patch in the transverse direction, compared to the layer spacing, $`a`$.
If $`\tau _{in}^1\tau _{}^1`$ so that the phase coherence length in the transverse direction is much larger than the layer spacing, which is the case considered previously , then one expects a continuum treatment of the system to be adequate. In this case, $`n_z^2=L^2/D\tau _{in}`$ and $`\delta g_{\mathrm{patch}}^21`$; in addition, $`n_x=C/v\tau _{in}`$. Combining these expressions and introducing the inelastic scattering length in the chiral direction, $`l_{in}=v\tau _{in}`$, one arrives at the conclusion (Eq. (4.1) of Ref. ) that
$$\delta g^2\left(\frac{l_{in}}{C}\right)^{1/2}\left(\frac{g}{N}\right)^{3/2}.$$
(1)
Experimental studies use a magnetic field with a component normal to the sample surface to generate conductance fluctuations in the chiral metal ; the field scale of conductance correlations is set by the flux quantum, $`\mathrm{\Phi }_0`$, and the area of a coherent region, being (Eq. (4.4) of Ref. )
$$\delta B_0\mathrm{\Phi }_0/(l_{in}^{3/2}a^{1/2}\sigma ).$$
(2)
The results we obtain in Sec. III substitute for Eqns. (1) and (2) the correlation function
$`\delta g(B_{})\delta g(B_{}+\delta B)=\left({\displaystyle \frac{g}{N}}\right)^{3/2}\left({\displaystyle \frac{l_{in}}{C}}\right)^{1/2}f(y)`$ (3)
(4)
where the function $`f(y)`$ is
$$f(y)=\pi ^{1/2}_{\mathrm{}}^{\mathrm{}}𝑑xe^{(x^2+yx^6)}$$
(5)
with the scaling variable
$$y=\frac{\pi ^2}{12}\sigma al_{in}^3\left(\frac{\delta B}{\mathrm{\Phi }_0}\right)^2.$$
(6)
The opposite limit of weakly coupled edges, $`\tau _{in}^1\tau _{}^1`$, has not previously been examined theoretically but seems to be important experimentally . Adapting the argument outlined above, we take, as elements of a classical resistor network, regions again of length $`l_{in}`$ in the chiral direction but now of width $`a`$ in the transverse direction. Then $`n_x=C/l_{in}`$ and $`n_z=N`$. The mean conductance of one such patch is small in this limit. We expect fluctuations to be of the same order, and take $`\delta g_{\mathrm{patch}}^2g_{\mathrm{patch}}^2=(gn_z/n_x)^2`$. In this way we obtain in place of Eq. (1)
$$\delta g^2\frac{g^2}{N}\frac{l_{in}}{C}.$$
(7)
The field scale of conductance correlations is again set by the flux quantum and the area of a coherent region, but Eq. 2 is replaced in this limit by
$$\delta B_0\mathrm{\Phi }_0/(l_{in}a).$$
(8)
We calculate in Sec. III the conductance correlation function for weakly coupled edges, finding in agreement with these estimates
$$\delta g(B_{})\delta g(B_{}+\delta B)=\frac{2g^2}{NC}\frac{l_{in}}{1+z^2}$$
(9)
where $`z=2\pi \delta Bl_{in}a/\mathrm{\Phi }_0`$.
The inelastic scattering length, $`l_{in}`$, appears as a phenomenological parameter in the expressions given above for the correlation function of conductance fluctuations. The contribution to inelastic scattering from electron-electron interactions is known in non-chiral dirty metals to have a characteristic temperature dependence, reflecting the form of the dynamically screened Coulomb interaction in a diffusive system , and it is interesting to study inelastic scattering microscopically for the chiral metal. We do this in Sec. IV, examining dynamical screening and using the results to calculate the inelastic scattering rate, $`\tau _{in}^1`$. We obtain in the regime $`\tau _{in}^1\tau _{}^1`$
$$\tau _{in}^1=c\frac{a}{D^{1/2}}\left(\frac{k_\mathrm{B}T}{\mathrm{}}\right)^{3/2},$$
(10)
where $`c1.5`$ is a dimensionless coefficient. A simple interpretation of this result can be given, following similar arguments established in the theory of diffusive metals . As a starting point, note that (in both chiral and non-chiral diffusive metals) the dynamically screened interaction strength at long wavelengths is independent of the bare interaction strength. In consequence, the inelastic scattering rate depends only on properties of the non-interacting system and temperature: $`\mathrm{}\tau _{in}^1`$ is of order the single-particle energy level spacing in a system with dimensions determined by the distance an electron moves in the time $`\mathrm{}/k_\mathrm{B}T`$. These dimensions are $`L_x(T)\mathrm{}v/k_\mathrm{B}T`$ and $`L_z(T)(\mathrm{}D/k_\mathrm{B}T)^{1/2}`$ in the chiral and transverse directions respectively, and the estimate $`\tau _{in}^1(\mathrm{}nL_x(T)L_z(T))^1)`$ is consistent with Eq. (10). An equivalent argument applied to non-chiral, diffusive conductors in $`d`$-dimensions yields $`\tau _{in}^1T^{d/2}`$, which is corrected in two dimensions to $`\tau _{in}^1T|\mathrm{log}(T)|`$ by detailed calculations . The factor of $`\mathrm{log}(T)`$ in this last expression reflects divergent contributions to scattering in the diffusive two-dimensional metal, from processes with small energy transfer. The same scattering processes are responsible for a difference, in diffusive two-dimensional metals, between the temperature dependence of the inelastic (or out-scattering) rate, which acts as a cut-off for conductance fluctuations, and the dephasing rate, which acts as a cut-off for the weak localisation effects that result in negative magnetoresistance. In the chiral metal, as for conventional, diffusive metals in more than two dimensions, small energy transfer processes are not dominant and there is a single relaxation rate.
Inelastic scattering in weakly coupled edge states has been discussed in Ref. , where it is noted that one expects $`\tau _{in}^1T`$ from perturbation theory in the interaction stength, applied to a single edge.
## III Conductance Fluctuations
### A Model
We study conductance fluctuations in the chiral metal using a single-particle description of the system, supplemented by a relaxation rate to represent inelastic scattering. We treat a sample with layers lying in the $`xy`$ plane and consider a surface in the $`xz`$ plane. Electrons on the surface then have a continuous coordinate, $`x`$, in the chiral direction parallel to the edges of layers, and an integer coordinate, $`n`$, labeling layers in the transverse ($`z`$) direction, which we combine as $`𝐫=(x,n)`$. We take the interlayer tunneling energy to be $`t`$ and represent a magnetic field component, $`B_{}`$, normal to the surface using the vector potential $`𝐀=B_{}an\widehat{𝐱}`$. The Hamiltonian $`H`$ acts on a wavefunction $`\psi _n(x)`$ according to
$`(H\psi )_n(x)=v(i\mathrm{}_x+eB_{}an)\psi _n(x)`$ (11)
$`t[\psi _{n+1}(x)+\psi _{n1}(x)]+V_n(x)\psi _n(x),`$ (12)
where $`V_n(x)`$ is a random potential arising from impurities and surface roughness. We choose $`V_n(x)`$ to be Gaussian distributed with short-range correlations, so that $`V_n(x)=0`$ and $`V_n(x)V_m(x^{})=\mathrm{\Delta }\delta _{nm}\delta (xx^{})`$.
Denoting the Green’s function for Eq. (12) by $`g(zH)^1`$, we require its disorder average
$$G(z;𝐫_1,𝐫_2)g(z;𝐫_1,𝐫_2),$$
(14)
the diffusion propagator
$$K(\omega ;𝐫_1,𝐫_2)g(\omega +i0;𝐫_1,𝐫_2)g(i0;𝐫_2,𝐫_1),$$
(15)
and its Fourier tranform
$$K(\omega ,𝐤)_{\mathrm{}}^{\mathrm{}}𝑑x\underset{n}{}e^{i(xk_x+ank_y)}K(\omega ;\mathrm{𝟎},𝐫).$$
(16)
These quantities are known from calculations using the usual expansion in powers of the Green’s function for the disorder-free system. Without disorder or a normal magnetic field component, eigenstates are plane waves and the electron dispersion relation is $`ϵ(𝐤)=\mathrm{}vk_x2tcos(k_za)`$. This dispersion along with the chiral motion leads to an open Fermi surface on which all electrons have the same $`x`$-component of velocity. In the presence of a normal magnetic field component of dimensionless strength $`beB_{}va/t`$, the eigenfunctions $`\psi _n(x)`$ in the pure system are
$$\psi _n(x)=\frac{1}{\sqrt{2\pi }}e^{ikx}\varphi _\alpha (n),$$
(17)
where $`\alpha `$ is an integer and $`\varphi _\alpha (n)=J_{n\alpha }(2/b)`$, the Bessel function of order $`n\alpha `$. The associated energy eigenvalues are $`\mathrm{}vk+ϵ_\alpha `$ with $`ϵ_\alpha =\alpha bt`$.
The treatment of disorder is simple in some respects because chiral motion prevents multiple scattering of an electron from any particular impurity. As a result, the average one-particle Green’s function is given exactly by the Born approximation. In this way, one arrives at
$`G^R(E;𝐫_1,𝐫_2)`$ $`=`$ (18)
$`{\displaystyle \frac{1}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑k{\displaystyle \underset{\alpha }{}}`$ $`{\displaystyle \frac{e^{ik(x_2x_1)}\varphi _\alpha (n_1)\varphi _\alpha (n_2)}{E+i\mathrm{\Delta }/(2\mathrm{}v)(\mathrm{}vk+ϵ_\alpha )}},`$ (19)
where we use the standard notation, $`G^R(E)G(E+i0)`$ and $`G^A(E)G(Ei0)`$. From the value of the self-energy, the elastic scattering time is $`\tau _{el}=\mathrm{}^2v/\mathrm{\Delta }`$, and the elastic scattering length is $`l_{el}=(\mathrm{}v)^2/\mathrm{\Delta }`$. In the absence of a normal magnetic field component, $`G^R(E;𝐫_1,𝐫_2)`$ is translationally invariant and has the Fourier transform
$$G^R(E;𝐤)=[E+i\mathrm{\Delta }/(2\mathrm{}v)(\mathrm{}vk_x2t\mathrm{cos}(k_za)]^1.$$
(20)
For non-zero $`B_{}`$, $`G^R(E;𝐫_1,𝐫_2)`$ acquires a phase under translations, which is central to calculations of conductance fluctuations with magnetic field:
$$G^R(E;\mathrm{𝟎},𝐫)=\mathrm{exp}(i\beta an^{}x)G^R(E;𝐫^{},𝐫+𝐫^{}),$$
(21)
where $`\beta =2\pi B_{}/\mathrm{\Phi }_0`$, $`𝐫=(x,n)`$ and $`𝐫^{}=(x^{},n^{})`$.
The diffusion propagator is similarly given by a a sum of ladder diagrams, and has a form that reflects the combination of chiral motion in $`x`$ and diffusion in $`z`$. At small wavevectors it is
$$K(\omega ,𝐤)=(\mathrm{}v)^1[\mathrm{}Dk_{z}^{}{}_{}{}^{2}i(\omega +\mathrm{}vk_x)]^1.$$
(22)
The diffusion constant has the value $`D=2(at)^2v/\mathrm{\Delta }`$ in the absence of a normal magnetic field component, and has the field dependence $`D(B_{})=D(0)/[1+(B_{}/B_0)^2]`$, where the field scale for magnetoresistance is $`B_0=\mathrm{\Phi }_0/al_{el}`$.
Because of the simplifications in the treatment of disorder that follow from chiral motion, the results above hold both at weak disorder ($`t\mathrm{\Delta }/(\mathrm{}v)`$) and at strong disorder ($`t\mathrm{\Delta }/(\mathrm{}v)`$). In a semiclassical description, electron motion follows a random walk in the transverse direction. At weak disorder, steps of the walk have speed $`at/\mathrm{}`$ and duration $`\tau _{el}`$; the rate for inter-edge tunneling is $`\tau _{}^1=t/\mathrm{}`$. At strong disorder, steps are of length $`a`$ and duration $`\tau _{}`$; from the value of $`D`$, we deduce that the rate for inter-edge tunneling is then $`\tau _{}^1=(t/\mathrm{})^2\tau _{el}`$.
### B Conductivity
The component of the current density operator in the transverse direction is
$$j_z(𝐫)=\frac{eat}{i\mathrm{}}[\psi _n^{}(x)\psi _{n+1}(x)\psi _{n+1}^{}(x)\psi _n(x)].$$
(23)
The Kubo formula for the transverse conductivity in a system of area $`A`$ is
$$\sigma =\frac{\mathrm{}}{2\pi A}Tr[j_zg^Rj_zg^A].$$
(24)
Evaluating the disorder average of this expression, one finds
$$\sigma =\frac{e^2}{h}\frac{D}{va}$$
(25)
as expected from the Einstein relation.
### C Conductance fluctuations in magnetic field
Conductance fluctuations have been studied previously for the two-dimensional chiral metal in the absence of inelastic scattering, in most detail by Gruzberg, Read and Sachdev , who discussed the variance and its dependence on sample geometry. Here we present calculations which include inelastic scattering, obtaining the correlation function
$$F(\delta B)=\delta g(B)\delta g(B+\delta B).$$
(26)
As a result of chiral motion, and provided that electrons do not circumnavigate the sample phase coherently ($`C/v>\tau _{in}`$ or $`C/v>L^2/D`$), the only one-loop diagram contributing to $`F(\delta B)`$ is the one shown in Fig. 1. It is most conveniently evaluated in real space, and represents the expression
$`F(\delta B)={\displaystyle \underset{n_i}{}K_{\delta B}(𝐫_1,𝐫_2)K_{\delta B}^{}(𝐫_4,𝐫_3)}`$ (27)
$`J(𝐫_1,𝐫_4)J(𝐫_2,𝐫_3)d\{x_i\}.`$ (28)
Here, $`K_{\delta B}(𝐫_1,𝐫_2)`$ stands for the zero-frequency diffusion propagator, generalised to the situation of interest, in which the two Green’s functions entering it are evaluated for different normal magnetic field strengths,
$$K_{\delta B}(𝐫,𝐫^{})=g_B^R(𝐫,𝐫^{})g_{B+\delta B}^A(𝐫^{},𝐫),$$
(29)
and in which inelastic scattering is included. The other factors arise from a combination of single-particle Green’s functions and current operators, and are short range: $`J(𝐫,𝐫^{})=(\mathrm{}\mathrm{\Delta }^2/2\pi L^2)|C(𝐫,𝐫^{})|^2`$, with
$$C(𝐫,𝐫^{})=\underset{n_1}{}G^A(𝐫,𝐫_1)j_z(𝐫_1)G^R(𝐫_1,𝐫^{})dx_1.$$
(30)
Our approach to evaluating $`F(\delta B)`$ is different in the two regimes, $`\tau _{in}^1\tau _{}`$ and $`\tau _{in}^1\tau _{}`$, which we consider separately.
#### 1 Strongly coupled edges: $`\tau _{}^1\tau _{in}^1`$
In this regime, the discreteness of the system in the transverse direction may be neglected. It is sufficient to approximate the short-range terms in Eq. (28) by
$$J(𝐫,𝐫^{})=J_0\delta (𝐫𝐫^{}),$$
(31)
where
$$J_0=\frac{e^2}{h}\frac{4(at)^2}{\mathrm{\Delta }L^2}(\mathrm{}v)^2.$$
(32)
The diffusion propagator, generalised to include the magnetic field difference $`\delta B`$ and inelastic scattering, is most easily calculated within a continuum treatment of the transverse direction by solving, for $`x>0`$, the differential equation
$$v_xK_{\delta B}(\mathrm{𝟎},𝐫)=[D_z^2+iv\delta \beta z\tau _{in}^1]K_{\delta B}(\mathrm{𝟎},𝐫)$$
(33)
with the initial condition $`K_{\delta B}(\mathrm{𝟎};x,z)=(\mathrm{}v)^2\delta (z)`$ at $`x=0`$. Here, the inelastic scattering rate, $`\tau _{in}^1`$, has been introduced,
and the magnetic field difference enters through $`\delta \beta (2\pi \delta B/\mathrm{\Phi }_0)`$. To describe a sample connected to contacts, absorbing boundary conditions ($`_zK_{\delta B}(\mathrm{𝟎},𝐫)=0`$) should be imposed at $`z=0`$ and $`z=L`$. However, provided the size of a phase-coherent region is much smaller than the sample size, so that $`D\tau _{in}L^2`$, it is sufficient to use the solution of Eq. (33) for a system of infinite extent in the $`z`$-direction when evaluating the right-hand side of Eq. (28). This solution is
$$K_{\delta B}(\mathrm{𝟎},𝐫)=(\mathrm{}v)^2\left(\frac{v}{4\pi Dx}\right)^{1/2}\mathrm{exp}(S)$$
(34)
where
$$S=\frac{vz^2}{4Dx}+\frac{x}{l_{in}}\frac{i\delta \beta xz}{2}+\frac{D(\delta \beta )^2x^3}{12v}.$$
(35)
Combining Eqns. (28), (31), and (34), we obtain
$$F(\delta B)=J_0^2LC_0^{\mathrm{}}𝑑x_{\mathrm{}}^{\mathrm{}}𝑑z|K_{\delta B}(\mathrm{𝟎},𝐫)|^2$$
(36)
and so find the correlation function for conductance fluctuations given in Eq.(4)
The form of the scaling function $`f(y)`$ is shown in Fig. 2. It is normalised so that $`f(0)=1`$, and decreases as $`fy^{1/6}`$ for $`y>>1`$, so that the conductance correlation function decays as $`F(\delta B)|\delta B|^{1/3}`$ for large $`|\delta B|`$. FIG. 2.: The scaling function $`f(y)`$
#### 2 Weakly coupled edges: $`\tau _{in}^1\tau _{}^1`$
If adjacent edges are sufficiently weakly coupled by tunneling, there is mesoscopic regime in which the inelastic scattering rate is small, in the sense that $`\tau _{in}^1\tau _{el}^1`$, but the interlayer tunneling rate is even smaller. It is sufficient in this regime to calculate the conductance correlation function at leading order in the interlayer coupling, $`t`$. Doing so, we retain $`t`$ in the current operators but use Green’s functions for a system of isolated edges. The disorder-averaged single-particle Green’s function in this limit
is
$`G^R(E;𝐫_1,𝐫_2)=\mathrm{\Theta }(x_2x_1)\delta _{n_1,n_2}\times `$ (37)
$`(i\mathrm{}v)^1\mathrm{exp}([(2l_{in})^1i\beta an_1][x_2x_1]),`$ (38)
where $`\mathrm{\Theta }(x)`$ is the step function. The short-range terms associated with current operators in Eq. (28) may then be approximated as
$$J(𝐫,𝐫^{})=\frac{J_0}{2}\delta (xx^{})[\delta _{n,n^{}+1}+\delta _{n,n^{}1}].$$
(39)
The diffusion propagator reduces for uncoupled edges to a function which has a phase determined by $`\delta B`$ and which decays in amplitude with distance only as a result of inelastic scattering. We take it to be
$`K_{\delta B}(𝐫_1,𝐫_2)=\mathrm{\Theta }(x_2x_1)\delta _{n_1,n_2}\times `$ (40)
$`(\mathrm{}v)^2\mathrm{exp}([l_{in}^1i\delta \beta an_1][x_2x_1]).`$ (41)
Combining Eqns. (28), (39), and (41), we obtain
$`F(\delta B)={\displaystyle \frac{1}{4}}J_0^2NC\times `$ (42)
$`{\displaystyle _0^{\mathrm{}}}dx[K_{\delta B}(0,n;x,n)K_{\delta B}^{}(0,n+1;x,n+1)+\mathrm{c}.\mathrm{c}.]`$ (43)
and hence arrive at the correlation function for conductance fluctuations given in Eq.(9).
### D Conductance fluctuations in energy
It is also of interest to consider the dependence of conductance fluctuations on Fermi energy. In particular, the range of conductance correlations in energy for a system at zero temperature will determine the extent to which the amplitude of fluctuations measured at finite temperature is reduced by thermal smearing of the electron distribution. We find that conductance fluctuations are perfectly correlated in energy. To show this, we return to Eq. (28) but replace the magnetic field difference $`\delta B`$ with an energy difference $`\delta E`$. The diffusion propagator at finite energy difference simply acquires a phase:
$$K_{\delta E}(𝐫,𝐫^{})=K(𝐫,𝐫^{})e^{i\delta E(x^{}x)/(\mathrm{}v)}.$$
(44)
This phase cancels when $`K_{\delta E}(𝐫_1,𝐫_2)`$ and $`K_{\delta E}^{}(𝐫_4,𝐫_3)`$ are combined, since the factors $`J(𝐫,𝐫^{})`$ associated with the current operators are short-range, so that $`𝐫_1𝐫_4`$ and $`𝐫_2𝐫_3`$ in Eq. (28). As a consequence, the correlation of conductance fluctuations is independent of $`\delta E`$. The physical reason for this behaviour is that, if $`\psi _n(x)`$ is an eigenfunction of Eq. (12) for some realisation of disorder, with energy $`E`$, then $`\psi _n^{}(x)e^{i\omega x/v}\psi _n(x)`$ is also an eigenfunction, with energy $`E^{}=E+\mathrm{}\omega `$, provided $`\omega C/v`$ is an integer multiple of $`2\pi `$, so that periodic boundary conditions in the chiral direction are satisfied. Hence, apart from the phase factor, which does not affect the conductance, states are perfectly correlated in energy. As a result, the only consequence of a change in temperature is a change in the inelastic scattering rate.
## IV Effects of Interactions
### A Polarization and screening
In this section we study several aspects of electron-electron interactions in the chiral metal. As a first step, it is necessary to examine screening of the Coulomb interaction between a pair of electrons, by other electrons in the Fermi sea. One can expect the frequency and wavevector dependence of the screened interaction to be significant. In conventional conductors, the consequence of disorder and the resulting diffusive motion of charge (with diffusion constant $`D`$) is that screening of a potential having wavevector $`q`$ is suppressed at frequencies higher than $`Dq^2`$. And in a one-dimensional chiral metal at the edge of a two-dimensional integer quantum Hall system, electron-electron interactions modify the dispersion relation for edge magneto-plasmons. We find in the two-dimensional chiral metal that dynamical screening at finite wavevector interpolates, as a function of the direction of the wavector, between these two types of behaviour.
Our results may be obtained either by using the random phase approximation within a diagrammatic calculation, or more transparently from a hydrodynamic approach, as follows.
Consider the response of the chiral metal to a space- and time-dependent external potential. Let $`\rho (𝐫,t)`$ be the screening charge density induced in the presence of an electric field that has components $`(E_x,E_z)`$ within the surface. We require an equation of motion for $`\rho (𝐫,t)`$, which we derive in the usual way, by considering $`𝐉(𝐫,t)`$, the deviation from equilibrium of the current density, and using the continuity equation. Treating the transverse direction as continuous, one has
$$J_x(𝐫,t)=v\rho (𝐫,t)$$
(45)
and
$$J_z(𝐫,t)=D_z\rho (𝐫,t)+\sigma E_z,$$
(46)
so that the effect of the transverse component of the electric field is simply to generate a transverse current density. The effect of the component of the electric field in the chiral direction is rather different: referring to the three-dimensional sample as a whole, $`E_x`$ generates Hall currents within each layer, which transport charge between the bulk and the surface. The resulting change in surface charge density may be represented by adding a term to the continuity equation, so that it becomes
$$_t\rho (𝐫,t)=𝐉(𝐫,t)+\frac{\sigma _H}{a}E_x,$$
(47)
where $`\sigma _H`$ is the quantised Hall conductance of a layer, and (as in preceding sections) $`a`$ is the layer spacing. It is helpful to express the conductivities in Eqns. (46) and (47) as $`\sigma =e^2nD`$ and $`\sigma _H=e^2/h`$, and to recall that the density of states is $`n=(hva)^1`$. Then, combining Eqns. (45), (46) and (47), the charge density evolves as
$`_t\rho (𝐫,t)=`$ (48)
$`v_x\rho (𝐫,t)`$ $`+D_z^2\rho (𝐫,t)+e^2n[vE_xD_zE_z].`$ (49)
The electric field appearing in these equations arises from a scalar potential $`\mathrm{\Phi }_{tot}(𝐫,t)`$ which is the sum of an externally imposed potential, $`\mathrm{\Phi }_{ext}(𝐫,t)`$, and the potential $`\mathrm{\Phi }_{scr}(𝐫,t)`$ due to the screening charge, $`\rho (𝐫,t)`$. These last two quantities are related in the standard fashion by the Poisson equation for a three-dimensional electrostatic problem, with coordinates $`𝐫=(x,z)`$ and $`y`$, in which the charge density is $`\rho (𝐫,t)\delta (y)`$ and $`\mathrm{\Phi }_{scr}(𝐫,t)`$ is found by evaluating the three-dimensional potential at $`y=0`$. Taking Fourier transforms, defined according to
$$\rho (\omega ,𝐤)d^2𝐫𝑑te^{i(𝐤𝐫+\omega t)}\rho (𝐫,t),$$
(50)
and similarly for other quantities, one finds
$$\mathrm{\Phi }_{scr}(\omega ,𝐤)=U_0(𝐤)\rho (\omega ,𝐤)/e^2,$$
(51)
where
$$U_0(𝐤)=\frac{e^2}{2ϵ_rϵ_0|𝐤|}$$
(52)
is the Fourier transform of the unscreened interaction potential betwen a pair of electrons. Expressing the electric field in terms of $`\mathrm{\Phi }_{tot}(𝐫,t)`$ and solving the Fourier transform of Eq. (49), we find
$$\mathrm{\Phi }_{tot}(\omega ,𝐤)=\frac{\mathrm{\Phi }_{ext}(\omega ,𝐤)}{1+U_0(𝐤)\mathrm{\Pi }(\omega ,𝐤)},$$
(53)
where
$$\mathrm{\Pi }(\omega ,𝐤)=n\frac{Dk_{z}^{}{}_{}{}^{2}ivk_x}{Dk_{z}^{}{}_{}{}^{2}i\omega ivk_x}.$$
(54)
This result simplifies in several limits. Static screening ($`\omega =0`$) is isotropic and as given by Thomas-Fermi theory, with an inverse screening length $`\kappa =e^2n/2ϵ_rϵ_0`$. Dynamical screening of a potential with Fourier components only in the transverse direction ($`k_x=0`$) is exactly as in a non-chiral disordered, two-dimensional metal. Response to a potential with Fourier components only in the chiral direction ($`k_z=0`$) is undamped, and excitations have the dispersion relation: $`\omega =v(k_x+\kappa )`$.
In the following, we shall treat interactions in the chiral metal using the Matsubara formalism. Within this approach, we find for the polarisation operator
$$\mathrm{\Pi }(i\omega _n,𝐤)=n\frac{Dk_{z}^{}{}_{}{}^{2}ivk_x\mathrm{sgn}(\omega _n)}{Dk_{z}^{}{}_{}{}^{2}+|\omega _n|ivk_x\mathrm{sgn}(\omega _n)}$$
(55)
with $`\omega _n=2\pi nk_BT`$, which has, as its analytical continuation, Eq. (54). Similarly, the screened Coulomb interaction potential is
$$U_{scr}(i\omega _n,𝐤)=U_0(𝐤)/[1+U_0(𝐤)\mathrm{\Pi }(i\omega _n,𝐤)].$$
(56)
### B Tunneling density of states
The enhanced interaction between slowly relaxing density fluctuations in a conventional, diffusive metal is responsible for a zero-bias anomaly in the tunneling density of states . Balents and Fisher have argued that such a zero-bias anomaly should not be expected in the chiral metal. Physically, charge that tunnels into the system is swept away at the chiral velocity; mathematically, the terms in perturbation theory associated with this anomaly yield only smooth contributions for the chiral metal.
In the present subsection we examine in more detail the reasons for this. It is sufficient for the purpose to consider a static interaction, $`U(𝐫)`$, with finite range, between electrons separated by a distance $`𝐫`$, and to calculate at first order in $`U(𝐫)`$ the disorder-averaged exchange energy . Let $`\mathrm{\Sigma }(E)`$ be the disorder-average of this quantity for an electron at energy $`E`$ interacting with a filled Fermi sea . The change in the density of states due to interactions is
$$\delta n(E)=n\frac{d\mathrm{\Sigma }(E)}{dE}.$$
(57)
Viewed in this way, the zero-bias anomaly in conventional, diffusive metals arises because states lying close in energy have wavefunctions which are highly correlated in space, enhancing interactions. Moverover, the wavefunction correlations depend strongly on the energy separation of the states, so that $`\mathrm{\Sigma }(E)`$ varies rapidly with $`E`$ if $`E`$ is close to the Fermi energy, $`E_\mathrm{F}`$.
Let $`\{\mathrm{\Psi }_\alpha (𝐫)\}`$ and $`\{E_\alpha \}`$ be single-particle eigenfunctions and eigenenergies in a disordered chiral metal. The disorder-averaged exchange energy is
$$\mathrm{\Sigma }(E)=\frac{1}{n}_{\mathrm{}}^{E_\mathrm{F}E}𝑑\omega d^2𝐫U(𝐫)A(\omega ,𝐫),$$
(58)
where the two-particle spectral function, $`A(\omega ,𝐫)`$, is defined by
$`A(\omega ,𝐫)={\displaystyle \underset{\alpha \beta }{}}\delta (E_\alpha )\delta (\omega E_\beta )\mathrm{\Psi }_\alpha ^{}(\mathrm{𝟎})\mathrm{\Psi }_\alpha (𝐫)\mathrm{\Psi }_\beta ^{}(𝐫)\mathrm{\Psi }_\beta (\mathrm{𝟎}).`$ (59)
This spectral function can be calculated from the real part of the diffusion propagator, Eq. (15), or by considering the time-Fourier tranform of the density associated with a spreading wavepacket. It has the form
$`A(\omega ,𝐫)`$ $`=`$ (60)
$`{\displaystyle \frac{1}{a(\mathrm{}v)^2}}`$ $`\left({\displaystyle \frac{v}{4\pi D|x|}}\right)^{1/2}\mathrm{exp}({\displaystyle \frac{vz^2}{4D|x|}}{\displaystyle \frac{i\omega x}{v}}).`$ (61)
Combining Eqns. (57), (58) and (61), one sees that, if the potential has range $`R`$, then $`\delta n(E)`$ is independent of $`E`$ for $`|EE_\mathrm{F}|<\mathrm{}v/R`$.
Inspecting Eq. (61), it is clear that eigenfunctions in the chiral metal do in fact have strong correlations: $`A(\omega ,𝐫)`$ with $`𝐫=(x,z)`$ diverges as $`|x|0`$ for $`z=0`$. These correlations, however, depend only weakly on energy, and so do not produce a large change in the tunneling density of states.
### C Inelastic scattering rate
The following subsection is devoted to a calculation of the contribution from electron-electron interactions to the inelastic scattering rate, $`\tau _{in}^1`$, which acts as a cut-off for conductance fluctuations.
We compute $`\tau _{in}^1`$ by finding at leading order the interaction contribution to the irreducible vertex that appears in the ladder sum for the diffusion propagator used in calculations of conductance fluctuations . The two single-particle (advanced and retarded) Green’s functions from which this diffusion propagator is constructed are associated with distinct measurements of the conductance. In consequence, one should not include in the irreducible vertex diagrams in which the interaction line connects the two single-particle propagators. The relevant diagrams are shown in Fig. 3, in which the full lines represent single-particle propagators, wavy lines are the screened Coulomb interaction, single dashed lines are impurity averages, and double dashed lines are diffusion ladders. We evaluate these diagrams treating the transverse direction as continuous, so that our results are restricted to the regime in which edges are strongly coupled: $`\tau _{}^1\tau _{in}^1`$. We set the energy and wavevector differences, $`\omega `$ and $`𝐪`$, at the vertex to zero, express sums on Matsubara frequencies as contour integrals, and combine terms. We expect that the dominant contribution to the resulting integral will be from small wavevectors and energies, and use the asymptotic form of the integrand appropriate for this regime, obtaining
$`\tau _{in}^1={\displaystyle \frac{1}{\pi \mathrm{}n}}{\displaystyle }{\displaystyle \frac{d^2𝐪}{(2\pi )^2}}{\displaystyle _{\mathrm{}}^+\mathrm{}}dx\times `$ (62)
$`[\mathrm{coth}(x/2k_\mathrm{B}T)\mathrm{tanh}((x)/2k_\mathrm{B}T)]\times `$ (63)
$`{\displaystyle \frac{x(\mathrm{}Dq_{z}^{}{}_{}{}^{2})^2}{[(x+\mathrm{}vq_x)^2+(\mathrm{}Dq_{z}^{}{}_{}{}^{2})^2][(\mathrm{}vq_x)^2+(\mathrm{}Dq_{z}^{}{}_{}{}^{2})^2]}}.`$ (64)
The dependence of this expression on $`T`$, $`D`$ and $`v`$ is made apparent by introducing the scaled variables: $`X=x/k_\mathrm{B}T`$, $`Q_x=\mathrm{}vq_x/k_\mathrm{B}T`$ and $`Q_z=(\mathrm{}D/k_\mathrm{B}T)^{1/2}q_z`$. The integral can be evaluated numerically, yielding the expression for $`\tau _{in}`$ given in Eq. (10). The physical interpretation of this result has been discussed in Sec. II.
### D Interaction corrections to conductivity
Electron-electron interactions are known in diffusive conductors to give rise to temperature-dependent contributions to the conductivity that are singular in the low-temperature limit The origin of these interaction corrections to the conductivity is similar to that of the zero-bias anomaly in the tunneling density of states. Since the chiral metal has no zero-bias anomaly, one expects it to have no singular interaction contribution to the conductivity. In the following subsection we outline calculations that confirm this at leading order. The presentation follows closely that of Altshuler, Khmel’nitskii, Larkin, and Lee for the diffusive metal.
To compute the interaction correction to the conductivity of the chiral metal using the Matsubara technique, it is convenient to use the Kubo formula in the form
$$\sigma =\underset{\omega 0}{lim}\frac{i\mathrm{\Lambda }(\omega )}{\omega }$$
(65)
where $`\mathrm{\Lambda }(\omega )`$ is the analytic continuation in $`\omega `$ of
$$\mathrm{\Lambda }(i\omega _n)=\frac{k_\mathrm{B}T}{A}\underset{ϵ_m}{}Tr[j_zg(i[ϵ_m+\omega _n])j_zg(iϵ_m)].$$
(66)
The diagrams that contribute to $`\mathrm{\Lambda }(i\omega _n)`$ at leading order in the screened interaction strength can be classified according to the number of diffusion propagators each one contains. Those with the most diffusion propagators are most singular in the low-temperature limit. They can be separated into two categories. Diagrams in the first category are illustrated in Fig. 4, corresponding to Figs. 5(a), 5(b) and 5(c) of Ref. . It has been shown for a diffusive conductor in Ref. that these cancel, and we find the same to be true in the chiral metal. Diagrams in the second category are illustrated in Fig. 5, corresponding to Figs. 5(d) and 5(e) of Ref. . For the two-dimensional diffusive metal, these diagrams make a contribution to the conductivity which is logarithmic in temperature. For the chiral metal, their contribution is formally proportional to $`T^{1/2}`$, but with a numerical coefficient which we find to be zero after integrating over the component in the chiral direction of the momentum transferred by the interaction.
## V Discussion
Our results, as set out in Sec. II, are a contribution to the understanding of transport in the chiral metal. In many senses, it is a particularly simple example of a disordered conductor. Quantum interference effects are suppressed by chiral motion, and interaction effects are less singular than in diffusive two-dimensional conductors. We have provided a rather detailed treatment of conductance fluctations, which we hope will be tested experimentally. We have also identified a distinct mesoscopic regime, for weakly coupled edges, in which the inter-edge tunneling rate, $`\tau _{}^1`$, is smaller than the inelastic scattering rate, $`\tau _{in}^1`$. In this regime, transport is controlled by the properties of an isolated, one-dimensional chiral metal.
Recent experiments appear, in fact, to probe behaviour of weakly coupled edges, though their analysis has necessarily so far been based on theory for the opposite regime. According to that analysis, the inelastic scattering length in the transverse direction, estimated from the variance of conductance fluctuations using Eq. (1), is shorter than the inter-layer spacing, even at the lowest temperatures investigated. Moreover, this estimate of the inelastic scattering length is inconsistent with the value deduced from the correlation field for condcutance fluctuations, using Eq. (2), as one might expect if the theory used is, in fact, inappropriate. It will be interesting to find whether these discrepancies are resolved by using the theory we have presented to analyse the experiments. If so, such measurements provide a way to investigate interaction effects in a one-dimensional system of chiral fermions.
## Acknowledgments
We are particularly grateful for discussions and correspondence with D. P. Druist and E. G. Gwinn. We also thank L. Balents, F. H. L. Essler, I. A. Gruzberg, R. A. Smith, and A. M. Tsvelik for useful discussions. J.B. acknowledges partial financial support from the European Union through the Marie Curie Grant; the work was also supported in part by EPSRC under Grant No GR/J8327.
## References
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# Quantum–Mechanical Detection of Non–Newtonian Gravity.
## 1 Introduction
General Relativity (GR) is one of the milestones of modern physics, and currently many of its predictions have been already tested. For instance, we already have the following results: gravitational time dilation measurement , gravitational deflection of electromagnetic waves , time delay of electromagnetic waves in the field of the sun , or the geodetic effect . Neverwithstanding, at this point it is also important to comment that all these impressive direct confirmations of GR are confirmations of weak field corrections to the Galilei–Newton mechanics.
The discovery of the first binary pulsar PSR1913+16 allowed to probe the propagation properties of the gravitational field , the results between theory and experiment agree at a level of $`10^3`$. The possibilities that binary pulsars offer do not finish here, they can also be used as laboratories for testing strong–field gravity . Concerning binary pulsars at this point it is noteworthy to mention that they are a confirmation of general relativity done at the classical level, here we mean that the observations and predictions comprise the orbital dynamics of a binary pulsar, for instance, orbital period, eccentricity .
Therefore, if GR is so successful, then why should we need analyze some possible deviation of the Newtonian inverse–square force law?. The answer stems from the fact that the agreement between general relativity and experiment might be compatible with the existence of a scalar contribution to gravity, such as a dilaton field .
This dilaton field emerges in several theoretical attempts that try to formulate a unified theory of elementary particle physics. As one of their consequences they predict the existence of new forces (which are usually refered to as “fifth force”), whose effects extend over macroscopic distances . In some ways, these new forces simulate the effects of gravity, but a crucial point is that they are not described by an inverse–square law, and even more, they, generally, violate the Weak Equivalence Principle (WEP) . Hence the presence of this kind of forces, coexisting with gravity, could be detected, in principle, by apparent deviations from the inverse–square law, or from the violation of WEP. Hence, a strong theoretical motivation for analyzing possible deviations from Newtonian gravity is to probe for new fundamental forces in nature.
To date, after more than a decade of experiments , there is no compelling evidence for any kind of deviations from the predictions of Newtonian gravity. But Gibbons and Whiting (GW) phenomenological analysis of gravity data has proved that the very precise agreement between the predictions of Newtonian gravity and observation for planetary motion does not preclude the existence of large non–Newtonian effects over smaller distance scales, i.e., precise experiments over one scale do not necessarily constrain gravity over another scale.
GW results conclude that the current experimental constraints over possible deviations did not severly test Newtonian gravity over the $`10`$$`1000`$m distance scale, usually called “geophysical window”.
Recently, a new test of the equivalence principle was carried out , the one sets new constraints on the possible ranges of a Yukawa term. This new experiment improves the current limit for ranges between 10km and 1000km. Neverwithstanding, in the short range it can say nothing about distances smaller than 1cm. Nevertheless, this experiment is performed on a classical system, namely, a 3 ton $`{}_{}{}^{238}U`$ attractor rotates around a torsion balance, which contains Cu and Pb macroscopical test bodies. In this experiment the differential acceleration of the test bodies toward the attractor was measured.
Finally, we must also mention the experiment been already carried out at the University of Padova . This proposal measures the displacements induced by an oscillating mass acting as a source of gravitational field on a micromechanical resonator. This device could give information about scalar interactions in the range below 1mm.
Among the models that in the direction of noninverse–square forces currently exist we have Fujii’s proposal , in which a “fifth force”, coexisting simultaneously with gravity, comprises a modified Newtonian potential with a Yukawa term, $`V(r)=G_{\mathrm{}}\frac{mM}{r}\left(1+\alpha e^{\frac{r}{\lambda }}\right)`$, here $`G_{\mathrm{}}`$ describes the interaction between $`m`$ and $`M`$ in the limit case $`r\mathrm{}`$, i.e., $`G=G_{\mathrm{}}(1+\alpha )`$, where $`G`$ is the Newtonian gravitational constant. This kind of deviation terms arise from the exchange of a single new quantum of mass $`m_5`$, the Compton wavelength of the exchanged field is $`\lambda =\frac{\mathrm{}}{m_5c}`$ , this field is usually called dilaton.
If we take a look at the experimental efforts that have been done in order to test the inverse–square law we will find that they can be separated into two large classes: (i) those experiments which involve the direct measurement of the magnitude $`G(r)`$, they compare preexisting laboratory Cavendish measurements of $`G`$ ; and (ii) the direct measurement of $`G(r)`$ with $`r`$ . A relevant characteristic of these efforts has to be mentioned, they remain always at the classical level, the action of the Yukawa term is always on classical systems, namely, classical test masses (Cavendish case), or in the case of mine and Borehole experiments, once again, classical test particles are employed. One of the exceptions around this topic is the use of the Casimir effect , here Planck constant, $`\mathrm{}`$, appears as a parameter in the experiment, another quantum analysis may be found in . Neverwithstanding, the existence of retardation forces, such as van der Waals forces, complicates these classical experimental constructions .
In this work we will try to explore the theoretical predictions, at quantum level, that a Yukawa term could have. This will be done resorting to an experimental proposal which is very similar to the Colella, Overhauser, and Werner (COW) construction . As we already know, COW allows the detection of the gravity–induced phase difference between the amplitudes of two wave packets arriving at a certain detection point. Having this in hindsight, we could wonder if the parameters $`\alpha `$ and $`\lambda `$, appearing in Fujii’s model, could render a detectable effect in this type of COW experiment. In other words, we will calculate the non–Newtonian gravity–induced interference between the amplitudes of two wave packets and compare this result with COW. It will be found that the difference depends on $`\alpha `$ and $`\lambda `$, and therefore could be detected, at least in principle. It gives also the possibility of measuring the mass of the dilaton field. Afterwards, we will consider the continuous monitoring of the position of the two particle beams and see that, in the context of the restricted path integral formalism (RPIF) , $`\alpha `$ and $`\lambda `$ do appear explicitly in the resulting interference term, and hence we obtain, comparing with the corresponding results of the Newtonian case , an additional method to determine these non–Newtonian parameters.
## 2 Non–Newtonian Gravity–Induced Interference
As was already mentioned above, let us now consider the case of a Yukawa modification to the Newtonian gravitational potential
$`V(r)=G_{\mathrm{}}{\displaystyle \frac{mM}{r}}\left(1+\alpha e^{\frac{r}{\lambda }}\right).`$ (1)
The Lagrangian of a particle with mass $`m`$, moving in this field, is
$`L={\displaystyle \frac{m}{2}}\dot{\stackrel{}{r}}^2+G_{\mathrm{}}{\displaystyle \frac{mM}{r}}\left(1+\alpha e^{\frac{r}{\lambda }}\right).`$ (2)
Let us now write $`r=R+l`$, where $`R`$ is the Earth’s radius, and $`l`$ the height over the Earth’s surface. Therefore, keeping terms up to second order in $`l`$, we find
$`L={\displaystyle \frac{m}{2}}\dot{\stackrel{}{r}}^2+G_{\mathrm{}}{\displaystyle \frac{mM}{R}}\left(\left[1+\alpha +{\displaystyle \frac{\alpha R}{\lambda }}\left({\displaystyle \frac{R}{2\lambda }}1\right)\right]\left[{\displaystyle \frac{1+\alpha }{R}}{\displaystyle \frac{\alpha R}{2\lambda ^2}}\right]l+{\displaystyle \frac{1+\alpha }{R^2}}l^2\right).`$ (3)
Let us now consider the case in which we perform an experiment similar to COW , i.e., two particles, starting at point $`P`$, move along two different trajectories, $`C`$ and $`\stackrel{~}{C}`$, and afterwards they are detected at a certain point $`Q`$. Here we assume that the size of the wavelengths of the packets is much smaller than the size in which the field changes considerably (i.e., we are always in the short wavelength limit), and in consequence we may consider a semiclassical approach in the analysis of the wave function.
Hence the wave function is given by the following expression
$`\psi (\stackrel{}{r},t){\displaystyle \frac{1}{[EV(\stackrel{}{r})]^{\frac{1}{4}}}}\mathrm{exp}\left\{\pm {\displaystyle \frac{i}{\mathrm{}}}{\displaystyle _{(P)}^{(Q)}}\sqrt{2m[EV(\stackrel{}{r})]}𝑑\stackrel{~}{L}{\displaystyle \frac{i}{\mathrm{}}}Et\right\},`$ (4)
where $`V(\stackrel{}{r})=G_{\mathrm{}}\frac{mM}{R}\left([1+\alpha +\frac{\alpha R}{\lambda }(\frac{R}{2\lambda }1)][\frac{1+\alpha }{R}\frac{\alpha R}{2\lambda ^2}]l+\frac{1+\alpha }{R^2}l^2\right)`$. Here the line integral appearing in expression (4) has to be calculated along $`C`$ and $`\stackrel{~}{C}`$, because we have two different trajectories.
Clearly, the interference term at the detection point, $`Q`$, is
$`I=\mathrm{cos}\left\{{\displaystyle \frac{1}{\mathrm{}}}{\displaystyle _{(C)}}\sqrt{2m[EV(\stackrel{}{r})]}𝑑\stackrel{~}{L}{\displaystyle \frac{1}{\mathrm{}}}{\displaystyle _{(\stackrel{~}{C})}}\sqrt{2m[EV(\stackrel{}{r})]}𝑑\stackrel{~}{L}\right\}.`$ (5)
Let us now consider the following trajectories. $`C`$ is defined as follows, the particle begins at point $`P`$, then moves horizontally to point $`A`$, and finally, in vertical form to point $`Q`$, which is the detection point. On the other hand, $`\stackrel{~}{C}`$ comprises the following cases, it also starts at $`P`$, but it moves, vertically, to point $`B`$, and afterwards, horizontally to $`Q`$. We also assume that the $`l`$–coordinate of point $`P`$ is zero, i.e., $`l_P=0`$, the horizontal distance between points $`B`$ and $`Q`$, and between points $`P`$ and $`A`$, is denoted by $`L`$, and finally $`l_Q`$ is the $`l`$–coordinate of $`Q`$.
Under these conditions we obtain
$`I=\mathrm{cos}[{\displaystyle \frac{L}{\mathrm{}}}(2mE+2G_{\mathrm{}}{\displaystyle \frac{m^2M}{R}}[1+\alpha +{\displaystyle \frac{\alpha R}{\lambda }}({\displaystyle \frac{R}{2\lambda }}1)]`$
$`+2G_{\mathrm{}}{\displaystyle \frac{m^2M}{R}}[{\displaystyle \frac{\alpha R}{2\lambda ^2}}{\displaystyle \frac{1+\alpha }{R}}]l_Q+2G_{\mathrm{}}{\displaystyle \frac{m^2M}{R^3}}[1+\alpha ]l_Q^2)^{\frac{1}{2}}`$
$`{\displaystyle \frac{L}{\mathrm{}}}(2mE+2G_{\mathrm{}}{\displaystyle \frac{m^2M}{R}}[1+\alpha +{\displaystyle \frac{\alpha R}{\lambda }}({\displaystyle \frac{R}{2\lambda }}1)])^{\frac{1}{2}}].`$ (6)
This last expression can be rewritten as
$`I=\mathrm{cos}\left\{{\displaystyle \frac{gm^2Ll_Q\mathrm{\Lambda }}{\mathrm{}^2}}\left[1{\displaystyle \frac{\alpha R^2}{2\lambda ^2(1+\alpha )}}{\displaystyle \frac{l_Q}{R}}\right]\right\}.`$ (7)
Here $`\mathrm{\Lambda }`$ denotes the initial reduced wavelength of the particles, we have also used the fact that $`g_{\mathrm{}}=g/(1+\alpha )`$, where $`g=\frac{GM}{R^2}`$.
Imposing the condition $`\alpha =0`$ enables us to rewrite expression (7) as
$`I_N=\mathrm{cos}\left\{{\displaystyle \frac{gm^2Ll_Q\mathrm{\Lambda }}{\mathrm{}^2}}\left[1{\displaystyle \frac{l_Q}{R}}\right]\right\},`$ (8)
which is the interference term that appears in the case of Newtonian gravity . As a matter of fact, the result of COW does not contain the term that is quadratic in $`l_Q`$, in our result it appears because we have introduced a less restricted approximation, to derive the results of COW we need only a homogeneous Newtonian gravitational field , and expression (3) includes the case of an inhomogeneous gravitational field, i.e., the term $`\frac{1+\alpha }{R^2}l^2`$. In other words, expression (8) is the interference term when we consider, in a Newtonian field, dependence, in the height above the surface of the Earth, up to second order in $`l`$.
We may now calculate the difference between the Newtonian and non–Newtonian cases, and divide this result by the Newtonian value, the outcome reads (approximately)
$`\mathrm{\Delta }={\displaystyle \frac{\alpha R^2}{2\lambda ^2(1+\alpha )}}\left(1+{\displaystyle \frac{l_Q}{R}}\right).`$ (9)
Introducing the expression for the Compton wavelength of our new quantum particle with mass $`m_5`$ ($`\lambda =\frac{\mathrm{}}{cm_5}`$), we find that
$`\mathrm{\Delta }={\displaystyle \frac{\alpha (Rcm_5)^2}{2\mathrm{}^2(1+\alpha )}}\left(1+{\displaystyle \frac{l_Q}{R}}\right).`$ (10)
We may rewrite (7) as follows
$`I=I_N\mathrm{cos}\left\{{\displaystyle \frac{gm^2Ll_Q\mathrm{\Lambda }\alpha R^2}{2\mathrm{}^2\lambda ^2(1+\alpha )}}\right\}\pm \sqrt{1I_N^2}\mathrm{sin}\left\{{\displaystyle \frac{gm^2Ll_Q\mathrm{\Lambda }\alpha R^2}{2\mathrm{}^2\lambda ^2(1+\alpha )}}\right\}.`$ (11)
## 3 Quantum Measurements
As has already been mentioned, in the attempts to solve the quantum measurement problem we may find RPIF . This formalism explains a continuous quantum measurement with the introduction of a restriction on the integration domain of the corresponding path integral. This last condition can also be reformulated in terms of a weight functional that has to be considered in the path integral.
Let us explain this point a little bit better, and suppose that we have a particle which shows one–dimensional movement. The amplitude $`A(q^{\prime \prime },q^{})`$ for this particle to move from the point $`q^{}`$ to the point $`q^{\prime \prime }`$ is called propagator. It is given by Feynman
$$A(q^{\prime \prime },q^{})=d\left[q\right]\mathrm{exp}\left(\frac{i}{\mathrm{}}S[q]\right),$$
(12)
here we must integrate over all the possible trajectories $`q(t)`$, $`S[q]`$ is the action of the system, which is defined as
$$S[q]=_t^{}^{t^{\prime \prime }}𝑑tL(q,\dot{q}).$$
(13)
Let us now suppose that we continuously measure the position of this particle, such that we obtain as measurement ouput a certain function $`a(t)`$. In other words, the measuring process gives the value $`a(t)`$ for the coordinate $`q(t)`$ at each time $`t`$, and this output has associated a certain error $`\mathrm{\Delta }a`$, which is determined by the experimental resolution of the measuring device. The amplitude $`A_{[a]}(q^{\prime \prime },q^{})`$ can be now thought of as a probability amplitude for the continuous measuring process to give the result $`a(t)`$. Taking the square modulus of this amplitude allows us to find the probability density for different measurement outputs.
Clearly, the integration domain in the Feynman path–integral should be restricted to those trajectories that match with the experimental output. RPIF says that this condition can be introduced by means of a weight functional $`\omega _a[q]`$ . This means that expression (12) becomes now
$$A_a=d\left[q\right]\omega _a[q]\mathrm{exp}\left(\frac{i}{\mathrm{}}S[q]\right).$$
(14)
The more probable the trajectory $`[q]`$ is, according to the output $`a`$, the bigger that $`\omega _a[q]`$ becomes . This means that the value of $`\omega _a[q]`$ is approximately one for all trajectories $`[q]`$ that agree with the measurement output $`a`$, and it is almost 0 for those that do not match with the result of the experiment. Clearly, the weight functional contains all the information about the interaction between measuring device and measured system.
Let us now consider the propagator of a particle whose Lagrangian is given by (3), (the particle goes from point $`P`$ to point $`Q`$)
$`U(Q,\tau ^{\prime \prime };P,\tau ^{})=\left({\displaystyle \frac{m}{2\pi i\mathrm{}T}}\right)\mathrm{exp}\left\{{\displaystyle \frac{im}{2\mathrm{}T}}L^2\right\}{\displaystyle }d[l(t)]\mathrm{exp}({\displaystyle \frac{i}{\mathrm{}}}{\displaystyle _\tau ^{}^{\tau ^{\prime \prime }}}[{\displaystyle \frac{m}{2}}\dot{l}^2`$
$`+G_{\mathrm{}}{\displaystyle \frac{mM}{R}}([1+\alpha +{\displaystyle \frac{\alpha R}{\lambda }}({\displaystyle \frac{R}{2\lambda }}1)][{\displaystyle \frac{1+\alpha }{R}}{\displaystyle \frac{\alpha R}{2\lambda ^2}}]l+{\displaystyle \frac{1+\alpha }{R^2}}l^2)]dt).`$ (15)
We now introduce a measuring process, namely we will monitor continuously the l–coordinate of the particle. Then expression (15) becomes now
$`U_{[a(t)]}(Q,\tau ^{\prime \prime };P,\tau ^{})=\left({\displaystyle \frac{m}{2\pi i\mathrm{}T}}\right)\mathrm{exp}\left\{{\displaystyle \frac{im}{2\mathrm{}T}}L^2\right\}`$
$`\times {\displaystyle }d[l(t)]w_{[a(t)]}[l(t)]\mathrm{exp}({\displaystyle \frac{i}{\mathrm{}}}{\displaystyle _\tau ^{}^{\tau ^{\prime \prime }}}[{\displaystyle \frac{m}{2}}\dot{l}^2`$
$`+G_{\mathrm{}}{\displaystyle \frac{mM}{R}}([1+\alpha +{\displaystyle \frac{\alpha R}{\lambda }}({\displaystyle \frac{R}{2\lambda }}1)][{\displaystyle \frac{1+\alpha }{R}}{\displaystyle \frac{\alpha R}{2\lambda ^2}}]l+{\displaystyle \frac{1+\alpha }{R^2}}l^2)]dt).`$ (16)
The modulus square of this last expression gives the probability of obtaining as measurement output (for the $`l`$–coordinate) function $`a(t)`$. The weight functional $`w_{[a(t)]}[l(t)]`$ contains the information concerning the measurement, and is determined by the experimental construction .
At this point, in order to obtain theoretical predictions, we must choose a particular expression for $`w_{[a(t)]}[l(t)]`$. We know that the results coming from a Heaveside weight functional and those coming from a gaussian one coincide up to the order of magnitude. These last remarks allow us to consider a gaussian weight functional as an approximation of the correct expression.
It will be supposed that the weight functional of our measuring device has precisely this gaussian form. We may wonder if this is not an unphysical assumption, and in favor of this argument we may comment that recently it has been proved that there are measuring apparatuses which show this kind of behaviour .
Therefore we may now choose as our weight functional the following expression
$$\omega _{[a(t)]}[l(t)]=\mathrm{exp}\left\{\frac{2}{T\mathrm{\Delta }a^2}_\tau ^{}^{\tau ^{\prime \prime }}[l(t)a(t)]^2𝑑t\right\},$$
(17)
here $`\mathrm{\Delta }a`$ represents the error in our measurement.
Hence with the introduction of a continuous quantum measurement the new propagator is
$`U_{[a(t)]}(Q,\tau ^{\prime \prime };P,\tau ^{})=\left({\displaystyle \frac{m}{2\pi i\mathrm{}T}}\right)\mathrm{exp}\left\{{\displaystyle \frac{im}{2\mathrm{}T}}L^2\right\}`$
$`{\displaystyle }d[l(t)]\mathrm{exp}\{{\displaystyle \frac{2}{T\mathrm{\Delta }a^2}}{\displaystyle _\tau ^{}^{\tau ^{\prime \prime }}}[l(t)a(t)]^2dt\}\mathrm{exp}({\displaystyle \frac{i}{\mathrm{}}}{\displaystyle _\tau ^{}^{\tau ^{\prime \prime }}}[{\displaystyle \frac{m}{2}}\dot{l}^2`$
$`+G_{\mathrm{}}{\displaystyle \frac{mM}{R}}\{[1+\alpha +{\displaystyle \frac{\alpha R}{\lambda }}({\displaystyle \frac{R}{2\lambda }}1)][{\displaystyle \frac{1+\alpha }{R}}{\displaystyle \frac{\alpha R}{2\lambda ^2}}]l+{\displaystyle \frac{1+\alpha }{R^2}}l^2\}]dt).`$ (18)
It can be rewritten as follows
$`U_{[a(t)]}(Q,\tau ^{\prime \prime };P,\tau ^{})=\left({\displaystyle \frac{m}{2\pi i\mathrm{}T}}\right)\mathrm{exp}\left\{{\displaystyle \frac{im}{2\mathrm{}T}}L^2\right\}`$
$`\times \mathrm{exp}\left\{{\displaystyle \frac{i}{\mathrm{}}}g_{\mathrm{}}mR\left[1+\alpha +{\displaystyle \frac{\alpha R}{\lambda }}({\displaystyle \frac{R}{2\lambda }}1)\right]T\right\}\mathrm{exp}\left\{{\displaystyle \frac{2}{T\mathrm{\Delta }a^2}}{\displaystyle _\tau ^{}^{\tau ^{\prime \prime }}}a^2(t)𝑑t\right\}`$
$`\times {\displaystyle }d[l(t)]\mathrm{exp}\left\{{\displaystyle \frac{i}{\mathrm{}}}{\displaystyle _\tau ^{}^{\tau ^{\prime \prime }}}[{\displaystyle \frac{m}{2}}\dot{l}^2+F(t)l{\displaystyle \frac{m}{2}}\omega ^2l^2]dt\right\}.`$ (19)
In this last expression we have introduced the following definitions, namely $`F(t)=mg[\frac{\alpha R^2}{2\lambda ^2(1+\alpha )}1]\frac{4i\mathrm{}}{T\mathrm{\Delta }a^2}a(t)`$, and $`\omega =i\mathrm{\Omega }`$, where $`\mathrm{\Omega }^2=2\frac{g}{R}(1+\frac{2i\mathrm{}R}{mgT\mathrm{\Delta }a^2})`$. It is readily seen that we have now the propagator of a driven harmonic oscillator, but now frequency and driving force have nonvanishing imaginary parts, which emerge as a direct consequence of our measuring process. We already know how to evaluate this kind of path integrals .
The result of the integration yields (here we do not assume that $`l_P=0`$)
$`U_{[a(t)]}(Q,\tau ^{\prime \prime };P,\tau ^{})=\left({\displaystyle \frac{m}{2\pi i\mathrm{}T}}\right)^{\frac{3}{2}}\sqrt{{\displaystyle \frac{\mathrm{\Omega }T}{\mathrm{sinh}(\mathrm{\Omega }T)}}}\mathrm{exp}\left\{{\displaystyle \frac{im}{2\mathrm{}T}}L^2\right\}`$
$`\times \mathrm{exp}\left\{{\displaystyle \frac{i}{\mathrm{}}}g_{\mathrm{}}mR[1+\alpha +{\displaystyle \frac{\alpha R}{\lambda }}({\displaystyle \frac{R}{2\lambda }}1)]T\right\}\mathrm{exp}\left\{{\displaystyle \frac{2}{T\mathrm{\Delta }a^2}}{\displaystyle _\tau ^{}^{\tau ^{\prime \prime }}}a^2(t)𝑑t\right\}`$
$`\times \mathrm{exp}({\displaystyle \frac{im\mathrm{\Omega }}{2\mathrm{}\mathrm{sinh}(\mathrm{\Omega }T)}}[(l_Q^2+l_P^2)\mathrm{cosh}(\mathrm{\Omega }T)2l_Ql_P`$
$`{\displaystyle \frac{8i\mathrm{}}{mT\mathrm{\Omega }\mathrm{\Delta }a^2}}\left\{l_QF^{(1)}(\tau ^{\prime \prime },\tau ^{})+l_PF^{(2)}(\tau ^{\prime \prime },\tau ^{})\right\}`$
$`+2g{\displaystyle \frac{l_Q+l_P}{\mathrm{\Omega }^2}}\left\{{\displaystyle \frac{\alpha R^2}{2\lambda ^2(1+\alpha )}}1\right\}\left\{\mathrm{cosh}(\mathrm{\Omega }T)1\right\}`$
$`2{\displaystyle \frac{g^2}{\mathrm{\Omega }^2}}\left\{{\displaystyle \frac{\alpha R^2}{2\lambda ^2(1+\alpha )}}1\right\}^2\left\{{\displaystyle \frac{1\mathrm{cosh}(\mathrm{\Omega }T)}{\mathrm{\Omega }^2}}+{\displaystyle \frac{T\mathrm{sinh}(\mathrm{\Omega }T)}{2\mathrm{\Omega }}}\right\}`$
$`+{\displaystyle \frac{8i\mathrm{}g}{m\mathrm{\Omega }^2T\mathrm{\Delta }a^2}}\left\{{\displaystyle \frac{\alpha R^2}{2\lambda ^2(1+\alpha )}}1\right\}{\displaystyle _\tau ^{}^{\tau ^{\prime \prime }}}F^{(1)}(\tau ,\tau ^{})\mathrm{sinh}(\mathrm{\Omega }(\tau ^{\prime \prime }\tau ))𝑑\tau `$
$`+{\displaystyle \frac{8i\mathrm{}g}{mT\mathrm{\Delta }a^2\mathrm{\Omega }^3}}\left\{{\displaystyle \frac{\alpha R^2}{2\lambda ^2(1+\alpha )}}1\right\}\left\{F^{(3)}(\tau ^{\prime \prime },\tau ^{})F^{(2)}(\tau ^{\prime \prime },\tau ^{})\right\}`$
$`+{\displaystyle \frac{32\mathrm{}^2}{m^2T^2\mathrm{\Omega }^2\mathrm{\Delta }a^4}}F^{(4)}(\tau ^{\prime \prime },\tau ^{})]),`$ (20)
where $`F^{(1)}(\tau ^{\prime \prime },\tau ^{})=_\tau ^{}^{\tau ^{\prime \prime }}a(\tau )\mathrm{sinh}(\mathrm{\Omega }(\tau \tau ^{}))𝑑\tau `$, we also have defined, $`F^{(2)}(\tau ^{\prime \prime },\tau ^{})=_\tau ^{}^{\tau ^{\prime \prime }}a(\tau )\mathrm{sinh}(\mathrm{\Omega }(\tau ^{\prime \prime }\tau ))𝑑\tau `$, $`F^{(3)}(\tau ^{\prime \prime },\tau ^{})=_\tau ^{}^{\tau ^{\prime \prime }}a(\tau )\mathrm{sinh}(\mathrm{\Omega }(\tau ^{\prime \prime }\tau ))\mathrm{cosh}(\mathrm{\Omega }(\tau \tau ^{}))𝑑\tau `$, and finally $`F^{(4)}(\tau ^{\prime \prime },\tau ^{})=_\tau ^{}^{\tau ^{\prime \prime }}𝑑\tau _\tau ^{}^\tau 𝑑sa(\tau )a(s)\mathrm{sinh}(\mathrm{\Omega }(\tau ^{\prime \prime }\tau ))\mathrm{sinh}(\mathrm{\Omega }(s\tau ^{}))`$.
## 4 Discussion
Expression (11) shows clearly that the interference pattern emerging in the case of a non–Newtonian gravity theory, here the deviation comprises a Yukawa term, does not match with the results of the inverse–square law situation. If we consider the experimental parameters of COW, and also the values $`\alpha 10^3`$ and $`\lambda 10^4`$cm (here $`\lambda `$ denotes the range of the Yukawa interaction), then we deduce that $`\frac{gm^2Ll_Q\mathrm{\Lambda }\alpha R^2}{2\mathrm{}^2\lambda ^2(1+\alpha )}10^7(cm)^1l_Q`$. Hence (11) reduces to
$`I=I_N\mathrm{cos}\left\{10^7(cm)^1l_Q\right\}\pm \sqrt{1I_N^2}\mathrm{sin}\left\{10^7(cm)^1l_Q\right\}.`$ (21)
The dependence in $`l_Q`$ of (21) could, in principle, be detected.
$`{\displaystyle \frac{\mathrm{\Delta }I}{\mathrm{\Delta }l_Q}}10^7(cm)^1\left(I_N\mathrm{sin}\left\{10^7(cm)^1l_Q\right\}\pm \sqrt{1I_N^2}\mathrm{cos}\left\{10^7(cm)^1l_Q\right\}\right)`$
$`+{\displaystyle \frac{\mathrm{\Delta }I_n}{\mathrm{\Delta }l_Q}}\left(\mathrm{cos}\left\{10^7(cm)^1l_Q\right\}{\displaystyle \frac{I_N}{\sqrt{1I_N^2}}}\mathrm{sin}\left\{10^7(cm)^1l_Q\right\}\right).`$ (22)
Knowing that these Yukawa terms emerge in a natural manner in some unified field theories (they are related to the existence of an intermediate–range new force coupled to baryon number or hypercharge), then our result could help to determine the phenomenological parameters $`\alpha `$ and $`\lambda `$. Parameter $`\alpha `$ could also be composition–dependent , this possibility could also be tested employing expression (11), i.e., performing the interference experiment with different type of materials. The relevance of this last point is related to the analysis of the validity of WEP and of the Strong Equivalence Principle (SEP), namely a composition–independent $`\alpha `$ would not violate WEP, but it might violate SEP .
Let us now consider the case in which our two beams start at point $`P`$ and afterwards are detected at point $`Q`$, here we do not consider any restriction on the wavelength of the beams. We also assume that the $`l`$–coordinate of the beams is being continuously measured, in other words, we may use expression (20) to calculate the corresponding wave function. The measuring process will take place under the following restrictions: (i) two functions are obtained as measurement outputs, namely $`a(t)`$ and $`b(t)`$ (each beam has its own function); (ii) we carry out this experiment using two devices (each beam has its own measuring device), whose errors are not the same, i.e., $`\mathrm{\Delta }a\mathrm{\Delta }b`$. We may then calculate the emerging interference pattern, the result is given by the real part of the following expression
$`I=\mathrm{exp}({\displaystyle \frac{im}{2\mathrm{}}}[(l_Q^2+l_P^2)\{{\displaystyle \frac{\mathrm{\Omega }}{\mathrm{tanh}(\mathrm{\Omega }T)}}{\displaystyle \frac{\mathrm{\Gamma }}{\mathrm{tanh}(\mathrm{\Gamma }T)}}\}`$
$`+2l_Ql_P\{{\displaystyle \frac{\mathrm{\Gamma }}{\mathrm{sinh}(\mathrm{\Gamma }T)}}{\displaystyle \frac{\mathrm{\Omega }}{\mathrm{sinh}(\mathrm{\Omega }T)}}\}+{\displaystyle \frac{8i\mathrm{}}{mT}}({\displaystyle \frac{l_QF^{(1)}(\tau ^{\prime \prime },\tau ^{})+l_PF^{(2)}(\tau ^{\prime \prime },\tau ^{})}{\mathrm{sinh}(\mathrm{\Omega }T)\mathrm{\Delta }a^2}}`$
$`{\displaystyle \frac{l_Qf^{(1)}(\tau ^{\prime \prime },\tau ^{})+l_Pf^{(2)}(\tau ^{\prime \prime },\tau ^{})}{\mathrm{sinh}(\mathrm{\Gamma }T)\mathrm{\Delta }b^2}})`$
$`+2g\stackrel{~}{\alpha }[\left\{{\displaystyle \frac{\mathrm{cosh}(\mathrm{\Omega }T)1}{\mathrm{\Omega }\mathrm{sinh}(\mathrm{\Omega }T)}}\right\}(l_Q+l_P+{\displaystyle \frac{g}{\mathrm{\Omega }^2}}\stackrel{~}{\alpha }){\displaystyle \frac{T}{2\mathrm{\Omega }^2}}`$
$`\left\{{\displaystyle \frac{\mathrm{cosh}(\mathrm{\Gamma }T)1}{\mathrm{\Gamma }\mathrm{sinh}(\mathrm{\Gamma }T)}}\right\}(l_Q+l_P+{\displaystyle \frac{g}{\mathrm{\Gamma }^2}}\stackrel{~}{\alpha })+{\displaystyle \frac{T}{2\mathrm{\Gamma }^2}}]`$
$`+{\displaystyle \frac{8i\mathrm{}g}{mT}}\stackrel{~}{\alpha }[{\displaystyle \frac{F^{(3)}(\tau ^{\prime \prime },\tau ^{})F^{(2)}(\tau ^{\prime \prime },\tau ^{})}{\mathrm{\Delta }a^2\mathrm{sinh}(\mathrm{\Omega }T)\mathrm{\Omega }^2}}{\displaystyle \frac{f^{(3)}(\tau ^{\prime \prime },\tau ^{})f^{(2)}(\tau ^{\prime \prime },\tau ^{})}{\mathrm{\Delta }b^2\mathrm{sinh}(\mathrm{\Gamma }T)\mathrm{\Omega }^2}}`$
$`+{\displaystyle \frac{F^{(5)}}{\mathrm{\Omega }\mathrm{\Delta }a^2sinh(\mathrm{\Omega }T)}}{\displaystyle \frac{f^{(5)}}{\mathrm{\Gamma }\mathrm{\Delta }b^2sinh(\mathrm{\Gamma }T)}}]`$
$`+{\displaystyle \frac{32\mathrm{}^2}{m^2T^2}}\{{\displaystyle \frac{F^{(4)}(\tau ^{\prime \prime },\tau ^{})}{\mathrm{\Omega }\mathrm{sinh}(\mathrm{\Omega }T)\mathrm{\Delta }a^4}}{\displaystyle \frac{f^{(4)}(\tau ^{\prime \prime },\tau ^{})}{\mathrm{\Gamma }\mathrm{sinh}(\mathrm{\Gamma }T)\mathrm{\Delta }b^4}}\}]).`$ (23)
Here $`\mathrm{\Gamma }^2=2\frac{g}{R}(1\frac{2i\mathrm{}R}{mgT\mathrm{\Delta }b^2})`$. Additionally, the following functions have been employed $`f^{(1)}(\tau ^{\prime \prime },\tau ^{})=_\tau ^{}^{\tau ^{\prime \prime }}b(\tau )\mathrm{sinh}(\mathrm{\Gamma }(\tau \tau ^{}))𝑑\tau `$$`f^{(2)}(\tau ^{\prime \prime },\tau ^{})=_\tau ^{}^{\tau ^{\prime \prime }}b(\tau )\mathrm{sinh}(\mathrm{\Gamma }(\tau ^{\prime \prime }\tau ))𝑑\tau `$, $`f^{(3)}(\tau ^{\prime \prime },\tau ^{})=_\tau ^{}^{\tau ^{\prime \prime }}b(\tau )\mathrm{sinh}(\mathrm{\Gamma }(\tau ^{\prime \prime }\tau ))\mathrm{cosh}(\mathrm{\Gamma }(\tau \tau ^{}))𝑑\tau `$. Also the following parameters have been defined $`f^{(4)}(\tau ^{\prime \prime },\tau ^{})=_\tau ^{}^{\tau ^{\prime \prime }}𝑑\tau _\tau ^{}^\tau 𝑑sb(\tau )b(s)\mathrm{sinh}(\mathrm{\Gamma }(\tau ^{\prime \prime }\tau ))\mathrm{sinh}(\mathrm{\Gamma }(s\tau ^{}))`$, $`F^{(5)}=_\tau ^{}^{\tau ^{\prime \prime }}F^{(1)}(\tau ,\tau ^{})\mathrm{sinh}(\mathrm{\Omega }(\tau ^{\prime \prime }\tau ))𝑑\tau `$, $`f^{(5)}=_\tau ^{}^{\tau ^{\prime \prime }}f^{(1)}(\tau ,\tau ^{})\mathrm{sinh}(\mathrm{\Gamma }(\tau ^{\prime \prime }\tau ))𝑑\tau `$, and finally $`\stackrel{~}{\alpha }=\frac{\alpha R^2}{2\lambda ^2(1+\alpha )}1`$.
It is readily seen that the emerging interference pattern depends not only on the parameters that define the Yukawa interaction, $`\alpha `$, $`\lambda `$, but also on the parameters that appear in RPIF, namely $`a(t)`$, $`b(t)`$, $`\mathrm{\Delta }a`$, and also on $`\mathrm{\Delta }b`$. Hence we could compare with the Newtonian situation , and therefore obtain an additional scheme that could render some restrictions upon the possible values of $`\alpha `$ and $`\lambda `$. If we consider the limit $`\alpha 0`$ and $`\lambda \mathrm{}`$, we obtain the results of the Newtonian situation .
Taking a look at expression (23) we may notice that that the mass of the test particle appears, explicitly, in the expression for the interference pattern, as happens also in the context of Newtonian gravity , and always as a function of $`\mathrm{}/m`$.
Expression (23), at the same time, also gives a new testing framework for the theoretical predictions of RPIF, which makes the work in this direction a little bit more complete .
This procedure could also be applied to some other possible modifications of the Newtonian gravity law, for instance, we could consider the case of the extra potential $`V(r)=\alpha /r^5`$ (this kind of terms arise in some modified gravitational theories in which the non–Newtonian behavior stems from antisymmetric terms in the metric tensor ), and then carry out the same analysis, but now using as Lagrangian
$`L={\displaystyle \frac{m}{2}}\dot{\stackrel{}{r}}^2+{\displaystyle \frac{1}{R}}\left(GmM+{\displaystyle \frac{\alpha }{R^4}}\right)\left({\displaystyle \frac{GmM}{R}}+5{\displaystyle \frac{\alpha }{R^5}}\right){\displaystyle \frac{l}{R}}+\left({\displaystyle \frac{GmM}{R}}+15{\displaystyle \frac{\alpha }{R^5}}\right){\displaystyle \frac{l^2}{R^2}}.`$ (24)
At this point the feasibility of the present proposal must be addressed. A monocromatic beam of particles could be used (as in COW case ), but here no restriction on the size of the corresponding wave packets is needed, in COW this condition emerges because a WKB approach describes the physical situation. Afterwards, the beam could be split in two parts, and then the vertical coordinate (in the present work it has been denoted by $`l`$) of each one of these beams has to be continuously monitored, and finally, both beams have to be brought together, and the corresponding interference pattern measured. In contrast to most of the existing experiments , the present idea could allow us to detect the effects of a Yukawa term upon quantum systems. Nevertheless, experimentally, the continuous monitoring of a moving particle lies, currently, outside the present technological posibilities . Neverwithstanding, the advances that in the topic of trapped ions have been achieved , allows us to consider the possibility of carrying out the needed experiments in a, hopefully, near future.
As was mentioned before, the “geophysical window” has not been tested severely , i.e., over the 10m–1000m distance scale. The lower limit of this range could be used in the present kind of proposals. Of course, regions smaller than 10m could also be explored, for instance, the 1cm–100cm distance scale. If we consider $`\lambda 10`$m, then the current experimental limit reads $`\alpha 10^1`$ . Looking at expression (23) we may notice that a crucial point in this kind of proposals concerns the resolution of the measuring devices, i.e., $`\mathrm{\Delta }a`$ and $`\mathrm{\Delta }b`$. Hence one of the points that determines the feasibility of the present idea is related to a, enoughly, small experimental error. The situation in which $`\mathrm{\Delta }a2\mu `$m (which is the resolution in the case of a particle in a Paul trap , and therefore at least in the context of motionless situations has already been achieved) and $`\mathrm{\Delta }b=10^3\times \mathrm{\Delta }a`$ could be an interesting case to consider.
Finally, a word must be said about some schemes which also imply a violation of the equivalence principle . In these models the contradictions (between the predictions of general relativity and the results in a Minkowskian spacetime) emerge as a purely quantum mechanical effect, i.e., they appear when a quantum system, the one has no classical counterpart, is embeded in a curved spacetime . This should be no surprise, indeed, even the kinematical description of quantum mechanical systems (without classical analogue) moving in a classical gravitational field shows conceptual difficulties , for instance, there is no consistent definition for the concept of time of flight probability distribution. This simple example also shows, very clearly, the danger of extrapolating, to quantum systems without classical analogue, the concepts of classical physics. Of course, these last arguments do not imply that the analysis of the quantum effects in a classical gravitational field shall not be carried out, they only assert (as it has already been pointed out ) that one should be very careful when addressing this issue.
In the present work, the possible violations of the equivalence principle are due to the presence of an additional interaction (the Yukawa term), a factor that is absent in the aforementioned models .
An interesting question at this point is the following one: let us suppose that we have performed the here proposed experiment, and that it implies a violation of the equivalence principle, how could we determine if this violation stems from the existence of a Yukawa term or it is quantum induced?, here the phrase “quantum induced” means the effects of a classical gravitational field upon a quantum mechanical system without classical analogue. The answer might come from the role that mass plays in the measurement outputs. Indeed, if we take a look at the dependence upon mass of the interference pattern, expression (23), it can be readily seen that it is not the same dependence that appears in some situations in which this violation is quantum induced (see expression (16) of and also of and hence it could be possible to determine the origin of this violation, i.e., it would suffice to perform the experiment several times, using each time a different mass.
Though the COW experiment has been performed with a very good precision , currently there are some works which endow COW with some discrepancies . Concerning these new effects, possibly the mixture of a quantum measurement process and gravitational effects renders an unavoidable modification to the de Broglie’s wave–particle duality . Of course, more work is needed in this direction, where not only the case of quantum demolition measurements (as position monitoring) have to been considered, but also the possibilities that quantum nondemolition measurements could offer in this issue have to be analyzed.
Acknowledgments
The author would like to thank A. A. Cuevas–Sosa and A. Camacho–Galván for their help, and D.-E. Liebscher for the fruitful discussions on the subject. The hospitality of the Astrophysikalisches Institut Potsdam is also kindly acknowledged. It is also a pleasure to thank R. Onofrio for bringing Ref. to my attention. The author is also indebeted to D. V. Ahluwalia for providing a copy of Ref. before its publication, and for the personal communication of reference . This work was supported by CONACYT (México) Posdoctoral Grant No. 983023.
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# Studies of the high luminosity quasar, PDS 456
## 1 Introduction
PDS 456 is a bright, radio-quiet QSO (V=14) recently discovered by Torres et al. (1997), in a search for young stellar objects. It lies fairly close to the Galactic plane ($`\beta =12`$) and is seen through an extinction of A$`{}_{V}{}^{}=1.5`$. PDS 456 is at a similar redshift (z=0.184) to 3C 273, but has a higher bolometric luminosity (by a factor of $`\times `$1.7, Simpson et al. 1999). Overall, PDS 456 is the most luminous object in the local (z $`<0.3`$) Universe (with M$`{}_{V}{}^{}=27`$, L$`{}_{BOL}{}^{}10^{47}`$ erg s<sup>-1</sup>).
## 2 Multi-Wavelength Properties of PDS 456
We have conducted an extensive campaign to observe PDS 456 from the radio through to the hard X-ray band. The spectral energy distribution (SED) is shown in Figure 1a and the optical spectrum in Figure 1b. PDS 456 shows strong H I emission lines, and like many NLS1s has strong optical FeII emission but weak OIII. Although by definition a broad-line quasar, PDS 456 has only a moderate width in H$`\beta `$ (FWHM = 3000 km s<sup>-1</sup>). VLA observations confirm that, unlike 3C 273, PDS 456 is radio-quiet (R<sub>L</sub>=-0.7) and has little extended radio emission. Overall the SED is dominated by the optical/UV ‘big blue bump’. The bolometric luminosity of PDS 456 is of the order 10<sup>47</sup> erg s<sup>-1</sup> (assuming $`H_0=50`$ km s<sup>-1</sup> Mpc<sup>-1</sup> and $`q_0=0.5`$).
## 3 X-ray Observations of PDS 456
### 3.1 The X-ray Spectrum
As part of our campaign, we observed PDS 456 with ASCA on 7-8 March 1998 and with RXTE on 7-10 March 1998. The hard X-ray spectrum of PDS 456 obtained by both ASCA and RXTE shows complex features (see Reeves et al. 2000 for details). The data/model residuals from a power-law fit ($`\mathrm{\Gamma }=2.4`$) to the RXTE data are shown in Figure 2a. Unsurprisingly a power-law gave an inadequate fit to the data in this band. We find that an unusually deep and ionised Fe K edge is observed, with best-fit parameters of E=8.7$`\pm `$0.2 keV and $`\tau =0.75\pm 0.15`$. The edge is detected in both the ASCA and RXTE data to $`>`$99.99% confidence. There is also some evidence in the X-ray spectrum for a broadened ($`\sigma 1`$ keV) line at 6 keV; this line may originate from the inner disk, as hypothesised in Seyfert 1s (Tanaka et al. 1995). A ‘warm’ absorber of lower ionisation may also be present at soft X-ray energies, whose properties are also similar to those observed in Seyfert 1s (e.g. Reynolds 1997).
A model consisting of reflection off a highly ionised accretion disk provides the best-fit to the hard X-ray spectrum; with disk solid angle, R=$`\mathrm{\Omega }/2\pi =1.0`$, ionisation parameter, $`\xi =6000`$ erg cm s<sup>-1</sup> and T$`{}_{disk}{}^{}=10^6`$ K. The high-ionisation of the disk reflection component can reproduce both the depth of the edge and its energy at 8.7 keV. Such high-ionisation reflection features are predicted in disk photoionisation models (e.g. Ross, Fabian & Young 1999, Nayakshin et al. 1999), particularly at high accretion rates when the primary X-ray emission is steep. Therefore the high ionisation of the reflector could imply a high accretion rate in PDS 456 (relative to the Eddington limit), particularly as $`\xi \dot{m}^3`$ for a photoionised accretion disk. This interpretation is consistent with the other X-ray properties of PDS 456, namely a steep underlying X-ray continuum and rapid X-ray variability, both of which are commonplace in NLS1s (Boller et al. 1996). NLS1s are also thought to be accreting near the Eddington limit (e.g. Pounds et al. 1995); indeed recent evidence has been found in one NLS1 (Ark 564) for a spectrum consistent with ionised disk reflection (Vaughan et al. 1999).
### 3.2 X-ray Variability
Both the ASCA and RXTE data were examined to search for X-ray variability. A strong hard X-ray flare is observed in the RXTE observation (figure 2b), well above any residual fluctuations in the detector background; the doubling time for the flare was $`15`$ ksec. (Note that, unfortunately, the shorter ASCA observation had ended by the time of the flare.) We also calculated that the probability of finding another contaminating source of comparable brightness in the RXTE beam was low ($`<2`$%), although not totally excluded. Additionally there is no other X-ray source detected in the RASS (ROSAT All Sky Survey) to within a degree of PDS 456.
Therefore, if confirmed, this would be unprecedented behaviour in such a high-luminosity source. This suggests, from simple light-crossing arguments, a maximum size of $`l=4.5\times 10^{12}`$ m for the varying region. For a black hole of mass $`10^9`$M (corresponding to PDS 456, with L$`{}_{BOL}{}^{}=10^{47}`$ erg s<sup>-1</sup>, at the Eddington limit), this implies that the X-ray flare occurs within a region of less than 2 Schwarzschild radii (2R<sub>S</sub>). A smaller mass black hole would loosen this requirement somewhat, but would then imply a super-Eddington accretion rate. Therefore one possible implication of the rapid variability is accretion near to or greater than L<sub>Edd</sub>. The variability also implies a (non-beamed) efficiency of converting matter to energy of $``$ 5%, close to the limit for a Schwarzschild black hole (see Fabian 1979).
## 4 Conclusions
In conclusion, PDS 456 is a remarkable object, showing clear features of a high ionisation reprocessor, one possible interpretation of which is through reflection off a highly ionised accretion disk. Overall the high ionisation spectral features, steep X-ray emission and the extreme rapid variability suggest that the super-massive black hole in PDS 456 could be running at an unusually high accretion rate.
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# The chemical evolution of gas-rich dwarf galaxies
## 1 Introduction
The chemical evolution of dwarf irregular (dIrr) <sup>1</sup><sup>1</sup>1In this paper a dwarf galaxy is defined to have absolute magnitude M$`{}_{\mathrm{B}}{}^{}17`$. and blue compact emission-line galaxies (BCGs) is of particular interest because a substantial body of observational data is available and some degree of simplicity exists because of the low level of ‘metal’ enrichment and absence of large abundance gradients. Furthermore, their wide range of intrinsic properties makes them suitable objects for testing certain expectations from stellar nucleosynthesis theory and the ‘Simple’ or other models of galactic chemical evolution, although at the same time there are complications associated with inflow of unprocessed material, outflow in homogeneous or selective galactic winds and bursting (or ‘gasping’) modes of star formation. Chemical evolution models attempt to apply all these concepts to account for the distribution of different elements, notably helium, oxygen and nitrogen, in relation to star formation rates and gas fractions. Because many parameters such as these last two are generally very poorly determined, the most convincing tests come from the comparison of different elements with one another.
Back in the 1970s, Smith , Peimbert and Edmunds & Pagel noticed a contribution of primary nitrogen to the N/O ratio in Galactic and extragalactic H II regions with low oxygen abundance and Edmunds & Pagel attributed the existence of scatter in N/O at a given O/H to the existence of a time delay in primary nitrogen production by intermediate-mass stars, combined with differing effective ages of the underlying stellar populations. Alloin et al. also noted the primary nitrogen and attributed scatter in N/O to variations in the initial mass function (IMF), whereas Lequeux et al. in their classic study of helium, nitrogen and oxygen in irregular galaxies and BCGs confirmed the primary nitrogen likewise, but were not convinced that there was any real scatter in their data. The models of Alloin et al. and Lequeux et al. assumed evolution to take place smoothly as a function of time; Matteucci & Chiosi were the first to incorporate into chemical evolution models for these systems the idea of bursting modes of star formation as prevously inferred by Searle & Sargent and Searle, Sargent & Bagnuolo , and interpreted on the basis of the SSPSF hypothesis by Gerola, Seiden & Schulman . The basic pattern of a primary (constant N/O) pattern at low metallicities in H II regions changing over to a secondary pattern (N/O $``$ O/H) at higher ones has been confirmed in many more recent investigations (e.g. Vila-Costas & Edmunds 1993; van Zee, Salzer & Haynes 1998).
The chemical evolution of dIrrs and BCGs has been studied in many more recent investigations. Matteucci & Tosi found good fits to the data with a Salpeter IMF, inflow, homogeneous outflow, bursting star formation and a choice of third dredge-up parameters from Renzini & Voli , attributing scatter in the N/O ratio to variations in $`M_{\mathrm{up}}`$, the upper limit to the masses of stars undergoing the third dredge-up with hot-bottom burning. Garnett indicated schematically how the occurrence of bursts could in itself lead to variations in the N/O ratio just as a result of observing systems at different stages in the burst cycle. Pilyugin developed similar ideas in quantitative numerical models involving self-enrichment of H II regions and selective galactic winds as well as bursting modes of star formation, and Marconi, Matteucci & Tosi also developed models with bursts and selective winds, while Carigi, Colín & Peimbert have investigated similar models, but prefer a ‘bottom-heavy’ IMF similar to one claimed in some globular clusters and giving rise to low true yields. However, part of the motivation for invoking selective winds was the apparent existence of a large $`dY/dZ`$ ratio suggested by Pagel et al. , which no longer seems valid , and the scatter in N/O also seems to have been overestimated in those investigations. In this paper, therefore, we investigate the problem again, making use of more recent data and models of stellar nucleosynthesis and exploring in particular the role of mixing processes and of differing burst phases in leading to scatter in the N/O, O/H relation and the relationship between oxygen abundance and gas fraction.
The structure of the article is as follows: Section 2 presents the adopted sample of abundance observations. Section 3 discusses the evolution of H ii-regions, wind-driven bubbles and supernova-driven supershells to investigate possible mixing scenarios. This leads to the description of our double-bursting models in section 4. The results of fitting the models to the sample of observations are presented in section 5, and discussed in section 6.
## 2 The observational data
It has been stressed that the sample should consist mainly of BCGs, as these objects are extreme in both star formation rates and metallicity, thus being attractive for a bursting model. Another demand has been to include the most recent data only, in order to get the observations as accurate as possible. Further, a few damped Ly-alpha systems (DLA) have been included, mainly for the sake of comparison. The motivation for including them is found in the growing opinion that they represent dwarf galaxies in early (high z) stages of evolution. Thus, they may offer some insight into the extreme low-metallicity environment in the early evolution of present-day dwarf galaxies.
The following sample is selected: BCGs and dIrrs from Pagel et al. , BCGs from Izotov, Thuan & Lipovetsky \[1997a\], the BCG SBS 0335-052 from Izotov et al. \[1997b\], the BCG IZw18 from Izotov & Thuan , and finally four DLAs from Lu, Sargent & Barlow . From the sample of Pagel et al. are excluded LMC, NGC5253 and NGC5455/NGC5461 (the last two are H ii-regions in the nearby spiral M101), and objects included from one of the other sources.
The total sample is shown in fig. 1 for N/O-O/H and fig. 2 for Y-O/H. The N/O abundances seem to have a constant level for 12+log(O/H) less than 8, implying the major part of the nitrogen to be produced as a primary element. For 12+log(O/H) higher than about 8, N/O is increasing as a function of metallicity, suggesting the major part of the nitrogen to be secondary. These trends are wellknown, see e.g. Garnett , Vila-Costas & Edmunds , Pettini, Lipman & Hunstead and van Zee, Salzer & Haynes \[1998a\].
Further, the scatter in N/O seems to be significant and our models are based on this assumption. If we fit a straight line to our selected observations, ignoring the DLAs, we find the standard deviation in N/O to be $`\sigma `$=0.11, which is of the same order of magnitude as the observational uncertainty \[1997a\], but a real scatter is evident in other data from a wider range of sources, see e.g. Kobulnicky & Skillman .
The DLA-systems pose severe problems when observing N and O abundances. Usually the N-lines are occuring on top of underlying absorption, and O-lines are almost always saturated. Because of these difficulties, Lu et al. found it useful to use Si or S instead of O, which makes sense as O/Si and O/S are found to be the same as solar in both the Galactic disc and halo as well as in H ii-regions in nearby galaxies. The question is whether the abundance ratios are equal to solar at the extreme low metallicity of DLAs. This is by now the largest uncertainty in this method. Thus, in fig. 1, the DLAs are represented by their N/Si and Si/H, assuming that these ratios are equivalent to N/O and O/H. Four systems have been selected from the Lu et al. sample namely those in front of QSO’s 0100+1300, 1331+1704, 1946+7658 and 2343+1232. For the rest of their sample, the quality of the observations restricts the abundance determination to be performed as higher or lower limits. Lu et al. did not make any corrections for dust. However, they noted that the depletion in dust is less than 0.4 dex. Though the observations of the four DLAs are included, they should not be taken too seriously and are not given much attention in this paper.
The helium mass fraction shows a linear dependency on metallicity. This linearity has been used for extrapolation back to O/H=0, giving the primordial He abundance. Izotov et al. \[1997a\] used their observations, included in the present sample as well (filled circles in fig. 2), to obtain 0.243$`\pm `$0.003, using a linear fit with slope dY/dZ=1.7$`\pm `$0.9 and assuming Z=20(O/H). Error bars are omitted in the figure to keep the clearness. However, the error bars are included in the figures presenting our results in section 5 and show no observational evidence for a scatter. Omitting SBS 0749+568, represented by the isolated point above 0.26 in fig. 2, but including all other galaxies in the sample, a linear least-square fit gives dY/dZ=2.63$`\pm `$2.21, again using Z=20 (O/H), and assuming a confidence interval of 95 per cent. This gives a primordial helium abundance $`Y_p`$=0.238$`\pm `$ 0.004, which is consistent with Izotov et al. \[1997a\] within uncertainties.
## 3 Ejecta dispersal and mixing
Many chemical evolution models assume a one-zone description with instantaneous mixing. Assigning the term one-zone to a BCG seems to be quite a poor approximation, and assuming the mixing to be instantaneous will always be doubtfull. However, if we choose to reject the assumptions, we are faced with the problem that no complete theory of mixing is available at present. The problem is that the gas dynamics following a starburst are very complex. The energy input into the ISM comes from photoionization, stellar winds from massive stars followed by SN type II. Different scenarios have been proposed for the dispersal and mixing processes. One of the propositions has been ’self-enrichment’, referring to the suggestion that only the H ii-region, surrounding the newly formed stellar cluster, is enriched with heavier elements, in particular oxygen. Thus, the observed abundances, using emission lines, may be considerably higher than they would be if observed in the neutral medium. The idea was originally proposed by Kunth & Sargent , suggesting that the enrichment is confined to take place within the Strömgren-sphere. However, the hypothesis is controversial, and a lot of opposing arguments exist. Use of self-enrichment implies the assumption of almost instantaneous mixing within the H ii-regions. The problem is that the correlated energy input of SNe changes the physical conditions of H ii-regions dramatically. In fact, the density is decreased by a factor $`10^310^6`$ and the temperature is increased by a similar factor behind the SN-driven shock front. The observed lines of single or double ionized N and O cannot arise from such extreme conditions. A future project could be X-ray abundance observations. If the observed emission lines cannot form within the superbubble/supershell (hereinafter a ’superbubble’ is a wind-driven bubble, while a ’supershell’ is a SN-driven shell), they may arise from a region outside the shock front. But according to detailed numerical hydrodynamical models , the ejecta will stay within the superbubble/supershell for a large part of their evolution, not enriching the surrounding medium. To test the viability of this statement, a few calculations are presented below, comparing radii of superbubbles/supershells and radii of H ii-regions. Further, observations give evidence against the self-enrichment hypothesis, see e.g. van Zee et al. \[1998b\], Kobulnicky , Kobulnicky & Skillman and Kobulnicky & Skillman .
### 3.1 H II-region evolution
In the following, the radii of H ii-regions are calculated for comparison with superbubble/supershell radii. The Strömgren-sphere is defined to be the sphere within which all ionizing photons are absorbed. Thus, setting the flux of ionizing photons equal to the number of recombinations, integrated over the entire Strömgren-sphere gives:
$`{\displaystyle \frac{4\pi }{3}}R(t)^3\alpha _Bn_en_pϵ=N(t)`$ (1)
$`ϵ`$ is the filling factor, $`\alpha _B`$ is the recombination rate equal to $`2.59\times 10^{13}\mathrm{cm}^3\mathrm{s}^1`$ (T=10000K) . $`n_p`$ and $`n_e`$ are the number densities of protons and electrons, respectively, and $`N(t)`$ is the flux of ionizing $`Ly_c`$-photons. Assuming that all hydrogen inside of the Strömgren-sphere is fully ionized, $`n_en_p`$. The filling factor is an indicator of the uniformity of matter, having a value of 1 when the matter is completely uniformly distributed. A typical value of the filling factor is of the order 0.01 . The electron density in H ii-regions is typically $`n_e100\mathrm{cm}^3`$ \[1997a\]. When inserting these two values into eq. 1, the corresponding rms density is $`n_e=10\mathrm{cm}^3`$, since $`ϵ=1`$ and $`n_e=10\mathrm{cm}^3`$ is equivalent to the insertion of the observed values. According to Spitzer \[1978, p.251\], the first phase in the existence of an H ii-region is characterized by an almost static Strömgren-sphere. Henceforth, the radius of this initial sphere will be referred to as ’the initial Strömgren radius’, $`R_0`$. During the second phase, the sphere is expanding until the H ii-region vanishes. If the flux of ionizing photons is assumed constant for the moment, Spitzer found the relation
$`R(t)=R_0\left(1+{\displaystyle \frac{7}{4}}{\displaystyle \frac{C_{II}t}{R_0}}\right)^{\frac{4}{7}}`$ (2)
on the basis that the expansion velocity nearly equates the velocity of the associated shock. $`t`$ is the time since the region started to expand, virtually equal to the age of the burst, and $`C_{II}`$ is the sound speed in the H ii-region, equal to $`17\mathrm{k}\mathrm{m}\mathrm{s}^1`$ (T=10000K, $`\gamma =5/3`$) . Eq. 2 is not realistic, though, because the photon flux is decreasing in time. In the following this is taken into account, extending the calculations by Spitzer. Differentiating eq. 1 with respect to time gives
$`{\displaystyle \frac{dN}{dt}}=3\beta R^2\rho _{II}^2{\displaystyle \frac{dR}{dt}}+2\beta R^3\rho _{II}{\displaystyle \frac{d\rho _{II}}{dt}}`$ (3)
where $`\beta =\frac{4\pi }{3}\eta ^2\alpha _Bϵ`$, $`n_e=\eta \rho _{II}`$, so $`\eta `$ gives the relation between the number density of electrons and the mass density. Below it is shown that $`\eta `$ cancels in the calculations, so we will not worry about its value. Finally $`\rho _{II}`$ is the mass density within the H ii-region. Consider a shell of ionized gas. Assuming uniform expansion, i.e.
$`{\displaystyle \frac{1}{r}}{\displaystyle \frac{dr}{dt}}={\displaystyle \frac{v_i}{R}}`$ (4)
where $`r`$ is the comoving radius of the shell and $`v_i`$ is the velocity of the ionized gas just within the ionization front, assumed to have the same radius $`R`$ as the Strömgren-sphere. The mass inside the comoving shell is conserved by definition, so $`r^3\rho _{II}`$ is constant leading to
$`{\displaystyle \frac{1}{r}}{\displaystyle \frac{dr}{dt}}={\displaystyle \frac{1}{3\rho _{II}}}{\displaystyle \frac{d\rho _{II}}{dt}}`$ (5)
Inserting these two equations into eq. 3 gives
$`v_i={\displaystyle \frac{1}{2}}{\displaystyle \frac{dR}{dt}}{\displaystyle \frac{\frac{dN}{dt}}{6\beta R^2\rho _{II}^2}}`$ (6)
and $`\frac{dR}{dt}=V_i`$, the velocity of the ionization front. The physical conditions across the shock, formed ahead of the ionization front may be described by the jump condition
$`\rho _IV_s^2=p_{II}+\rho _{II}u_s^2`$ (7)
$`\rho _I`$ being the mass density of the surrounding neutral medium, $`V_s`$ is the velocity of the shock, $`p_{II}`$ is the pressure within the H ii-region and $`u_s`$ is the inward velocity of matter with respect to the shock. It has been assumed that the density and pressure between the shock and ionization fronts are the same as within the ionization front, and the pressure of the surrounding neutral medium has been neglected. Assume $`V_s=V_i`$, so $`u_s=u_i=V_iv_i`$, where $`u_i`$ is the inward velocity of matter with respect to the ionization front. Inserting this and the equation of state $`p_{II}=\frac{\rho _{II}C_{II}^2}{\gamma }`$ with $`\gamma =5/3`$, $`\rho _{II}=\sqrt{\frac{N}{\beta R^3}}`$ (from eq. 1) and $`\rho _I=\sqrt{\frac{N_0}{\beta R_0^3}}`$ (i.e. the density before expansion starts) into eq. 7 and 6, one finally arrives at the second order equation
$`\left[\left({\displaystyle \frac{NR_0^3}{N_0R^3}}\right)^{1/2}{\displaystyle \frac{1}{4}}\right]\left({\displaystyle \frac{dR}{dt}}\right)^2{\displaystyle \frac{\frac{dN}{dt}}{6\beta R^2\rho _{II}^2}}\left({\displaystyle \frac{dR}{dt}}\right)`$ (8)
$`\left({\displaystyle \frac{\frac{dN}{dt}}{6\beta R^2\rho _{II}^2}}+{\displaystyle \frac{C_{II}^2}{\gamma }}\right)=0`$
It is straightforward to solve this equation for $`\frac{dR}{dt}`$, obtaining a differential equation, which is solved numerically, using timesteps equal to 1 Myr. For the calculations we used $`n_e=10\mathrm{cm}^3`$ ($`\beta =1`$)as the initial density. When inserted into the second order equation instead of $`\rho _{II}`$, the factor $`\eta `$ cancels remembering that it is included in $`\beta `$ as well. For every timestep, the density $`n_e(t)`$ is calculated from eq. 1. The fluxes of ionizing photons are taken from the models by Stasinska & Leitherer , giving the fluxes from 1 Myr to 10 Myr after the burst in intervals of 1 Myr. An instantaneous starburst of mass $`10^6M_{}`$ was assumed, comparable to those used in our model. Further, a Salpeter IMF with the same upper mass limit is used, only demanding a simple scaling from their lower mass $`1M_{}`$ to our lower mass, either $`0.1`$ or $`0.01M_{}`$. If $`m_L=0.1`$ the normalization constant is 0.17, see section 4.1. Integration of the IMF yields a mass fraction of 0.39 above 1 $`M_{}`$, so the adopted fluxes are multiplied by 0.39. Correspondingly, using $`m_L=0.01`$ the normalization constant is 0.07, giving a mass fraction of 0.17 above 1 $`M_{}`$. The decrease of continuum photon flux calculated by Stasinska & Leitherer using their parameter values is shown in figure 3.
$`R_0`$ is found from eq. 1 by inserting the initial flux of photons and the initial density $`n_e=10\mathrm{cm}^3`$. The initial flux, $`N_0=N(t=0)`$ is the flux just before the H ii-region starts to expand. No information exist on the photon flux before 1 Myr after the burst. Thus, $`N_0=N(1\mathrm{M}\mathrm{y}\mathrm{r})`$ is assumed so $`R_0`$ is equal to R(1 Myr) found from eq. 1. The calculations are performed for both of the IMF’s and the results are shown in table 1.
### 3.2 Supershell evolution
In the case of a single supernova event, it is rather simple to follow the four main evolutionary phases of its remnant analytically. First, the ejecta and the leading shock move almost undecelerated. Second, the velocity of the shock decreases, as a shell of swept-up ISM builds up and grows larger. In this phase, radiation from the hot cavity is not important, and the total energy within the shock is conserved. Third, radiative cooling of all matter within the shock become important. The shell is now driven by the pressure still existing in the cavity. Finally the PdV work from the cavity vanishes, so that the shell is now only being driven by the momentum. At some moment, the velocity of the shell has dropped to a value comparable to the rms velocity of the ISM, thus being non-detectable.
When one deals with the correlated effect of many supernovae, which is the case in starburst galaxies where the number of massive stars is large, additional considerations have to be made, since supernovae keep feeding energy to the cavity for about 50 Myr. The idea with the following calculations is to get a feeling of the timescales and radii of supershells. The analytical expressions from McCray & Kafatos are adopted to do the simple calculations, and the equations below are taken from their article, unless another reference is given. Strictly speaking, the analytical expressions may be used only if the ambient density is constant. However, the supershell first expands into the wind-driven bubble environment, then into a region with density depending on the H ii-region expansion, and finally into a non-expanding region outside the H ii-region. Hence, the evolution of the supershell is only followed while it is smaller than the H ii-region, and the density within the superbubble is assumed to be the same as outside. This last assumption is not realistic, as the passage of the wind-blown shock sweeps a large fraction of the gas into a superbubble. Hence, we may expect a more accurate model to show that the supershell reaches the superbubble at an earlier time. The gas gathered by the wind-driven shock cools to form the superbubble after relatively short time
$`t_b=(2.3\times 10^4\mathrm{yr})n_e^{0.71}L_{38}^{0.29}`$ (9)
, and the radius of the superbubble is
$`R_s=269pc(L_{38}/n_e)^{1/5}t_7^{3/5}`$ (10)
where $`L_{38}`$ is the characteristic wind luminosity in units of $`10^{38}\mathrm{ergs}\mathrm{s}^1`$ and $`t_7`$ is the elapsed time since the starburst in units of $`10^7`$yr. Strictly speaking, the time is counted from the start of superbubble expansion, but since O-stars start their wind-phase soon after the burst, $`t_7`$ is counted from the burst. $`n_e`$ is the matter density within the H ii-region, assumed to be homogeneous. The wind luminosity is found in Leitherer & Heckman from their fig. 55 (instantaneous starburst, metallicity is one tenth of solar, mass of burst is $`10^6M_{}`$, Salpeter IMF with upper limit 100 $`M_{}`$ and lower limit 1$`M_{}`$). At $`10^6`$ yr the wind luminosity is read-off to be $`10^{39}\mathrm{ergs}\mathrm{s}^1`$. Doing the same IMF scaling as in the photon flux discussion above, $`L_{38}=3.9`$ is obtained if $`m_L=0.1M_{}`$ and $`L_{38}=1.7`$ if $`m_L=0.01M_{}`$. The wind phase is important for as long as O-stars exist ($`5`$ Myr). However, after about 3 Myr the first supernovae explode and soon dominate the total energy output. The superbubble radius has been calculated 1, 2 and 3 Myr after the burst using eq. 10. The density typical for an H ii-region during its early evolution is $`10\mathrm{cm}^3`$, as discussed above.
After 3 Myr, the first SN appear. Hence, the expansion of the supershell is set to start at t=3 Myr. Strictly speaking, the energy output from supernovae appears as discrete events. However, it can be treated as continuous, as long as the interval between explosions is sufficiently short, at most in the order of $`10^5`$ yr (Tomisaka, Habe & Ikeuchi 1981). If the total mass of the burst is $`10^6M_{}`$, the number of stars with mass greater than or equal to 8$`M_{}`$ (lower limit of a SNII progenitor, according to Woosley & Weaver ) may be calculated using our IMF. The result using the 0.1 $`M_{}`$ lower mass limit is 8145 stars, whereas it is 3452 stars using $`m_L=0.01M_{}`$. The event lasts for about $`5\times 10^7`$yr so the mean interval between explosions is in the order of $`6\times 10^3`$ yr and $`1\times 10^4`$ yr, respectively. Thus, it is meaningful to assume continuous energy injection. Using this, eq. 10 may be used with just one replacement, namely the insertion of the mean supernova power instead of the wind luminosity. The mean power is calculated by dividing the total SN energy $`N_{}E_{51}\times 10^{51}`$ ergs with the total duration of the event. $`N_{}`$ is the number of stars with mass greater than 8 $`M_{}`$ and $`E_{51}`$ is the energy in units of $`10^{51}`$ ergs. Here, 0.4 is adopted as a representative value, see Woosley & Weaver \[1986, their table 1\]. The energy ejection is assumed not just to be continuous, but also constant in time. This is probably quite a good approximation, see Leitherer & Heckman (1995). The considerations lead to
$`R_s`$ $`=`$ $`97\mathrm{pc}(N_{}E_{51}/n_e)^{1/5}t_7^{3/5}`$ (11)
$`V_s`$ $`=`$ $`5.7\mathrm{km}\mathrm{s}^1(N_{}E_{51}/n_e)^{1/5}t_7^{2/5}`$
The time of shell formation is given by an expression similar to eq. 9, obtained by replacing the wind luminosity with the mean supernova power.
The density of the H ii-region is assumed constant in time and equal to 3 $`\mathrm{cm}^3`$, a reasonable value during late stage H ii-region evolution (see table 1). Eqs. 11 are used in the adiabatic phase only (no cooling), since they implicitly assume that the loss of energy within the shell is negligible. Cooling, and hence energy loss, from the interior becomes important at a time
$`t_c=4\times 10^6\mathrm{yr}\xi ^{1.5}(N_{}E_{51})^{0.3}n_e^{0.7}`$ (12)
At this time, the shell has a radius
$`R_c=50\mathrm{pc}\xi ^{0.9}(N_{}E_{51})^{0.4}n_e^{0.6}`$ (13)
found by inserting eq. 12 into eq. 11. $`\xi `$ is the metallicity in units of the solar metallicity. For the objects in question, $`\xi =0.1`$ is a representative value. Even if a small starburst mass of 10 $`M_{}`$ is assumed and the density of the surrounding medium is 3 $`\mathrm{cm}^3`$, it turns out that the supershell reaches the end of the adiabatic phase at a time $`2.1\times 10^7`$ yr if $`m_L`$=0.1 is used and $`1.6\times 10^7`$ yr if $`m_L`$=0.01. Hence, the interior of the supershell starts to cool after the H ii-region has vanished. This tells us that we may use eqs. 11 to calculate the size of the supershell, as we are only interested in the evolution during H ii-region existence. The results are given in table 1.
The measured emission lines may indeed originate in the part of the H ii-region that is still outside the superbubble, since the size of the H ii-region always exceeds the size of the superbubble as seen in table 1. It also exceeds the size of the supershell until 6 Myr after the burst, though only significantly within the first 5 Myr. The above calculations do not prove that the emission lines arise outside the superbubble/supershell - they just confirm that it is a possibility. If emission lines really are measured this way, they are not affected by enrichment from the present burst, thus opposing self-enrichment.
### 3.3 Propagating star formation
Propagating star formation means that the energy deposited in the ISM by an evolving star formation event, is initiating new star formation. As a supershell grows, it will become Rayleigh-Taylor unstable and fragments. The situation has been treated by Elmegreen , who found an expression for the timescale of cloud-collapse in the supershell:
$`t_{cloud}=103\left({\displaystyle \frac{n_0}{\mathrm{cm}^3}}\right)^{1/2}\mathrm{Myr}`$ (14)
where $``$ is the expansion speed of the shell divided by the rms velocity dispersion in the shell. For an adiabatic shell, $``$ is equal to 1.8 . If $`n_0=10\mathrm{cm}^3`$, this gives the result that star formation is expected to start not earlier than $`24`$ Myr after shell formation. On the other hand, if $`n_0=3\mathrm{cm}^3`$ as used in the calculations above for a typical density in the late stages of H ii-region evolution, the timescale is more like 44 Myr. The mass of the supershell is given as
$`M=\rho _0V_{ss}1.3n_0m_p{\displaystyle \frac{4\pi }{3}}R^3`$ (15)
assuming all of the ISM originally within $`R`$ to be incorporated in the shell. $`\rho _0`$ is the mass density of the ambient medium, $`V_{ss}`$ is the volume occupied by the hot phase within the supershell, $`m_p`$ is the proton mass and 1.3 is the approximate mass per particle in units of the proton mass. The radius is of the order a few 100 pc as shown above, so eq. 15 gives a shell mass of the order $`10^6M_{}`$. Hence the induced star formation is of the same order of magnitude as the original central burst.
Observational indications on the existence of star forming supershells are numerous. One recent example is NGC 2537 included in the sample of Martin . In $`H_\alpha `$, it shows a very clear spherical distribution of starformation sites. As a curiosity, it may be mentioned that Mori et al. calculated the effects of star formation in an expanding supershell, using a 3-D hydrodynamical code including a dark matter halo. The results were in remarkable accordance with available observations of dE’s, such as exponential surface brightness profile, positive metallicity gradient and inverse color gradients.
If the ejected metals from the central burst mix with the material of the supershell, the H ii-regions of the induced burst will show abundances differing from the central H ii-region. Although there have been some doubt whether the metals are allowed to mix into the supershell , and if so when this will happen, we have been inspired by this possibility to let our model have a double burst nature.
## 4 The model
We have designed a model to fulfil certain demands, namely that the observed trend of constant N/O for low metallicities and increasing N/O for higher metallicities should be reproduced and also that the observed scatter in N/O should be explained. Fixing all the parameters by fitting the N/O-O/H evolution, the model should also be able to explain the observed helium mass fraction as a function of O/H and O/H as a function of gas fraction.
The included elements are H,He,C,N,O. The production of N is still a hot topic, since the degree of primary production at various metallicities is unclear. Further, N is produced mainly by intermediate mass stars while O is produced by massive stars, making it necessary to include stellar lifetimes. However, the inclusion of both N and O provides the opportunity of constraining stellar parameters and mixing scenarios.
To account for the scatter in N/O, it is suggested that the bursts are instantaneous and ordered in pairs, hence using the delay between the ejection of N and O. The principle of time delay was used by Garnett , though employing single bursts only. It is doubtful whether single bursts produce scatter, since at least some abundance observations have to be done at the time of N-release. However, this is between two bursts, where no giant H ii-regions are present, providing no possibility for abundance observations by emission lines. Single bursts may produce scatter only if IMF parameters or enriched wind efficiencies vary from galaxy to galaxy. However, there is no indication of IMF variations, and the existence of enriched winds, on the whole, is questionable, and may be important in extreme low-mass galaxies only - see below. Double bursts produce the scatter quite naturally: It was found reasonable above that a localized burst results in another burst, this time in the expanding supershell, surrounding the original burst. The timescale for star formation in the supershell is found to be $`2.4\times 10^7`$ yr, comparable to the timescale for O-ejection, but shorter than the timescale for N-ejection, hence producing the desired scatter.
In our model we assume that the very first burst is a single one, all other bursts appear in pairs. The interburst period between the two bursts of a pair is tuned to give maximum scatter, found below to be 30 Myr. The time between two pairs is set to 1 Gyr.
### 4.1 The IMF
The IMF used throughout this paper is the single-power Salpeter-IMF:
$`\varphi (m)=\varphi _0m^{2.35}`$ (16)
. $`\varphi _0`$ is the normalization constant, found by $`_{m_L}^{m_U}m\varphi (m)𝑑m=1`$. Using a lower mass cutoff $`m_L=0.1M_{}`$ and an upper mass limit $`m_U=100M_{}`$, $`\varphi _0`$ is equal to 0.17. These limits are standard values, very often used in the literature. A Salpeter IMF still fits the observations (above about $`1M_{}`$) quite well despite of its age and simple appearance .
The lower-mass cutoff is an important parameter because the yield from a generation of stars is a function of the adopted lower mass. The reason is that stars less massive than $``$1 $`M_{}`$ do not eject metals, thus locking-up all the material from which they are formed. Hence, lowering the lower-mass cutoff implies a lower recycling fraction, giving a lower yield. As will be apparent from our results, it becomes attractive to invoke another value of $`m_L`$ namely 0.01 $`M_{}`$. In this case the normalization constant is found to be 0.07. Thus, the yield will be a factor 2-3 lower when using the lower cutoff.
Throughout this paper, a metallicity invariant Salpeter IMF has been used, since there are at present several observational indications of an abundance invariance . The number of stars with mass \[$`m_1`$,$`m_2`$\] is given by
$`N([m_1,m_2])`$ $`=`$ $`\varphi _0M_{burst}{\displaystyle _{m_1}^{m_2}}m^{2.35}𝑑m`$ (17)
$`m_1`$ $`=`$ $`m_j{\displaystyle \frac{\mathrm{\Delta }m}{2}}`$
$`m_2`$ $`=`$ $`m_j+{\displaystyle \frac{\mathrm{\Delta }m}{2}}`$
where $`M_{burst}`$ is the mass of a burst, i.e. the mass of gas turned into stars. The stellar mass grid may be as fine as one wishes. For use in the present model, a grid of 60 different stellar masses has been adopted, since this was found to give sufficient mass resolution for the present purpose.
The mass of each burst may be constrained by the available observations of SFRs. The SFR for BCGs is typically 0.1 - 1 $`M_{}\mathrm{yr}^1`$. For the mean SFR of a double burst to be within this range, the masses of the bursts may be estimated from $`M_{burst}(1)+M_{burst}(2)/t_{ib}0.11`$, where $`t_{ib}`$ is the time between the two bursts of a pair. If $`t_{ib}=3\times 10^7`$ yr, and assuming $`M_{burst}(1)M_{burst}(2)`$, this gives $`M_{burst}1.515\times 10^6M_{}`$, in agreement with Marlowe et al. \[1995, their table 7\].
### 4.2 The equations
The calculation of abundances is carried out just before every burst. The reason for doing this is simply that we wish to compare with the observed abundances, being measurable in H ii-regions only. The adopted equations are described in the following. The notation is similar to the one in Pilyugin . The first step is to calculate the masses of each element present in the gas phase just before every burst. For the $`j^{th}`$ burst:
$`M_i(t_j)`$ $`=`$ $`M_i(t_{j1})\mathrm{\Delta }M_{j1}X_i(t_{j1})`$
$``$ $`\mathrm{\Delta }W_i(t_{j1})+\mathrm{\Delta }W^{Inf}(\tau _{j,j1})X_i^{Inf}`$
$`+`$ $`{\displaystyle \underset{k=1}{\overset{j1}{}}}\mathrm{\Delta }M_k(Q(\tau _{j,k})Q(\tau _{j1,k}))X_i(t_k)`$
$`+`$ $`{\displaystyle \underset{k=1}{\overset{j1}{}}}\mathrm{\Delta }M_k(Q_i(\tau _{j,k})Q_i(\tau _{j1,k}))`$
$`t_j`$ refers to the time just before the $`j^{th}`$ burst. $`X_i(t_k)`$ is the abundance of element $`i`$ just before the $`k^{th}`$ burst. $`\tau _{j,k}`$ is defined as $`t_jt_k`$, so it is the time elapsed since the $`k^{th}`$ burst. Thus the element yields from stars with lifetimes shorter than $`\tau _{j,k}`$ has to be included in the two summations. $`Q(\tau _{j,k})`$ is the mass fraction of gas ejected from the $`k^{th}`$ generation of stars, just before starburst $`j`$. The composition is left unchanged by stellar nucleosynthesis, thus being the same as in the gas from which the stars were formed. $`Q_i(\tau _{j,k})`$ is correspondingly the mass fraction ejected of element $`i`$, but newly synthesized. It is clear from this notation that $`Q(\tau _{j,j})=Q_i(\tau _{j,j})=0`$, since stellar lifetimes are finite. The first term on the right side is the mass of element $`i`$ just before the $`(j1)^{th}`$ burst. The second term is the mass of element $`i`$ that has been turned into stars at burst $`j1`$. The first summation term is the mass of element $`i`$ ejected, without being changed, by the $`k^{th}`$ burst in the period between the $`(j1)^{th}`$ and $`j^{th}`$ burst. The second summation term is corresponding to the first one, except for the fact that the mass of element $`i`$ was newly synthesized.
$`\mathrm{\Delta }W_i(t_{j1})`$ is the mass of element $`i`$, leaving the galaxy as a result of starburst $`j1`$. This term consists of two parts: The part belonging to the ordinary wind, and the one corresponding to the enriched wind. Thus, the wind term may be written as
$`\mathrm{\Delta }W_i(t_j)`$ $`=`$ $`W_{ISM}\mathrm{\Delta }M_jX_i(t_j)`$
$`+`$ $`W_{SN}(Q_i(t_w)+Q(t_w)X_i(t_j))\mathrm{\Delta }M_j`$
. $`W_{ISM}`$ and $`W_{SN}`$ are the efficiencies of the ordinary and enriched winds, respectively. Both efficiencies are zero if the model is closed. The physics behind enriched winds is based on the principle that the supershell following a burst, breaks up, allowing the hot ejecta to blow out and escape from the galaxy. It is assumed that the wind is caused by supernovae type II, so the end of the wind phase $`t_w`$ is the lifetime of the least massive star exploding as a SNII, which is set to $`8M_{}`$, consistent with the mass adopted for our supershell calculations. Hence, for stars with masses above this limit, a fraction $`W_{SN}`$ of the ejected mass of element $`i`$ is leaving the galaxy. For stars less massive, $`W_{SN}=0`$. It is important to notice, in accordance with the physics involved that this efficiency factor is the same for all elements considered, but since oxygen is dominating the ejecta from SNII, the winds will be enhanced in oxygen.
The other possibility is ordinary galactic winds, arising as a consequence of a general heating of the ISM, causing a fraction to leave the galaxy, i.e. the composition of the gas leaving the galaxy is the same as the composition of the ISM. The first term of eq. 4.2 is the instantaneous-burst representation of Hartwick-outflow , giving a mass-loss proportional to the burst mass. By multiplying the mass of the wind with the fraction of the $`i^{th}`$ element in the ISM, taken just before the burst that is responsible for the wind, one obtains the mass removed of element $`i`$ due to an ordinary wind. $`W_{ISM}`$,$`W_{SN}`$ and the burst masses are treated as free parameters, though one restriction is made, namely that the yields of the first burst should be sufficient to place the second burst approximately at the abundances of IZw18.
$`\mathrm{\Delta }W^{Inf}(\tau _{j,j1})X_i^{Inf}`$ is the increase in mass since the last burst of element $`i`$ in the ISM due to inflow of gas. The composition of the infalling gas, given by $`X_i^{Inf}`$, is taken to be primordial. The inflow rate is given by
$`\dot{M}(t)={\displaystyle \frac{M_0}{\tau _{inf}}}e^{t/\tau _{inf}}`$ (20)
. $`M_0`$ is the total mass accreted for $`t\tau _{inf}`$ and $`\tau _{inf}`$ is the accretion timescale. The mass accreted at time $`t`$ is found by integrating eq. 20 from 0 to $`t`$ finding $`M(t)=M_0(1e^{t/\tau _{inf}})`$, finally giving
$`\mathrm{\Delta }W^{Inf}(\tau _{j,j1})`$ $``$ $`M(t_j)M(t_{j1})`$ (21)
$`=`$ $`M_0(e^{t_{j1}/\tau _{inf}}e^{t_j/\tau _{inf}})`$ (22)
This is unfortunately giving two free parameters further, namely $`\tau _{inf}`$ and $`M_0`$. In models not including inflow $`\mathrm{\Delta }W^{Inf}(\tau _{j,j1})`$=0. The mass of gas just before starburst number $`j`$ is
$`M_g(t_j)`$ $`=`$ $`M_g(t_{j1})\mathrm{\Delta }M_{j1}\mathrm{\Delta }W(t_{j1})`$
$`+`$ $`\mathrm{\Delta }W^{Inf}(\tau _{j,j1})`$
$`+`$ $`{\displaystyle \underset{k=1}{\overset{j1}{}}}\mathrm{\Delta }M_k(Q(\tau _{j,k})Q(\tau _{j1,k}))`$
Finally the mass of the entire galaxy is
$`M_{dw}(t_j)=M_{dw}(t_{j1})+\mathrm{\Delta }W^{Inf}(\tau _{j,j1})\mathrm{\Delta }W(t_{j1})`$ (24)
For the closed models $`M_{dw}(t_j)=M_{dw}(0)`$.
$`Q_i`$ and $`Q`$ are adopted from models of stellar evolution as described in section 4.3. However, $`Q_H(t)`$ is calculated from
$`{\displaystyle \underset{i=H,He,C,N,O}{}}Q_i(t)=0`$ (25)
assuming all newly consumed H to go into the production of He, C, N and O. This simplification is realistic because the involved elements are by far the most important.
From the above equations, it is now possible to calculate the abundances by mass in the ISM, just before a new burst:
$`X_i(t_j)={\displaystyle \frac{M_i(t_j)}{M_g(t_j)}}`$ (26)
However, the observed abundances are not given by mass, but by number. The abundance by number of, say $`O(t_j)`$ relative to $`H(t_j)`$ is
$`{\displaystyle \frac{O}{H}}(t_j)={\displaystyle \frac{M_O(t_j)}{M_H(t_j)}}{\displaystyle \frac{A_H}{A_O}}`$ (27)
where A represent the atomic masses and M is the gas phase masses calculated by eq. 4.2.
In the calculations it is assumed that the ejecta mix into the entire interstellar medium before the new burst appear. It should be noted that this assumption is doubtful for the short interburst period between the two bursts of a pair. The fact that the metallicity of second generation stars is slightly higher, is included in the sense that the adopted yields are metallicity dependent. The calculations are terminated if the amount of gas is insufficient for a new burst, or the age of the dwarf galaxy is more than 15 Gyr.
The initial conditions to be put into the equations are the following. The very first burst (the single burst) takes place at primordial gas composition, i.e. Z=0, $`X_{He}`$=0.243 \[1997a\], $`X_H`$=0.757, $`X_C`$=$`X_N`$=$`X_O`$=0.
Before the first burst, the galaxy consists of gas only. Hence, the total mass of the galaxy is equal to the mass of gas, in all calculations adopted to be $`10^8M_{}`$. However, the situation is a little different for the inflow models. In this case, the initial mass of a galaxy is adopted to be zero. Gradually, the galaxy increases its mass by inflow of primordial gas, until the gas mass reaches a limit, sufficient to start the first burst. To be consistent, this mass limit has been adopted to be $`10^8M_{}`$. So the only difference from models not including inflow is that the first burst is delayed by the time it takes to build up a sufficiently large gas cloud.
### 4.3 The adopted yields and their implementation
Since the present models follow the evolution from the very first burst at Z=0 to about solar metallicity, metallicity-dependent yields have been adopted. The Renzini & Voli yields for low- and intermediate mass stars are the most widely used for the purpose of chemical evolution modelling. Their strength is that they are giving the yields for several choices of convection and wind parameters. Furthermore they give the primary yield of nitrogen (and $`{}_{}{}^{13}C`$) separately, which is very important for our purpose. Indeed, one of the major problems to solve for low metallicity objects is the extent of primary produced N, compared to the secondary production. Marigo, Bressan & Chiosi (1998) made also a separation between primary and secondary production in the mass interval 4-5$`M_{}`$. The question whether stars more massive than 5$`M_{}`$ produce primary N is controversial. The question is whether it is necessary at all to introduce such a source for primary N production. This is one of the questions we address in this paper. The yields by van den Hoek & Groenewegen , also for low- and intermediate mass stars, do not include separate results for primary nitrogen, so the yields given are a mixture of primary and secondary components. This is the reason, why we have not adopted their yields for more than just a reference, see section 5.
For use in the models, yields are selected in a consistent way, i.e. yields from Geneva tracks and yields from Padova tracks (Marigo, Bressan & Chiosi 1996; Marigo et al. 1998; Portinari, Chiosi & Bressan 1997) are kept in separate sets of yields. The adopted yields are organized as shown in table 2.
Note that not all yields in sets 1, 2a and 2b are deduced directly from the Geneva tracks, as van den Hoek & Groenewegen , Renzini & Voli and Woosley & Weaver use different indirect methods.
The problem is to calculate the $`Q`$ and $`Q_i`$ terms just before every burst, given the IMF and the ages of the previous bursts from which the stars eject the elements. The mass ejected of an element $`i`$, at a time $`\tau `$ after a burst is calculated as a sum of contributions from each star down to a stellar mass $`m_\tau `$, corresponding to a lifetime $`\tau `$. The lifetime of each star is found by linear Lagrange interpolations in both mass and metallicity between the values given by the stellar tracks. The following expression is used for calculating $`Q_i`$
$`Q_i={\displaystyle \underset{m_j>m_\tau }{}}q_i(m_j)N(m_j)`$ (28)
where $`N(m_i)`$ is the number of stars with mass $`m_i`$, found with the adopted IMF, and $`q_i(m_i)`$ is the mass ejected of element $`i`$ from stars having this mass. The $`q_i`$s are found from the references in table 2 in two steps. First, a linear interpolation is performed between the values for the available metallicities. If the metallicity is lower than the metallicities, for which the yields are known, primary and secondary components of N are treated differently. The yields for primary elements are taken to be equal to the yields at the lowest metallicity with known yields, except for set no. 3, in which the primary N is obtained by extrapolating linearly from the two metallicities with known yields. For this set, an extrapolation should be more correct than just assuming the yield to be constant below the range of metallicities covered by the yield sources, since the range does not extend below Z=0.008 for the critical mass interval 4 - 5 $`M_{}`$. The price to pay is an extrapolation reaching far beyond the covered metallicity range. Note that stars with masses less than 4 $`M_{}`$ or larger than the upper mass for hot bottom burning (free parameter in the model for set 2a and 2b, and 5 $`M_{}`$ for set 3) are only producing secondary N. The secondary N yield is interpolated between the lowest metallicity with known yields and Z=0, since the yield of secondary N is 0 at Z=0, according to the definition of secondary production. For H, He, C and O the yields interpolated/extrapolated are the sum of primary + secondary yields. This is without influence on the results.
Second, the yields are interpolated linearly with respect to initial stellar mass, finally giving the $`q_i`$s. For 2a and 2b, the mixing length parameter $`\alpha `$, defined as the convection mixing length divided by the pressure scale height, is a free parameter. We know the yields for two different values of $`\alpha `$ from Renzini & Voli , namely $`\alpha `$=0 and $`\alpha `$=1.5. Thus when implementing the yields, the $`q_i`$s are found in three steps. Before the two steps described above are carried out, the yields given in the mass interval 4 $`M_{}`$ to the upper mass of HBB, $`m_{HBB},`$variable between 5 and 8 $`M_{}`$, are linearly interpolated with respect to $`\alpha `$.
The implementation of the mass ejected unprocessed is very similar. The expression for the total mass ejected since the starburst is written as
$`Q={\displaystyle \underset{m_i>m_\tau }{}}(m_im_{rem})N(m_i)`$ (29)
$`m_{rem}`$ is the mass of the stellar remnant (white dwarf, neutron star or black hole). Again the value of $`m_{rem}`$ is extracted from the sources by linear interpolation, first with respect to metallicity, then between initial stellar masses.
It is very useful to calculate the true yields of selected elements because they allow for an easy comparison between the nucleosynthetic outcome of different stellar models. The true yield of an element is defined to be the mass of an element ejected from a star generation, divided by the fraction locked up in stellar remnants and low-mass stars.
The total mass ejected of an element, e.g. O is given by $`_{m_i}q_O(m_i)N(m_i)`$, where $`q_O(m_i)`$ is found from the sources following the proces described above. The lock-up fraction is equal to $`1`$R, where R is the return fraction, i.e. the mass ejected in units of the total mass of the burst $`Q/M_{burst}`$. The resulting true yields are given in table 3 for He and O, as we shall need these quantities later. Note that the true He yields for set 3 are higher than those of 2a and 2b, in particular at the higher metallicities. Note also the lower yields, when using $`m_L`$=0.01, by a factor of $``$2.8.
### 4.4 Stellar lifetimes and element timescales
In the case of set no. 1, 2a and 2b, the stellar lifetimes are adopted from Schaller et al. (Geneva-tracks), while in set no. 3 they are from Portinari et al. (Padova-tracks). To compare the two sources, lifetimes are shown for some selected metallicities and masses in table 4.
It is possible directly to compare the stellar lifetimes at solar metallicity. Low- and intermediate mass stars seems to have shorter lifetimes using the Padova-tracks (between $``$40 to $``$70 per cent), while the opposite is true for massive stars (but $``$10 per cent longer). This age difference may be a result of different stellar wind-efficiencies, and the inclusion of overshooting in the Padova-tracks.
By combining our knowledge of stellar yields with the respective lifetimes as function of mass, we obtain information on the ejection timescales of different elements. For the present purpose, it is the difference between the timescales of N and O ejection that is particularly interesting, since the time when the ejected N/O has its minimum is equal to the time, where the second burst of a pair should appear if maximum scatter is to be obtained. The timescales are found for a fixed metallicity at which the ejected masses of O and N of one burst are calculated in small timesteps (of the order of $`10^6`$ yr) until all stars down to 1 $`M_{}`$ have ended their life cycle. The time evolution of the element ejection has been calculated for set no. 2a and 3 for two metallicities. Fig. 4 shows the timescales for set 2a. The results for set 3 are very similar.
The production of N in massive stars is different for the two metallicities. For solar metallicity, it is released much faster for $`\alpha `$=1, i.e by more massive stars. This is only secondary N though, since the contributors have lifetimes shorter than 50 - 80 Myr, hence being more massive than 8 $`M_{}`$ (compare with the stellar lifetimes given in table 4). For set 2a, it is not surprising that the ejection timescale of N is much shorter at low metallicity if $`\alpha `$=0.0 (dashed lines in fig. 4) than if $`\alpha `$=1.0, because no primary nitrogen is ejected from intermediate mass stars when $`\alpha `$=0.0. The major N production sets in at about 100 Myr at the low metallicity. For solar metallicity a large fraction is produced within the first 50 Myr, about 40 (2a) to 60 per cent(3). Almost all of the oxygen is produced in $``$30-40 Myr for all sets and metallicities.
Hence, in order to choose the interburst period between the bursts of a pair, a minimum value of N/O would appear somewhat between 30 and 80 Myr for the low metallicity and between 30-40 Myr for the high metallicity. If a choice has to be made, 30 Myr would be the best to make, since, in general, it is easier to reproduce scatter at low metallicities due to the relative effect of ejected metals to metals already present in the interstellar medium. Following this conclusion, the time between the two bursts of a pair is set to 30 Myr.
## 5 Results
The model is applied to fit the observations in the N/O-O/H plane with the use of the yields in table 2. The accompanying results on Y-O/H and O/H vs. gas fraction are displayed as well. The order of presentation starts with the results of the closed model.
### 5.1 The closed model
The closed model has been applied to all of the yield sets. In most plots, the IMF with $`m_L`$=0.1 has been used. For those employing the 0.01 low-mass limit, it is mentioned separately in the text. It is not necessary to do all calculations for both IMF’s, since the results are the same, except for the gas fractions. For all elements, the only difference is a lower yield using the low value of $`m_L`$, hence more massive bursts are needed to produce the same abundances as when the higher low-mass limit is used. This statement is confirmed in the following.
#### 5.1.1 N/O-O/H
Fig. 5 shows the evolutionary path of a double-bursting dwarf galaxy, as it has been calculated by the model using yield set 1. The point where the evolutionary path starts (at 12+log(O/H)$``$7.1) is where the second burst appears, 1 Gyr after the first burst. The next point is where the third burst would be observed, 30 Myr after the second burst. Thus, the upper points of the saw-tooth pattern are corresponding to the first bursts of each pair, and the lower points to the second pair-bursts. The scatter is more pronounced at low metallicities because the mass of O ejected relative to the one already present in the ISM is higher, the lower the metallicity. It is clear from this figure, why set no. 1 is used as a reference only. The level of N/O is much too high, in particular at low metallicities. The only way one can lower the ratio is by assuming a top heavy IMF, since the yields from van den Hoek & Groenewegen do not allow for changing the convection parameters or provide possibilities for distinguishing between primary and secondary components. Thus, it is not possible to scale the secondary components separately to Z=0. Further, neither outflow nor inflow is able to cure the problem, in particular not if selective winds are considered, since this would increase N/O even more.
It is interesting to note that the parameters (convection, wind etc.) used to fix the yields in van den Hoek & Groenewegen have been found by using a synthetic model on AGB stars in the LMC and in the Galactic disc. The LMC is not represented in fig. 5, but according to Pagel et al. 12+log(O/H)= 8.36 and log(N/O)= -1.22 for LMC, which is fitted fairly well after 4 pairs of bursts (but may be fitted by 2 pairs of bursts if the mass of each burst is larger), thus being consistent. The problem seems to arise at lower metallicities.
Next, set 2a is applied (fig. 6). The upper mass limit of hot bottom burning (HBB) is set to $`m_{HBB}=5`$$`M_{}`$, compatible with set no. 3. In this case, a value of $`\alpha `$=1.1 is necessary to obtain the right level of N/O. A larger value of the convection parameter would produce too much primary N and vice versa.
Both theoretical predictions and observations of AGB stars in the Large Magellanic Cloud are suggesting a level of HBB, roughly corresponding to $`\alpha `$=2.0 \[1997, and references therein\], which is well above the value obtained here. If $`\alpha `$ really is so large in AGB stars in the LMC, an explanation might be that the value of the convection parameter depends on metallicity, since the metallicity of the LMC is high (about half-solar). A higher value of $`\alpha `$ at high metallicities is not ruled out by the model as long as it is equal to 1.1 during the first few bursts.
A large spread is evident and sufficient to explain the observed scatter. Here, the plusses, representing the DLA-systems, are neglected. A similar run is presented in fig. 7, only employing the IMF having $`m_L`$=0.01 $`M_{}`$. As expected, it is seen that the result resembles the situation in fig. 6, if one assumes all bursts to involve twice the mass of the bursts in fig. 6.
Figures 6 and 7 give a good comprehension of the outcome of the model, but some considerations should be made before any solid conclusions are drawn. Firstly, abundances are measured in H ii-regions only, and therefore only at the ’positions’ of the bursts, not in between. Secondly, the evolutionary path represents the evolution of one dwarf galaxy. Another dwarf galaxy, having different parameters such as total mass, burst masses etc., would have another evolution, and thus another path in the N/O-O/H plot. Thus, fig. 8 shows the evolutionary paths of dwarf galaxies having different fractional burst masses. Yield set 2a has been used, and all of the parameters, except the mass of each burst, are the same.
To make the plot more clear, lines have been drawn through the first bursts of each pair, and through the second bursts. The former lines are almost coincident, because the time interval between two pairs of bursts is $``$ 1 Gyr, enough to release almost all of the nitrogen. Thus, larger bursts increase both N and O, adjusting N/O to be roughly unchanged. The case is different for the lines through the second pair-bursts, because the N from the first pair-burst has not yet been released after 30 Myr. Thus, a larger burst mass gives an increase in O abundance only.
Two results may be deduced from fig. 8. Firstly, the scatter of the observations is perfectly explained, because an observational scatter of $``$0.1 dex has to be taken into account. Furthermore, the area between the lower line corresponding to 5 per cent burst masses and the upper line is ’filled’ in the sense that the distribution of burst masses may be everything between zero and 5 per cent to match the observations. Thus, the distribution of abundances is accounted for by using the models of dwarf galaxies having various reasonable relative burst masses. The other result is the slight increase of the evolutionary path toward higher metallicities. However, the small upturn may not be a consequence of dominating secondary N production, but rather the decreasing O yield at higher metallicities, see table 3.
Set no. 2b is different from no. 2a only in that a modest mass-loss from massive stars is assumed. The effect of this change on the evolutionary path is seen by comparing fig. 9 with fig. 8. It is possible directly to compare the model outputs of the two figures, since all parameters, including the burst masses, are identical. When using set 2b instead of 2a, the N/O ratio has an even smaller upturn at higher metallicities, because the O-yield is higher. Thus, even when N production is increasing at higher metallicities, so is the oxygen, keeping the ratio down and thereby hiding the secondary behavior of N. Still, the conclusions are the same: both the level of N/O and the scatter is explained, independent of $`m_L`$.
Using yield set no. 3, the same kind of plot is constructed as those above. Here there is no explicit assumption on the convection parameter, but for reasons explained below, the primary component of N (and O) for Z$`<`$0.008 has been extrapolated linearly, using the yields of Z=0.008 and Z=0.02 instead of just adopting the primary yield at Z=0.008. The run of the model is presented in fig. 10.
At least two features have to be mentioned. The first is that the scatter is somewhat smaller than for set 2a (or 2b) and the second is an upturn in N/O at higher metallicities, being more pronounced, than it was for set 2a and 2b. Unfortunately, the upturn is taking place for higher oxygen abundances than those of the observations.
The primary component in set 3 has also been modelled in another way: For the same choice of parameters as in fig. 10 (solid line), the model has been applied with primary yields equal to the one at the lowest metallicity for which the yields are given (Z=0.008), with the resulting evolutionary path plotted in fig. 11. A comparison shows that an adoption of the primary yields at Z=0.008 will produce too little primary N. It seems to be necessary to make primary production more effective, the lower the metallicity, as it is done by the linear extrapolation, because the primary N yield at Z=0.008 is slightly larger than it is for solar metallicity. This is actually an important result, because it provides a constraint on the produced primary N, independent of stellar models. The explanation for this result might be that the number of thermal pulses in the AGB phase of low metallicity stars is higher, thus converting more C and O into N via the CNO-cycle. Hence, in the following, all model runs using set 3, will use linear extrapolations of primary N yields, below Z=0.008. Of course, one could make the IMF steeper thereby producing less oxygen. However, observations do not support a steeper IMF.
When comparing the outputs of the model, using set 2a (2b) and 3, it is important to remember that in set 2a and 2b, the lowest metallicity for which yields are available is Z=0.004, at least in the critical, and very important mass interval 4 to 5 $`M_{}`$, while it is as high as Z=0.008 in the same mass interval for set 3. Thus, whether primary N yields at Z $`<`$0.008 are extrapolated from Z=0.02 over Z=0.008 or are adopted to be equal to the yield at Z=0.008, the result is quite uncertain, and stellar yields at much lower metallicities are definitely desirable.
#### 5.1.2 Y-O/H
It is worth emphasizing that the helium abundances are fitted with input parameters not different from those used to fit the N/O ratio.
For the closed model presented in fig. 8 by the solid line, Y has been calculated as well. Y against O/H is shown in fig. 12. The evolutionary Y path does not show any scatter, though the dwarf galaxy is double bursting, because about half of the helium is produced in massive stars, together with the oxygen. The fit is very good and well within the error bars. For the two most metal deficient objects in the plot, IZw18 and SBS0335-052, it seems to be necessary to assume smaller burst masses, say 1 per cent of the total mass - the model with the dashed lines in fig. 8. Thus, these objects can be explained by experiencing their second burst and having burst masses equal to 1 per cent of their total mass. The output of the model gives the abundances just before every burst. Hence, to calculate the model value of dY/dZ, the Y,O/H values were fitted using a linear least-square fit, assuming Z=20(O/H). The result is 1.0, lower than the (very uncertain) value of 2.6 derived from the observations, see section 2. The corresponding result using set 2b is not shown, since it resembles the one by set 2a. Certainly, it is true that the He yield is smaller for modest mass loss, but only for very massive stars and high metallicities .
The dependency of He on O is more pronounced, when using the Padova yields, as seen in fig. 13.
Calculating dY/dZ, using this yield set, in the same way as done for set 2a one obtains 3.1, higher than the value calculated for the observations (2.6). From table 3 it is clear that the true He yields for set 3 are a factor of $``$3 higher than those of sets 2a and 2b, but it is unclear why the Padova tracks produce that much helium. This feature makes the fit not quite as good as it was for set 2a (2b). However, the primordial helium determination is only certain within 0.003 \[1997a\]. Setting the primordial He abundance equal to the lower limit (0.240) in the initial conditions of the model does not change the N/O-O/H evolutionary path, but it does lower the level of Y in fig. 13. The new evolutionary path of Y has the same slope as the solid one of course, but sufficiently lower to be a fairly good fit, as it appears well within the observational uncertainties.
#### 5.1.3 O/H - $`\mu `$
Gas fractions may be useful in serving as a further constraint on chemical evolution models. It has been pointed out by several authors that closed models, such as the one presented above, are unable to explain the observed gas fractions . Thus, gas fractions have been calculated just before every burst using eqs. 4.2 and 24 allowing for comparison with some observed values from the literature.
Problem lies in the interpretation of the observations. The total dynamical mass of a dwarf galaxy also includes dark matter (DM). This dark is unlike to participate in the chemical evolution. Our model is including baryonic matter only. Thus, in order to compare the observations with the calculated gas fractions, one has to be sure that no significant amount of non-baryonic DM is included in the total mass estimates. However, as pointed out by several authors dwarf galaxies, including BCGs, are DM-dominated. Consequently, dynamical estimates of total masses are useless, when considering dwarf galaxies. Carigi et al. tried to solve the problem by using dynamical estimates within the visible Holmberg radius, to avoid inclusion of a DM halo. Unfortunately, this leaves us with a very small sample of observations, and still the amount of DM within the Holmberg radius is uncertain.
However, gas fractions are one of the few possible ways of constraining chemical evolution models, so a discussion of gas fractions, using the same models as in the preceding paragraphs should be included. To avoid erroneous conclusions, we used lower and upper limits of the gas fractions.
The gas fractions used are presented in table 5. The gas masses included are HI-masses multiplied by 1.3 to account for the helium content. The amount of $`H_2`$ is ignored. All dynamically estimated gas fractions have been scaled to a Hubble constant $`H_0=65\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ using $`M_{gas}(distance)^2`$ and $`M_{tot}(distance)`$, so the gas fraction $`\mu (distance)H_0^1`$.
DM is implicitly included, and molecular hydrogen is ignored, hence these gas fractions must be regarded as lower limits. Unfortunately, even the lower limits seem to be very uncertain. For instance, the gas fraction of IZw18 was first calculated using the data from Staveley-Smith, Davies & Kinman (1992). However, after He-correction and $`H_0`$-scaling, the gas fraction became larger than 1. This is also the case for some of the galaxies in the Thuan & Martin sample (not included in the present sample). In the case of IZw18, the value 15 per cent gas has been adopted from Matteucci & Chiosi .
Extreme upper limits have been obtained assuming that each of the galaxies are experiencing their first burst, and are observed during maximum luminosity. Indeed, this is an extreme upper limit, and should be regarded as such. Adopting the evolutionary stellar models of Leitherer & Heckman , the $`M/L_B`$ ratio may be estimated for a burst, $`M`$ being the mass of stars. Leitherer & Heckman calculated the blue luminosity for a $`10^6M_{}`$ starburst as a function of time. To be consistent, their Salpeter IMF has been scaled to ours. Indeed, two different IMF’s have been used so far, but to be sure that upper gas fraction limits are obtained, the high lower limit, $`m_L`$=0.1$`M_{}`$, is used, as this gives the lowest $`M/L_B`$, hence the highest gas fraction. The scaling is very simple, since the only difference is to scale the lower mass limit from 1 to 0.1 $`M_{}`$. Hence, all luminosities from Leitherer & Heckman are multiplied by 0.39. From their fig. 9 (Z=0.1 $`Z_{}`$), the minimum magnitude is read off to be -16.4 corresponding to a maximum luminosity $`5.6\times 10^8L_{}`$ using the absolute blue solar magnitude 5.48. After IMF scaling this gives $`2.2\times 10^8L_{}`$, so $`M/L_B0.004`$. In table 5 the observed blue luminosity, scaled to h=0.65, is given for each galaxy. Table 5 gives the corresponding gas fraction, as $`\mu =M_{gas}/M_{gas}+M_{stars}`$. If IZw18 is experiencing its first burst, this upper limit should be realistic as its true gas fraction. The large difference from the dynamically deriven gas fraction is caused by a combination of influence from DM, poor total dynamical mass estimation or/and that IZw18 is not experiencing its first burst.
To constrain the models more, three more galaxies have been included, for which we know that they are definitely not experiencing their first burst. One of them is the LMC. Strictly speaking, the LMC is not a dwarf galaxy, but in many ways it behaves as such. For these three galaxies, $`M/L_B`$ is assumed to be 0.5, not as extreme as above, thus obtaining upper limit gas fractions lower than for the other galaxies. Does this assumption still provide an upper limit for the gas fractions? The answer is confirmatory, if the assumed $`M/L_B`$ is a lower limit. The ratio is estimated using the star formation histories of LMC \[1998, fig.7(c)\], SMC , and NGC6822 \[1996, fig.12\]. The general procedure is quite simple. First we find $`L_B`$ of the objects by using a stellar population synthesis model and the SFH. The blue magnitudes in Charlot & Bruzual \[1991, fig.5, dashed line\] are calibrated by using Leitherer & Heckman \[1995, fig.10, 0.25$`Z_{}`$, solid line\]. When we know the age of a starformation epoch, we may find the blue magnitude, hence the luminosity by reading off figure 5 of Charlot & Bruzual. The results are valid for $`m_L`$=1$`M_{}`$, $`m_U`$=100$`M_{}`$ and a starformation rate of 1$`M_{}yr^1`$, since the magnitudes are calibrated by using Leitherer & Heckman. By doing this we ignore the very few stars with masses between 100$`M_{}`$ and 125$`M_{}`$ originally included by Charlot & Bruzual. This will not have any significant effect on the results.
For LMC we have the SFR in relative units. We have splitted the SFH into two star formation epochs, one that started 2 Gyr ago and one 12 Gyr ago. Integration of the SFR and weighting the two read-off luminosities with respect to SFR leads to $`M/L_B`$=0.4. By scaling the luminosities to $`m_L`$=0.1$`M_{}`$ and $`m_L`$=0.01$`M_{}`$, we obtain $`M/L_B`$=1.0 and 2.4, respectively.
For NGC 6822 we are given the SFR in absolute units, but we restrict ourselves to use the relative SFR only. Again we split the SFH into two epochs. A very recent one that started 200 Myr ago and a very old one that started 15 Gyr ago. As we ignore the epoch that stopped 5 Gyr ago, we will obtain a lower estimate of $`M/L_B`$. Using the method outlined above, we get $`M/L_B`$=0.16, and 0.42, 0.96 for the 0.1 and 0.01 lower mass limits, respectively.
The lack of knowledge of the SFH in the SMC is remarkable. Hence, we have used the SFH from the model by Pagel & Tautvaisene , which has succes in explaining e.g. \[Fe/H\]. The SFR is given in absolute units in their table 2, so integration yields the total mass of stars. Using $`L_B=0.99\times 10^9M_{}`$ from table 5, we arrive at $`M/L_B`$=0.43.
In conclusion we have that $`M/L_B`$=0.5 is a conservative lower limit, hence giving an upper limit on the gas fraction.
The results from the calculations are shown for yield set 2a (top plots) and 3 (bottom plots) in fig. 14. In both cases, using the IMF with $`m_L=0.1M_{}`$, the fits are always between the lower and upper limits, except for the three ’moderate upper limit’ galaxies. Lowering the yield by adopting $`m_L=0.01M_{}`$, helps somewhat, but only one of the three moderate galaxies is fitted, namely NGC 6822, the two others being out of range.
Since the model fails in explaining the more restrictively chosen galaxies, the yields have to be lowered further. The number of low-mass stars could be increased, but it would be parameter-gambling, not appealing very much to a physical control of the model. However, it should be noted that Carigi et al. use this possibility to solve the problem. Keeping the IMF, the only way of lowering the yield is to open the model and allow for galactic winds, ordinary or enriched or/and inflow .
### 5.2 The model including enriched winds
From this place forth, the model is opened, allowing for gas exchanges with the intergalactic medium. All simulations have used $`m_L=0.01M_{}`$.
The results of including enriched winds are shown in fig. 15 for set 2a and fig. 16 for set 3. As seen, the wind efficiencies employed are the lowest acceptable for fitting the gas fractions of the Magellanic clouds. Using the same model, the outputs have been plotted with the observations for both N/O-O/H and Y-O/H.
However, removing that much oxygen raises the N/O ratio high above the level of the observations. One could try to reduce N/O by lowering the value of $`\alpha `$, in the case of set 2a, but a value much lower than 1, seems to be quite unrealistic, seen in the light of recent work on stellar evolution, e.g. Marigo .
For set 3 (fig. 16), it is clearly seen that setting the primordial He fraction equal to its lower uncertainty value, as done when using set 3 in the closed model (see fig. 13), is not capable of explaining the He fractions, when introducing these high-efficient enriched winds - the fits are much too poor.
Hence, introducing enriched winds makes it impossible to fit gas fractions and N/O simultaneously. Further, in the case of set 3, problems also arise in fitting the He fractions.
### 5.3 The model including ordinary winds
Though keeping the mass of each burst as a free parameter, it turns out that $`2\times 10^6M_{}`$ is fulfilling the requirement on placing the second burst close to the abundances of IZw18. Thus, it is appropriate to keep this burst mass constant in the following discussion. The results are shown, in fig. 17 for set 2a and fig. 18 for set 3.
It was shown above that the scatter in N/O could be explained by our model. Hence, the scatter will not be given much attention in the following, as the work will be concentrated on getting the right level of N/O, coincident with the gas fraction fitting, if possible.
From the top plots, it is evident that a mass removal of 5 times the burst mass does not appreciably change the appearance of the fits from the closed model. Actually, the O-yield is lowered only by $``$0.3 dex as seen when comparing to fig. 14, using 2a yields. The point may be that we are dealing with a discontinuous star formation. Note for instance the upper linear part of the model curve. The increase in O/H along this part is due to one burst only, as seen when comparing to the corresponding N/O plots. The explanation is that the smaller the absolute gas mass, the larger the effect of oxygen enrichment. Thus, even when adopting higher wind efficiencies, it is not possible to bring the oxygen yield sufficiently down. The same mechanism is responsible for increasing the scatter in N/O for higher wind efficiencies. Hence, it is possible to produce large scatter without having a large burst mass.
Note also that the evolutionary tracks stay close to a gas fraction of 1 for a relatively large part of the evolution. A simple calculation will show, how difficult it is to bring the gas fraction down. For instance, assume the first burst to turn $`2\times 10^6M_{}`$ into stars, and set $`W_{ISM}=5`$. Then, just before the second burst, the gas fraction will be
$`\mu {\displaystyle \frac{M_{gas}}{M_{tot}}}`$ $`=`$ $`{\displaystyle \frac{M_{gas}(0)2\times 10^6M_{}10^7M_{}}{M_{tot}(0)10^7M_{}}}`$
$`=`$ $`0.98,`$
since the initial mass of the dwarf galaxy is $`M_{gas}(0)=M_{tot}(0)=10^8M_{}`$. Using the equation again, one finds $`\mu =0.95`$ just before the third burst. These values are lower limits (!), since the increase in gas from stellar ejecta is disregarded.
Both the level of N/O and slope of Y are satisfied. Note the similarity between the Y plots presented here and those presented by the closed model. This is not surprising, as about half of the helium is produced by massive stars, hence following oxygen. Status is that all observations are matched reasonably well by including ordinary winds, but no better than for the closed model. In particular, it is found that it is impossible to bring the yields sufficiently down to explain the gas fraction intervals of the Magellanic clouds.
### 5.4 Comparing the ordinary wind model with an analytical model
A few calculations have been performed, using a simple analytical model, employing continuous star formation and Hartwick-outflow. The outputs of the analytical model and the numerical model are then compared. The comparison is only performed for yield set 2a, since the inclusion of winds is identical using set no. 3.
The general formulation of Hartwick-outflow is
$`{\displaystyle \frac{dM_{dw}}{dt}}=W_{ISM}{\displaystyle \frac{dM_s}{dt}}`$ (31)
where $`M_{dw}`$ is the mass of the dwarf galaxy and $`M_s`$ is the mass of stars. The next step is to use
$`{\displaystyle \frac{dZ}{dM_s}}={\displaystyle \frac{p}{M_g}}`$ (32)
obtained by considering the changes in Z as a result of stellar ejecta, star formation and outflow, see e.g. Pagel \[1997, his eq. 7.37\]. Instantaneous recycling is assumed, and so is the absence of inflow. Separating the variables and using $`\delta M_g=\delta M_sW_{ISM}\delta M_s`$, one finds
$`{\displaystyle \frac{Z}{p}}={\displaystyle \frac{1}{1+W_{ISM}}}\mathrm{ln}\left({\displaystyle \frac{M_g(Z=0)}{M_g}}\right)`$ (33)
where the yield is assumed constant. The total mass of the dwarf galaxy is written
$`M_{dw}`$ $`=`$ $`M_{dw}(Z=0)W_{ISM}M_s`$ (34)
$`=`$ $`M_g(Z=0)W_{ISM}(M_{dw}M_g)`$ (35)
Isolating $`M_g(Z=0)`$ and inserting it into eq. 33 gives
$`\mu `$ $`=`$ $`{\displaystyle \frac{1+W_{ISM}}{\mathrm{exp}\left(\frac{Z}{p}(1+W_{ISM})\right)+W_{ISM}}}`$
$`=`$ $`\mathrm{exp}\left({\displaystyle \frac{Z}{p}}\right){\displaystyle \frac{1+W_{ISM}}{\mathrm{exp}\left(W_{ISM}\frac{Z}{p}\right)+W_{ISM}\mathrm{exp}\left(\frac{Z}{p}\right)}}`$
where $`\mu `$ is the gas fraction $`\frac{M_g}{M_{dw}}`$. It is seen that for a closed model, $`\mu =\mathrm{exp}\left(\frac{Z}{p}\right)`$. For our numerical model, the yield is metallicity-dependent. Hence, $`\mathrm{exp}\left(\frac{Z}{p}\right)`$, equal to the gas fraction at Z for a characteristic value of p, is found from the closed numerical model. Inserting this into eq. 5.4, finally gives the gas fraction of the analytical model, now employing metallicity-dependent yields and outflow.
The results are shown in fig. 19 for $`W_{ISM}`$=5. The numerical fit (solid line) is the same as in fig. 17. Note that for a large part of the evolution, the two models are coincident, but at the last part they are differing from each other. It is likely that the difference is a result of using instantaneous bursts instead of a continuous SFR. The point may be that the mass of gas is small at the last part of the evolution. Thus, the relative amount of ejected oxygen from one burst, is high compared to the gas mass, hence giving a large increase in O/H.
To make a further check on this hypothesis, a continuous version of the numerical model has been made. This is done in an approximate way, introducing the following changes to the above numerical outflow model:
1. The bursts are all single bursts, i.e. all interburst periods are the same.
2. The interburst periods are made short. The shorter the interburst period, the more continuous the star formation. An interburst period of 3 Myr was sufficient for our purpose.
3. The masses of the bursts are calculated by assuming a mean SFR. The calculations here assume $`<\mathrm{SFR}>`$=0.002 $`M_{}yr^1`$, corresponding to one burst of mass $`2\times 10^6M_{}`$ every Gyr. A reasonable value compared to the calculations above. The masses of the bursts are then equal to 0.002 $`M_{}yr^1`$ 3 Myr = 6000 $`M_{}`$.
All other calculations are performed exactly the same way, as they were, using the double bursting model. The calculations are performed for $`W_{ISM}`$=5 and compared to the analytical model from fig. 19. From fig. 20, it is seen that the two models are almost identical. The small offset may be caused by the calculation of the analytical model, adopting the yields from the closed numerical bursting model. The similarity between the two models confirms the suspicion that instantaneous bursting models behave differently from continuous models, as the numerical model is the same that produced the fit in fig. 19, except for the semi-continuous star formation.
### 5.5 The model including ordinary winds and inflow
The introduction of inflow gives two new parameters: the total mass accreted onto the galaxy $`M_0`$ and the duration of the inflow event $`\tau _{inf}`$. Only results for a few parameter choices are shown and interpreted by close examination of the results.
The gas fractions of the three selected galaxies provide strong restrictions on the possible parameter space, since inflow and outflow balance each other for obtaining the right gas fraction values. Hence, if one increases the rate of inflow, one has to increase the rate of outflow as well. Otherwise, the modelled gas fractions would be too high, and the three selected gas fraction intervals would never be reached.
It is assumed that the mass of the original gas cloud is 0 at t=0. However, inflow of gas increases the mass of the gas cloud, until it starts to form stars at a mass threshold $`M(t^{})`$. The time $`t^{}`$ is calculated using eq. 20, giving
$`t^{}=\tau _{inf}\mathrm{ln}\left({\displaystyle \frac{M_0}{M_0M(t^{})}}\right)`$ (37)
For consistency with previous calculations, the threshold is assumed to be $`10^8M_{}`$. Thus, after a time $`t^{}`$, the first burst appears. As for all the previous models, the first burst is a single one, and all other bursts appear in pairs. Hence, the appearance of the first burst is delayed by the time $`t^{}`$, when comparing to the models not including inflow. This time difference is included, when calculating the age of the system. Remember that the model terminates if the system is older than 15 Gyr.
For yield set 2a, the results are displayed in figs. 21, 22 and figs. 23, 24.
For the parameter values chosen very small scatter results in the N/O diagram in fig. 21. The plots were made using a high inflow/outflow rate to lower the yield sufficiently to fit the gas fractions of NGC 6822 and the Magellanic clouds. The price to pay seems to be a very small scatter in N/O. Hence, in fig. 22 a low rate of inflow has been assumed, in particular at late stages of evolution, as $`\tau _{inf}`$ is decreased to 2 Gyr. Now, it is possible to explain the scatter in N/O, but the result is an unsatisfactory gas fraction fit. Obviously, a large inflow rate of primordial gas buffers the impact of O-rich ejecta on the ISM abundances, resulting in small scatter.
To investigate further, whether it is possible or not to explain the scatter in the N/O-O/H diagram and gas fractions simultaneously, the outflow parameter is kept at 8, but the burst masses are increased to $`8\times 10^6M_{}`$ to give a larger scatter. To balance the high outflow rate (remember that the outflow rate is proportional to the burst mass) a rather high inflow rate is needed, ($`\tau _{inf}`$,$`M_0`$)=(5 Gyr,$`10^9M_{}`$).
The results are shown in figs. 23 and 24. Note that the high outflow/inflow rate forces N/O to increase dramatically with increasing O/H, because inflow of primordial gas decreases O/H, but not N/O. At the same time, the strong outflow decreases the O abundance more than the N abundance, since the composition of the ISM just before each second pair-burst is O enhanced, hence giving an O enhanced wind (though defined and treated as an ordinary wind).
The corresponding gas fraction plot fits NGC 6822 and SMC, but the LMC gas fraction and abundance are never reached. Note the ’sawtooth’ behavior in the gas fraction plot. During the long interburst period, the high inflow rate deposits a large amount of H in the ISM, but no O, hence decreasing O/H until the next burst enriches the ISM again. However, the gas fraction does not increase, because the inflow is balanced by the high outflow rate. From these last results, it seems to be possible to explain both the scatter in N/O and the gas fractions simultaneously, but note that in our struggle to explain the gas fractions of the three selected galaxies, some 3-4 other dwarf galaxies are not fitted within the intervals. Note also that the closed model explained those gas fractions quite well!
The results using yield set 3 are displayed in fig. 25. The gas fractions of the three well-known galaxies are properly fitted, though the same remark as above has to be made, namely that the fits are outside the intervals of 3-5 other dwarf galaxies.
The right level of N/O is obtained. Note in particular the increasing N/O in the left plot, following the trend of the observations. As for set 2a the scatter in N/O is rather small, and insufficient to explain the observed scatter.
The four Y-O/H plots have not yet been discussed, but it may be done in a few words. The model predictions match the observations quite well, but they are of course still showing the difference in slope between set 2a (or 2b) and 3.
## 6 Discussion
Numerical models calculating the chemical evolution of gas-rich dwarf galaxies have been presented. The models have been fitted to a sample of abundance and gas fraction observations. A chemical evolutionary model has to fit these observations simultaneously. It is clear that a realistic model is not just a couple of equations including a lot of parameters that one can change until all data are fitted. One has to remember that the equations are applied to a physical system, obeying physical laws. Hence, before doing any calculations, some considerations were made about ejecta dispersal and mixing processes. It was found that the processes involved are complicated and no complete theory exists yet. However, the four most important features in the enrichment process are found to be stellar ejecta, giant H ii-regions, wind-driven superbubbles and SN-driven supershells. Wind-driven superbubbles arise shortly after a burst due to the strong stellar winds of massive stars. The radii of these superbubbles are in general smaller than the radii of the H ii-regions. As the massive stars explode as SNe, a supershell is swept-up, soon catching up on the superbubble. However, a simple calculation showed that their radii become comparable to the radii of H ii-regions only in the late stages of H ii-region existence. Hence, it is suggested that observed emission lines, used for abundance determinations, arise in the ionized medium outside of the superbubble/supershell. Numerical hydrodynamical models indicate that SN-ejecta always stay within the supershell. If this is true, the observed abundances are not affected by the SN-ejecta from the stars producing the H ii-region, hence being typical for the ISM.
The details of supershell evolution are still not known, but both theory and observations support star formation in expanding shells. From theory, it is expected to start not earlier than about 20 Myr after the starburst, resulting in star formation involving a mass comparable to the mass of the burst that initiated the supershell, hence appearing as a double burst. The double bursting mode of our numerical model ensures the appearance of scatter in the N/O-O/H plane, according to the time-delay idea, namely that N is released some time after O. The time interval between the two bursts of a pair is tuned to give maximum scatter. It is found that this timescale is comparable to the timescale of star formation in an expanding supershell. The requirement for the ’shell burst’ to be the second pair-burst is that the O-rich ejecta mix into the supershell. Because of the poorly understood physics of supershells and mixing processes, this assumption should be seen so far as a working hypothesis. All bursts are assumed to be instantaneous, hence representing short but intense star formation events. The closed model is able to explain the observations of N/O-O/H in both scatter and level. If assuming the upper limit of hot bottom burning to be 5$`M_{}`$, $`\alpha `$=1.1 is used to explain the observations, where $`\alpha `$ is the mixing length parameter. This is a rather low value compared to the results of recent works however , favouring a value close to 2.
Using yield set 3 (the Padova set), it is necessary to extrapolate the primary N yields below Z=0.008 to obtain a higher N yield. Otherwise, the level of N/O becomes too low compared to the observations. The primary N production has to be increasing with decreasing metallicity. Calculations of stellar yields for metallicities lower than Z=0.008 are definitely desired to quantify this..
A very important conclusion is that no primary N production in massive stars is needed to explain the observations. Intermediate mass stars are in position to produce a sufficient amount of primary N.
For the Y-O/H observations, a linear trend is visible, and a linear fit gives $`dY/dZ=2.63\pm 2.21`$ and $`Y_p=0.238\pm 0.004`$. These values are consistent with those of Izotov et al. \[1997a\] within the uncertainties. The closed model is found to explain the Y-O/H observations perfectly, only with different slopes, depending on the yield set in use. In all cases the slopes are within the uncertainties. As seen by inspection of the true He yields for the three sets in table 3, one finds the explanation for the slope difference to be that the He yields of set 3 are 2-3 times higher than those of set 2a or 2b.
It is noteworthy that the Y-O/H relation was fitted using exactly the same parameters as for the N/O-O/H fitting. The problem arises when fitting O/H-$`\mu `$ data. It is argued that the observed gas fractions are actually lower limits, because dark matter is implicitly included in dynamical mass estimates and molecular hydrogen ignored. Hence, it is found useful to calculate upper limits using M/L estimates from starburst evolutionary models. For most objects, very extreme upper limits have to be used, assuming the galaxies to experience maximum luminosity of their first burst, except for three galaxies where the known star formation histories allow us to adopt more moderate and realistic upper limits. If star formation histories are found for a larger sample of dIrrs and even BCGs constraining chemical evolution models with observed gas fraction intervals may eventually turn out to be extremely useful. The closed models do not reproduce the gas fractions of the three well-known objects, even if a lower IMF-cutoff, equal to 0.01 $`M_{}`$ instead of 0.1 $`M_{}`$ is adopted. Hence, open models are considered, allowing gas to escape or to accrete on to the galaxy. Two kind of winds, enriched and ordinary, have been used. The results when incorporating enriched winds are not in accordance with the observed level of N/O and the Y-O/H fitting is not satisfactory, when the gas fractions of the three well-known systems are fitted. Hence, it is concluded that our models employing enriched winds is in conflict with the observations.
The next step is to include ordinary instead of enriched winds. Both N/O-O/H and Y-O/H are fitted, with results resembling those of the closed model. Unfortunately, only the gas fraction of one of the three selected galaxies is fitted, not differing much from the results of the closed model. To check the inclusion of ordinary winds, the outcome of the model is compared to the results of a simple analytical model, employing continuous star formation. A close resemblance is found between the two models at low metallicities, but at higher metallicities, the numerical model seems to have problems in getting the yield down. The difference may be caused by the behavior of starbursts at a low absolute gas mass. This is confirmed when comparing the numerical model, changed slightly to employ continuous star formation, to the analytical model, displaying almost identical outputs. Thus, it is concluded that it is important to specify clearly whether instantaneous bursts or continuous star formation is used, when including ordinary winds. Instantaneous bursts resemble the intense bursts of BCGs, whereas the more moderate bursts of dIrrs are better explained using a continuous SFR.
Finally, inflow and ordinary winds were included. The results are in accordance with all observations, except that it is difficult to obtain the right gas fractions and N/O scatter simultaneously. Only for quite extreme parameter choices as in fig. 23 and 24, one may be succesful. It may be important to note that the results, when using the combined inflow/ordinary wind model, show the upturn in N/O to be more pronounced, than it was for the closed model.
One question is unavoidable: which model is preferred? It is impossible to give an unambiguous answer. Dwarf galaxies are different, both in mass and appearance. Some are explained well by a closed model, others need a combination of ordinary winds and inflow. This is true for both dIrrs an BCGs. However, for dIrrs one should prefer to adopt a more continuous SFR before fitting the observations.
In all cases, the model including enriched winds seems to be ruled out, since it is in direct conflict with the observations, as also found by e.g. Carigi et al. .
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# MAXIMA-1: A Measurement of the Cosmic Microwave Background Anisotropy on angular scales of 10′ to 5°
## 1 Introduction
Measurements of the anisotropy of the cosmic microwave background (CMB) can discriminate between cosmological models and determine cosmological parameters with high accuracy (Kamionkowski & Kosowski, 1999, and references therein). Inflationary dark matter models, for example, predict a series of peaks in the angular power spectrum of the anisotropy. The collected results from many experiments show the existence of a first peak at angular scales corresponding to the spherical harmonic multipole number $`\mathrm{}200`$. These results have been interpreted as evidence for a flat universe (de Bernardis et al., 2000; Lange et al., 2000; Dodelson & Knox, 2000; Tegmark & Zaldarriaga, 2000; Pierpaoli, Scott, & White, 2000). Additional observations probing a broad range of angular scales would greatly increase confidence in these results and further constrain cosmological parameters.
MAXIMA is a balloon-borne experiment optimized to map the CMB anisotropy over hundreds of square degrees with an angular resolution of 10′. In this paper we report results from the MAXIMA-1 flight which took place on August 2, 1998. These include a 100 square degrees map of the CMB anisotropy and the resulting power spectrum over the range $`36\mathrm{}785`$, which is the largest range reported to date. Despite several common team members, the data analysis was independent of that leading to the recently reported BOOMERANG results (de Bernardis et al., 2000). A companion paper, Balbi et al. (2000), discusses the cosmological significance of the MAXIMA-1 results.
## 2 Instrumentation
Lee et al. (1999) gives a detailed description of the MAXIMA system. It is based on a well-baffled, under-filled, off-axis Gregorian telescope with a 1.3 m primary mirror, mounted on an attitude-controlled balloon-borne platform. A well-baffled liquid-Helium-cooled optics box is lined with absorbing material (Bock, 1994) and contains two reimaging mirrors, low-pass filters, field and aperture stops, feed horns for the 16 photometers, and a focal-plane stimulator. Eight conical single-mode horns at 150 GHz and four multi-mode Winston horns each at 240 GHz and 410 GHz provide 10′ beams at all three frequencies. The frequency bands are defined with absorptive and metal-mesh filters. Radiation is detected with spider-web bolometers (Bock et al., 1996) operated at 0.1 K with an adiabatic demagnetization refrigerator (Hagmann & Richards, 1995). The bolometers are AC-biased to avoid low-frequency amplifier noise. Additional channels with a constant resistor, a thermometer, and a dark bolometer are used to monitor electromagnetic interference, cross-talk, and drifts in electronic gain and temperature.
The gondola azimuth is driven by a reaction wheel using information from a two-axis magnetometer and a three-axis rate gyroscope. The telescope elevation is set using information from an angle encoder. Observations were carried out at fixed elevation with the primary mirror scanning $`\pm 2`$° in azimuth at 0.45 Hz and the gondola also scanning in azimuth but at a frequency of $`0.02`$ Hz. Both scans were triangle functions of time with smoothed turnarounds.
## 3 Observations
The MAXIMA-1 flight was launched from the National Scientific Balloon Facility in Palestine, Texas at 1 UT on August 2, 1998. Observations of the CMB dipole for the purpose of calibration began at 3.6 UT when the payload reached an altitude of 32 km and ended at 4.2 UT after $`100`$ rotations at 3.3 rpm. The elevation angle was set to 51°. The payload reached float altitude of $`38.4\pm 0.4`$ km at 4.6 UT.
The 1.6 hour CMB-1 observation began at 4.35 UT with a telescope elevation of 46.3°. The gondola was scanned $`\pm 4.1`$° in azimuth at 16.1 mHz centered at 321.5°. The 1.4 hour CMB-2 observation began at 6.0 UT with a telescope elevation of 32.3°. The gondola was scanned $`\pm 2.9`$° at 21.3 mHz centered at 323°. Because of sky rotation, the combination of these observations covered a nearly square region of the sky with an area of 124 square degrees of which 45% is cross-linked at an angle of $``$22°.
Observations of Jupiter were carried out from 7.5 to 8.1 UT to map the telescope beams and provide additional calibration information. The elevation was fixed at 44.2° while sky rotation and the primary mirror modulation provided $`200`$ transits across each beam.
## 4 Pointing Reconstruction and Calibration
We identified the stars which moved through the field of a CCD camera aligned with the center of the primary mirror scan by using the balloon location, telescope elevation, and the position of Polaris in an offset CCD camera. Interpolations using an angle encoder on the primary mirror, rate gyroscopes, and the known star positions permitted pointing reconstruction to better than 1′ RMS. Less than 0.1% of the data had pointing uncertainty larger than 2′ and were not used.
A full beam calibration of the 150 and 240 GHz photometers was obtained from observations of the CMB dipole. The data from each rotation were $`\chi ^2`$-fitted to a linear combination of a dipole model (Lineweaver et al., 1996), a galactic-dust emission model (Schlegel, Finkbeiner, & Davis, 1998), data from one 410 GHz photometer, an offset, and a gradient. The amplitude of each of these components was treated as a free parameter. A monotonic change in the detector calibration of less than 9% throughout the CMB observations, due to an increasing detector temperature, was monitored by illuminating the focal plane with the stimulator lamp. Estimated $`1\sigma `$ calibration uncertainties were less than 4% for each of the 150 and 240 GHz photometers. The uncertainties in the dipole calibration and the time dependent calibration contributed about equally to the total error, and systematic sources contributed about 25%.
Beam maps and an independent calibration were obtained from observations of Jupiter. We estimate a $`1\sigma `$ uncertainty of $`\pm 0.5`$′ in the size of the beams. The beam profiles were integrated and used with the angular diameter and brightness temperature of Jupiter (Goldin et al., 1997) and the optical bandpass functions to calibrate all 16 photometers. For the data reported here, the errors in the calibration from Jupiter were between 12 and 14% (with about 10, 5, 5, and 2% coming from uncertainties in the beam solid angle, frequency bands, Jupiter’s flux, and measured signal, respectively). The absolute calibrations from the dipole agreed with those from Jupiter to within $`1\sigma `$, with the Jupiter calibration predicting larger temperature fluctuations by 11 to 14%.
## 5 Map and Angular Power Spectrum
In this paper we report on the analysis of data from the three 150 GHz, one 240 GHz, and one 410 GHz photometer described in Table 1. At 150 and 240 GHz these photometers are the most sensitive in the MAXIMA array and give the highest sensitivity of any CMB instrument reported to date.
The raw data for each photometer consisted of 2.3 million samples of which about 16% were not used. We removed the stimulator calibration events and other events with an amplitude larger than $`6\sigma `$. This procedure broke the data into 20 segments that were treated as independent observations of the sky. Samples in each of the segments which were in excess of $`4\sigma `$, such as cosmic ray hits and short telemetry drop-outs, were removed. For the data of one of the photometers, we repeated the data analysis by using a threshold of $`3\sigma `$ with no significant change in the resulting angular power spectrum. We deconvolved the transfer functions of the bolometers and readout electronics and estimated the noise power spectrum from sections of the time stream that had no gaps (Stompor et al., in preparation). We used the procedure of Ferreira & Jaffe (2000) to confirm that the time-domain data are dominated by noise. We marginalized over frequencies lower than 0.1 Hz and higher than 30 Hz, where we did not expect appreciable optical signals.
The calibrated time stream data were combined with the pointing solution to produce a maximum likelihood pixelized map of temperature anisotropy and a pixel-pixel noise correlation matrix for each photometer (Wright, 1996; Tegmark, 1997; Bond et al., 1999). An area of $`20`$ square degrees of the map was not well cross-linked and is not included in the present analysis. The data showed a signal that was phase-synchronous with the primary mirror scan which was due to radio-frequency interference from on-board transmitters modulated by the mirror motion. This signal was constant within each data segment, varied between different photometers, and had an amplitude of 100 - 300 $`\mu \text{K}`$. We removed it by allocating fictitious map pixels to values of the primary mirror angle and determined the maximum likelihood map in these pixels simultaneously with the temperature anisotropy map (Stompor et al., in preparation).
We verified that there are no noise correlations between maps of different photometers by producing difference maps from the data of pairs of photometers. The angular power spectra of these maps were consistent with no signal. We also showed that histograms of the temperatures in the difference maps were consistent with the distributions expected for no noise correlations at a Kolmogorov-Smirnov significance level larger than 10%. A combined temperature anisotropy map was then produced by adding individual maps with a weight inversely proportional to their noise correlation matrices. A Wiener filtered version of this map is shown in Figure 1. We assign a calibration uncertainty of 4% to the magnitude of temperature fluctuations in the combined map.
We calculated the angular power spectrum $`C_{\mathrm{}}`$ of the combined map using the MADCAP (Borrill, 1999) implementation of the maximization of the likelihood function following Bond, Jaffe & Knox (1998). This implementation assumes that the beam shape has axial symmetry. We produced an effective beam for the analysis of the combined map by noise-weight averaging the individual beams. We followed the procedure of Wu et al. (2000) to find a symmetric approximation for the effective beam and included the small smoothing provided by the pixelization. We tested this procedure for the MAXIMA-1 beams and data and found no systematic bias of the $`C_{\mathrm{}}`$ estimates (Wu et al., 2000) . We calculated the power spectrum of the temperature fluctuations using 15 bins in $`\mathrm{}`$ over the range $`3\mathrm{}1500`$ assuming a constant $`\mathrm{}(\mathrm{}+1)C_{\mathrm{}}`$ band power in each bin, and marginalizing over the CMB monopole and dipole. We further marginalized over the bins at $`\mathrm{}<35`$ and $`\mathrm{}>785`$ and diagonalized the $`\mathrm{}`$-bin correlation matrix using a variant of techniques discussed in Bond, Jaffe & Knox (1998). The correlations between the dominant bin and adjacent bins were typically less than 10%. Table 2 lists the dominant bins, the $`C_{\mathrm{}}`$ estimates, and the $`\mathrm{\Delta }T=\sqrt{\mathrm{}(\mathrm{}+1)C_{\mathrm{}}/2\pi }`$ estimates for the corresponding uncorrelated linear combinations of bins. We quote $`1\sigma `$ errors on the $`C_{\mathrm{}}`$ estimates assuming 68% confidence intervals using the offset log-normal distribution model of Bond, Jaffe & Knox (2000). These errors do not include two additional independent sources of uncertainty. Expressed as a percentage of $`\mathrm{}(\mathrm{}+1)C_{\mathrm{}}/2\pi `$, the $`1\sigma `$ calibration error is a constant 8% and the $`\mathrm{}`$-dependent error due to the beam-diameter uncertainty, which has been shown to be fully correlated between $`\mathrm{}`$ bins, is given in Table 2. Information on the shape of the bin-power likelihood functions and window functions will be made available on the MAXIMA web site (http://cfpa.berkeley.edu/maxima).
The top panel of Figure 2 shows the maximum likelihood power spectrum, an inflationary adiabatic model that best fits the MAXIMA-1 and COBE/DMR power spectra, and a $`\mathrm{\Lambda }`$CDM model. The $`\chi ^2`$ for the best fit and for the $`\mathrm{\Lambda }`$CDM models are 41 and 50, respectively, using all 38 data points. If we only use the 10 data points of MAXIMA-1 we obtain $`\chi ^2=7`$ and 9 for the best fit and $`\mathrm{\Lambda }`$CDM, respectively (Balbi et al., 2000).
## 6 Foregrounds
Foreground sources include emission from the earth, the atmosphere, interstellar dust, free-free radiation, synchrotron radiation, and point sources, and scattering due to the Sunyaev-Zeldovich effect. The two CMB scans, performed with a separation of 1.5 hours at different telescope elevations, show consistent structure. This temporal stability is inconsistent with an atmospheric or ground-based origin for the signal. We extrapolated the 100 µm, 10′ resolution Schlegel, Finkbeiner, & Davis (1998) map of our observed region to lower frequencies using the two-component dust model favored by Finkbeiner, Davis & Schlegel (1999). The predicted RMS dust temperature is only 2.3 and 9.3 $`\mu \text{K}`$, for the 150 and 240 GHz bands, respectively. We found a statistically significant correlation between the dust model and our maps. The ratio of the detected to predicted RMS dust signal was statistically consistent with unity. When we subtracted the correlated dust signal from the maps the change in the measured angular power spectrum was negligible. A catalog search (Sokasian, Gawiser, & Smoot, 2000; Gawiser & Smoot, 1997) yielded no detectable radio or infra-red sources in any of the frequency bands. Estimates of bremsstrahlung and synchrotron radiation (Bouchet & Gispert, 1999) yielded contributions of less than 1 $`\mu \text{K}`$ at 150 and 240 GHz and no subtractions were made.
Integrating the measured power spectrum of one photometer at 150 GHz and the one at 240 GHz we find for a thermodynamic temperature fluctuations ratio of $`0.91\pm 0.18`$ compared to unity for the CMB, 0.06 for emission from dust, 2.1 for synchrotron, and 3.7 for free-free radiation, assuming the Tegmark et al. (2000) “middle” foregrounds model. For the ratio between 150 and 410 GHz, we found a lower limit of $`1.110^2`$ compared to $`3.210^4`$ for dust and $`5.710^5`$ for atmosphere. The observed anisotropy is not consistent with the thermal SZ effect, which would produce anti-correlated structure in the 150 and 240 GHz bands.
## 7 Tests for Systematic Errors
Because of computational limitations, tests for systematic errors were done with maps which had square pixels of 8′ and 10′ on a side. The power spectra calculated from the 8′-pixel and 5′-pixel combined maps were statistically consistent. The following combinations of the data were analyzed and produced a power spectrum consistent with no signal: (1) a dark bolometer, (2) the data from the 410 GHz photometer, (3) the difference between the overlapping part of the combined map from the CMB-1 and -2 scans, (4) the differences between the maps produced by different photometers. We also weight-averaged the maps of the second and third photometers in Table 1 and the first and fourth. The maximum likelihood estimate of the power spectrum of the difference between these independent maps is consistent with no signal as shown by the open circles in the top panel of Figure 2.
We compared the estimate of the angular power spectrum to that obtained using: (1) only the sections of the map where the CMB-1 and -2 scans overlap, (2) a map of each of the CMB scans alone. We made maps and calculated $`C_{\mathrm{}}`$ estimates from the data of each photometer alone and using: (1) only a sub-section of the time stream data, (2) a high-pass filtered version of the time stream where the high-pass was a time domain box-car with a width of 10 sec, (3) various combinations of frequency marginalizations between 30 and 70 Hz, and 0.05 and 0.3 Hz, respectively. In all these cases the computed power spectra agreed among themselves and with the power spectrum presented in this paper. The Kolmogorov-Smirnov test (see Section 5), as applied to difference maps and to a map of the dark bolometer, confirms that we correctly estimated the noise in the experiment and that the pixel-domain noise is Gaussian.
We used simulations to test the algorithm used to subtract the signal that was phase-synchronous with the primary mirror modulation. We found that the power spectrum estimate was not biased with phase-synchronous signals larger than those observed in the data. If we make the maps without removing the phase-synchronous signal the power spectrum estimate changes only at $`\mathrm{}200`$.
The computer programs used to generate the maps and power spectra were tested extensively using simulations of the time domain data and noise. Maps were produced by two independent computer codes and the power spectra calculated from these maps were consistent. We used one of the map-making codes to make maps of Jupiter and found them consistent with those obtained with a simple data-binning technique.
## 8 Discussion
We have observed temperature anisotropy on the sky at 150 and 240 GHz that is consistent with fluctuations in the cosmic microwave background radiation, and inconsistent with any known foreground. The observations were carried out with photometers that give the highest CMB sensitivity reported to date. Our measurements cover a range of angular scales corresponding to the multipole range $`36\mathrm{}785`$, which is the largest yet reported by a single experiment. The measured angular power spectrum shows a clear peak at $`\mathrm{}220`$, and an amplitude varying between $`40`$ $`\mu \text{K}`$ and $`50`$ $`\mu \text{K}`$ at $`400\mathrm{}785`$. The power spectrum is well fit by an inflationary adiabatic model over the entire range of $`\mathrm{}`$. The best-fit model has a total energy density close to unity and a non-zero cosmological constant. The MAXIMA-1 power spectrum appears consistent with that of the BOOMERANG experiment (de Bernardis et al., 2000) once the power spectra of the two experiments are scaled by factors equal to their respective $`1\sigma `$ calibration uncertainties, see Figure 2. A detailed analysis of the combined data sets which includes a determination of the calibration factors that bring the experiments to agreement is presented in Jaffe et al. (2000).
We thank Danny Ball and the other staff at NASA’s National Scientific Balloon Facility in Palestine, TX for their outstanding support of the MAXIMA program. MAXIMA is supported by NASA Grants NAG5-3941, NAG5-6552, NAG5-4454, GSRP-031, and GSRP-032, and by the NSF through the Center for Particle Astrophysics at UC Berkeley, NSF cooperative agreement AST-9120005, and a KDI grant 9872979. The data analysis used resources of the National Energy Research Scientific Computing center which is supported by the Office of Science of the U.S. Department of Energy under contract no. DE-AC03-76SF00098. PA acknowledges support from PPARC rolling grant, UK.
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# K-Theory Torsion
## 1 Introduction
It has been shown that K-theory classifies the topological charge of the D-brane gauge bundle (or the associated vector bundle). The crucial observation for this was that adding a brane-antibrane pair with the same gauge bundle does not change the total charge. Or in other words, you may add a D-p brane with the trivial bundle, then try to straighten out any “windings”, in the bigger bundle, bring the bundle back in the form where the trivial bundle is one summand, and then remove the brane you added.
Mathematically, this is called stabilization, and the charge of the gauge bundle are stable isomorphism classes. It is not difficult to find examples of vector bundles that are stable isomorphic but not isomorphic, for example $`TS^2`$ and the trivial bundle $`S^2\times ^2`$ (as real vector bundles).
Another way to look at K-theory is that it is a generalized cohomology theory, that is it satisfies all the usual axioms except that higher cohomology groups of a point may not vanish. Since we can express usual field theory in terms of differential forms and de Rahm cohomology, it seems natural that a generalization of field theory leads to a generalized cohomology theory.
Now for all well-behaved spaces $`X`$ (such as topological manifolds or finite CW complexes), $`K(X)`$ is a finitely generated abelian group, i.e. of the form $`^n\text{Torsion}`$. Interestingly, torsion charges can appear. In ordinary field theory you could also have torsion in integral cohomology $`H^{}(X,)`$, but physical fields must be represented by differential forms, and this prohibits torsion. But on the K-theory side torsion charges are apparently physical charges. The purpose of this paper is to better understand the relation between integral cohomology and K-theory.
For concreteness, I will restrict myself in the following to IIB string theory with spacetime manifold $`X`$, where the possible D-brane charges are $`K\mathbf{(}X\mathbf{)}`$, the Grothendiek group of complex vector bundles.
Of course we want $`X`$ to be a $`10`$–manifold (where Poincaré duality holds). In the following I will investigate compact topological manifolds of lower dimension which exhibit torsion in cohomology. The physical motivation for this is a spacetime of the form $`M_d\times ^{1,9d}`$, which on the K-theory side is just the $`(10d)`$-th suspension of $`M_d`$. So the torsion of K-theory comes purely from the compact dimensions, and not from the Minkowski part of spacetime.
During the preparation of this paper another work appeared that also discusses the implications of the Atiyah–Hirzebruch spectral sequence .
I would like to thank Philip Candelas, Thomas Friedrich, Albrecht Klemm, Dieter Lüst, André Miemiec, and Ulrike Tillmann.
Here is a quick outline of the following sections:
1. This introduction
2. The Chern isomorphism and the Atiyah–Hirzebruch spectral sequence, which are the main tools used in this paper, are introduced. To demonstrate their utility I prove that the order of the torsion part is the same in integer cohomology and K-theory for odd $`\mathrm{P}^n`$. The sequence provides a necessary criterion for K-theory torsion.
3. Given some K-theory element, I determine the dimensionality of the corresponding D-brane (That is the minimum dimension needed to carry the charge). This can also be calculated from the Atiyah–Hirzebruch spectral sequence.
4. An example with different order of the torsion part in integer cohomology an K-theory is analyzed in detail.
5. By considering line bundles I find a sufficient criterion for K-theory torsion. This makes it possible to give an example for a Calabi-Yau manifold with torsion.
6. Conclusion
## 2 From $`H^{}(X,)`$ to $`K^{}(X)`$
The most important result is the Chern isomorphism:
$`K^0(X)_{}`$ $``$ $`{\displaystyle H^{2i}(X,)}\stackrel{\mathrm{def}}{=}H^{\mathrm{ev}}(X,)`$ (1)
$`K^1(X)_{}`$ $``$ $`{\displaystyle H^{2i+1}(X,)}\stackrel{\mathrm{def}}{=}H^{\mathrm{odd}}(X,)`$
which is induced by the Chern character $`ch:K(X)H^{\mathrm{ev}}(X,)`$.
This means that we can compute the free part of K-theory directly from ordinary de Rahm cohomology. Or in physical language, K-theory without torsion is just a reformulation of what one already knows from calculations on the level of differential forms. On the other hand side the torsion part of $`H^{\mathrm{ev}}(X,)`$ and $`K(X)`$ do in general differ, for example<sup>2</sup><sup>2</sup>2$`\mathrm{P}^5`$ is not spin, and therefore not a phenomenologically viable background spacetime. $`\mathrm{P}^7`$ would be a counterexample that is spin. Remember the relevant Stiefel-Whitney classes $`w_1(\mathrm{P}^n)=n+1`$, $`w_2(\mathrm{P}^n)=\frac{n(n+1)}{2}`$ (mod 2)
$`K(\mathrm{P}^5)`$ $`=`$ $`_4`$ (2)
$`H^{\mathrm{ev}}(\mathrm{P}^5,)`$ $`=`$ $`_2_2`$
It has been noted that — although there is no surjective group homomorphism — the order of the torsion part is equal. Unfortunately, this is caused by peculiarities in the cohomology of real projective spaces and not a generic feature. A counterexample will be presented in section 4.
For now, lets use the Atiyah–Hirzebruch spectral sequence to understand why the order is indeed equal for $`\mathrm{P}^n`$, with $`n`$ odd so that the manifold is orientable. This spectral sequence stems from the filtration of the space by its CW-skeleton. It has the second term
$$E_2^{p,q}=\{\begin{array}{cc}H^p(\mathrm{P}^n,)& q\text{even}\hfill \\ 0& q\text{odd}\hfill \end{array}$$
(3)
and converges towards the associated graded complex of $`K(\mathrm{P}^n)`$.
For simplicity, take $`n=5`$:
$$E_2^{p,q}=\begin{array}{ccccccc}& 0& 0& 0& 0& 0& 0\\ q& & 0& _2& 0& _2& \\ & 0& 0& 0& 0& 0& 0\\ & & 0& _2& 0& _2& \\ & 0& 0& 0& 0& 0& 0\\ & & 0& _2& 0& _2& \\ & & & & & & \\ \multicolumn{7}{c}{p}\end{array}$$
(4)
The differential $`d_2^{p,q}:E_2^{p,q}E_2^{p+2,q1}`$ either has domain or range $`0`$, thus
$$E_3^{p,q}=kerd_2^{p,q}/imgd_2^{p2,q+1}=E_2^{p,q}$$
(5)
In this sequence the $`d_{\mathrm{even}}`$ are obviously irrelevant, and $`E_{2k}=E_{2k+1}`$.
The only<sup>3</sup><sup>3</sup>3Of course the table is cyclic of order 2 in $`q`$ $`d_3^{p,q}:E_3^{p,q}E_3^{p+3,q2}`$ with nonvanishing domain and range is<sup>4</sup><sup>4</sup>4In it is noted without proof that $`d_3=Sq^3`$, the third Steenrod Square $`d_3^{2,2}:_2`$. Since there is no nonzero group homomorphism from $`_2`$ to $``$, $`d_3=0`$.
So far we found $`E_5=E_2`$, and again there is only one $`d_5`$ with nonvanishing domain and range, $`d_5^{0,4}:`$. But the Chern isomorphism tells us that after tensoring everything with $``$ the spectral sequence already degenerates at level 2. Thus $`d_5_{}=0`$, and since domain is torsion free this implies $`d_5=0`$ (This also proves that torsion in cohomology is necessary for torsion in K-theory).
Thus $`E_{\mathrm{}}=E_2`$, but this is not enough to compute $`K(\mathrm{P}^5)`$. All that it tells us is that there is a filtration
$`K^0(\mathrm{P}^5)`$ $`=`$ $`F_6^0F_5^0F_4^0F_3^0F_2^0F_1^00`$ (6)
$`K^1(\mathrm{P}^5)`$ $`=`$ $`F_6^1F_5^1F_4^1F_3^1F_2^1F_1^10`$
such that the successive cosets are the even respectively odd diagonals of $`E_{\mathrm{}}`$:
$$\begin{array}{cccccc}F_6^0/F_5^0=& F_5^0/F_4^0=0& F_4^0/F_3^0=_2& F_3^0/F_2^0=0& F_2^0/F_1^0=_2& F_1^0/0=0\\ F_6^1/F_5^1=0& F_5^1/F_4^1=0& F_4^1/F_3^1=0& F_3^1/F_2^1=0& F_2^1/F_1^1=0& F_1^1/0=\end{array}$$
(7)
Obviously $`K^1(\mathrm{P}^5)=F_6^1=F_5^1=F_4^1=F_3^1=F_2^1=F_1^1=`$. But for $`K^0(\mathrm{P}^5)`$ we find $`F_1^0=0`$, $`F_3^0=F_2^0=_2`$ and then hit the extension problem: either $`_4/_2=_2`$ or $`(_2_2)/_2=_2`$. So $`F_5^0=F_4^0=_4\text{or}_2_2`$. Since $`F_5^0`$ is pure torsion each possibility determines a unique $`K(\mathrm{P}^5)=F_6^0`$, either $`_4`$ or $`_2_2`$.
However, the ambiguity is between groups of equal order (since the ambiguous extension was between finite abelian groups), and moreover one of the possibilities was $`H^{\mathrm{ev}}(\mathrm{P}^5)`$, since the spectral sequence already degenerates at $`E_2`$.
The same argument can be used for all real projective spaces to prove that the order of the torsion part of cohomology and K-theory are equal, but as we have seen the proof depends on the special properties of $`\mathrm{P}^n`$.
## 3 Dimension of D-branes
### 3.1 Filtering $`K(X)`$
In flat space one can explicitly construct vector bundles that carry a nontrivial topological charge (using the Clifford algebra, see and ). The bundle is trivial everywhere except on a hyperplane of even codimension, which is identified with the D-brane. One can extend this construction to general submanifolds with $`\mathrm{Spin}_\mathrm{c}`$ normal bundle. This fits nicely to the fact that D-branes in IIB are even dimensional.
But to understand what the charges are one should rather understand which submanifold can carry a given K-theory element. The intuitive answer would be: An arbitrary submanifold $`YX`$ can carry the charge $`x=[E][F]K(X)`$ if there exists an isomorphism $`E|_{XY}F|_{XY}`$. Of course this is not well-defined, since the same K-theory element could be represented by different vector bundles $`E^{}`$, $`F^{}`$ that are stably isomorphic but not isomorphic. So we should really ask whether there exists an isomorphism $`(E|_{XY}^k)(F|_{XY}^k)`$ for some $`k`$.
One would like to use the inclusion map $`i:XYX`$ to pull back $`x`$, and thus automatically include stabilization as an element of $`K(XY)`$, but unfortunately in general $`i^{}(E)i^{}(F)K(XY)`$ since the complement $`XY`$ is not compact (Remember that K-theory on noncompact spaces are differences of vector bundles that are isomorphic outside a compact subset).
So instead take compact submanifolds $`j:ZX`$ as probes: If their dimension is too low, they will generically miss the D-brane and the pullback $`j^{}(x)=0K(Z)`$. Since $`j^{}`$ depends only on the homotopy class of $`j`$, we do not have to worry about degenerate cases. If we cannot detect $`x`$ with submanifolds of a given dimension $`p`$ then we conclude that $`x`$ is carried by a D-brane of codimension greater than $`p`$.
But the total charge $`0K(X)`$ could also be carried by a brane-antibrane pair that is separated in spacetime. Probing only in the neighborhood of one brane one would falsely find a charge. So our probe submanifold must somehow be big enough. Discussing this in terms of submanifolds is very cumbersome, so instead think of spacetime $`X`$ as a cell complex (simplicial complex or CW complex). Then take the $`p`$-skeleton $`X^p`$ as probe; it can easily be seen that this is independent of the chosen cell structure. Any cell complex embedded in $`X`$ is a subcomplex for some cell structure on $`X`$, in that sense $`X^p`$ probes the whole space.
Let $`K_p(X)`$ be the subgroup of $`K(X)`$ of charges that live on a brane<sup>5</sup><sup>5</sup>5More precisely a stack of coincident branes, although I will not make that distinction in the following of codimension $`p`$ or higher, that is D-$`(dim(X)p1)`$-branes or lower. According to the previous arguments
$$K_p(X)=ker\left(K(X)K(X^{p1})\right)$$
(8)
where the map is the one induced by the inclusion $`X^{p1}X`$.
This yields a filtration
$$K(X)=K_0(X)\stackrel{~}{K}(X)=K_1(X)K_2(X)\mathrm{}K_{dimX+1}(X)=0$$
(9)
where the successive quotients $`K_p(X)/K_{p+1}(X)`$ are the D-$`(dim(X)p1)`$-brane charges.
### 3.2 Remarks
Lets try to understand eq. (8) better. $`K_{dimX+1}(X)=0`$ means that there are no D-$`(2)`$-branes or less, which is correct.
The $`9`$-brane charges ($`p=0`$) are $`K(X)/\stackrel{~}{K}(X)=`$, which is the virtual rank of the bundle pair. This we can also understand: If we do not start with the same number of $`9`$ and $`\overline{9}`$-branes, then there will always be a $`10`$-dimensional brane left. On the other hand side if the virtual rank is $`0`$ (as required by tadpole cancellation), then the vector bundles are isomorphic over sufficiently small open sets (since they are locally trivial), which one could use to localize the nontrivial windings at a subspace of codimension $`1`$.
Fortunately there is a way to calculate the quotients $`K_p(X)/K_{p+1}(X)`$. First note that one can extend eq. (8) to the higher $`K`$-groups straightforwardly:
$$K_p^n(X)=ker\left(K^n(X)K^n(X^{p1})\right)$$
(10)
And the associated graded complex to this filtration is precisely the limit of the Atiyah–Hirzebruch spectral sequence
$$E_{\mathrm{}}^{p,q}=K_p^{p+q}(X)/K_{p+1}^{p+q}(X)$$
(11)
If there is no torsion in integer cohomology then the spectral sequence degenerates at level 2, and (compare eq. (3))
$$K_p^{}(X)/K_{p+1}^{}(X)H^p(X,)$$
(12)
where the isomorphism is just the Chern character. This confirms the interpretation of the dimensionality of the K-theory elements.
The odd rows in $`E_{\mathrm{}}^{p,q}`$ all vanish, so for odd $`p`$ and odd $`q`$
$$K_p^{p+q}(X)/K_{p+1}^{p+q}(X)=K_p(X)/K_{p+1}(X)=0$$
(13)
which just means that there are no topological charges for odd dimensional D-branes.
A word of caution: even if $`K(X)`$ is torsion-free, one of the successive quotients can be torsion, as in the example $`/2=_2`$. Physically, this means that there can be an apparent torsion charge on a D-brane in the sense that multiple copies of that brane can decay to something lower-dimensional, which a single brane cannot. But the lower-dimensional remnant then carries an ordinary (non-torsion) charge that keeps track of the number of branes we started with.
## 4 Examples
### 4.1 $`K(\mathrm{P}^n)`$
The best way to construct manifolds with K-torsion is to use quotients of well-understood manifolds (like the sphere) by free group actions. At the example $`\mathrm{P}^n`$ I will review the necessary tools (See e.g. ).
Let $`_2`$ act on $`S^n`$ ($`n`$ odd) via the antipodal map, a free group action. In general (for free group actions) K-theory on the quotient is equal to equivariant K-theory on the covering space $`K^{}(S^n/_2)=K__2^{}(S^n)`$. Writing down the (cyclic) long exact sequence associated to the inclusion $`S^nD^{n+1}`$, we find:
$$\begin{array}{ccccc}K__2^1(S^n)& & K__2^1(D^{n+1})& & K__2^1(D^{n+1},S^n)\\ & & & & \\ K__2^0(D^{n+1},S^n)& & K__2^0(D^{n+1})& & K__2^0(S^n)\end{array}$$
(14)
where the $`_2`$ action on the disk $`D^{n+1}`$ is the obvious extension of the $`_2`$–action on $`S^n`$.
Now identify $`K__2^0(D^{n+1},S^n)`$, virtual differences of vector bundles on $`D^{n+1}`$ that are isomorphic over the boundary, with $`K__2^0(^{n+1})`$, virtual differences on $`^{n+1}`$ with isomorphism outside a compact subset. The associated $`_2`$–action on $`^{n+1}`$ is again $`xx`$. Since $`n+1`$ is even, we can interpret $`^{n+1}=^{(n+1)/2}\stackrel{\mathrm{def}}{=}^m`$ with a linear $`_2`$–action on $`^m`$. And this is a $`_2`$–equivariant vector bundle over a point.
Then use the Thom isomorphism, that is $`K_G(E)=K_G(X)`$ for any $`G`$–vector bundle $`E`$ over $`X`$ (as abelian groups, the multiplication law is different):
$$K__2^0(D^{n+1},S^n)=K__2^0(^m)=K__2^0(\{\text{pt}\})=R(_2)=[x]/x^21$$
(15)
$`R(_2)`$ are the formal differences of representations of $`_2`$ (with the obvious ring structure induced by the tensor product of representations), and $`x`$ denotes the unique nontrivial irreducible representation of $`_2`$. If one is only interested in the underlying abelian group, this is of course $``$.
Doing the same for $`K^1`$ and using the homotopy $`D^{n+1}\{pt\}`$, we evaluate eq. (14):
$$\begin{array}{ccccc}K__2^1(S^n)& & 0& & 0\\ & & & & \\ [x]/x^21& \stackrel{f}{}& [x]/x^21& \stackrel{g}{}& K__2^0(S^n)\end{array}$$
(16)
Since $`[x]/x^21`$ is torsion free as abelian group, so must be $`K__2^1(S^n)=K^1(\mathrm{P}^n)`$. From the Chern isomorphism then follows that $`K^1(\mathrm{P}^n)=`$. But to determine the torsion part of $`K^0(\mathrm{P}^n)`$, we need to identify the map $`f`$. Tracing everything back to the Thom isomorphism, one can show that $`f`$ is multiplication with $`(x1)^m`$. Using exactness ($`imgf=kerg`$) we find
$`K^0(\mathrm{P}^n)`$ $`=`$ $`K__2^0(S^n)=[x]/x^21,(x1)^m=`$ (17)
$`=`$ $`[z]/(z+1)^21,z^m=[z]/z^2+2z,z^m`$
Up to the given relations, each ring element can be represented as $`az+b`$, $`a,b`$. While $`b`$ is not subject to any relation, we can use $`z^2+2z=0`$ and $`z^m=0`$ to show $`2^{m1}z=0`$. Therefore (ignoring the ring structure):
$`K^1(\mathrm{P}^n)`$ $`=`$ $``$
$`K^0(\mathrm{P}^n)`$ $`=`$ $`_{2^{m1}}`$ (18)
### 4.2 $`K(\mathrm{P}^3\times \mathrm{P}^5)`$
Here is the promised example of a space where the order of the torsion subgroup in K-theory and ordinary cohomology differs.
Of course we use the Künneth formula to calculate the cohomology<sup>6</sup><sup>6</sup>6In this section, $`H^{}(X)`$ is always cohomology with integer coefficients of a Cartesian product:
$$0\underset{i+j=m}{}H^i(X)H^j(Y)H^m(X\times Y)\underset{i+j=m+1}{}Tor(H^i(X),H^j(Y))0$$
(19)
The cohomology of real projective space is
$$H^i(\mathrm{P}^3)=\{\begin{array}{cc}& i=3\hfill \\ _2& i=2\hfill \\ 0& i=1\hfill \\ & i=0\hfill \end{array}H^i(\mathrm{P}^5)=\{\begin{array}{cc}& i=5\hfill \\ _2& i=4\hfill \\ 0& i=3\hfill \\ _2& i=2\hfill \\ 0& i=1\hfill \\ & i=0\hfill \end{array}$$
(20)
Thus eq. (19) contains the exact sequences
$$\begin{array}{ccccccccc}0& & & & H^8(\mathrm{P}^3\times \mathrm{P}^5)& & 0& & 0\\ 0& & _2_2& & H^7(\mathrm{P}^3\times \mathrm{P}^5)& & 0& & 0\\ 0& & _2& & H^6(\mathrm{P}^3\times \mathrm{P}^5)& & 0& & 0\\ 0& & _2& & H^5(\mathrm{P}^3\times \mathrm{P}^5)& & _2& & 0\\ 0& & _2_2& & H^4(\mathrm{P}^3\times \mathrm{P}^5)& & 0& & 0\\ 0& & & & H^3(\mathrm{P}^3\times \mathrm{P}^5)& & _2& & 0\\ 0& & _2_2& & H^2(\mathrm{P}^3\times \mathrm{P}^5)& & 0& & 0\\ 0& & 0& & H^1(\mathrm{P}^3\times \mathrm{P}^5)& & 0& & 0\\ 0& & & & H^0(\mathrm{P}^3\times \mathrm{P}^5)& & 0& & 0\end{array}$$
(21)
Using Poincaré duality ($`\mathrm{P}^3\times \mathrm{P}^5`$ is an orientable 8-manifold since each factor is), $`H_{\mathrm{tors}}^5H_{\mathrm{tors}}^4`$ and $`H_{\mathrm{tors}}^3H_{\mathrm{tors}}^6`$. This fixes the extension ambiguities, and we find
$$H^i(\mathrm{P}^3\times \mathrm{P}^5)=\{\begin{array}{cc}& i=8\hfill \\ _2_2& i=7\hfill \\ _2& i=6\hfill \\ _2_2& i=5\hfill \\ _2_2& i=4\hfill \\ _2& i=3\hfill \\ _2_2& i=2\hfill \\ 0& i=1\hfill \\ & i=0\hfill \end{array}\{\begin{array}{c}H^{\mathrm{ev}}(\mathrm{P}^3\times \mathrm{P}^5)=^2_2^5\\ H^{\mathrm{odd}}(\mathrm{P}^3\times \mathrm{P}^5)=^2_2^5\end{array}$$
(22)
For K-theory there is the following analog to the ordinary Künneth formula:
$$0\underset{i+j=m}{}K^i(X)K^j(Y)K^m(X\times Y)\underset{i+j=m+1}{}Tor(K^i(X),K^j(Y))0$$
(23)
where all indices are modulo $`2`$. Thus
$$\begin{array}{ccccccccc}0& & \left[\left(_4\right)\right]\left[\left(_2\right)\right]& & K^1(\mathrm{P}^3\times \mathrm{P}^5)& & _2& & 0\\ & & & & & & & & \\ & & _4_2& & & & & & \\ 0& & \left[\right]\left[\left(_2\right)\left(_4\right)\right]& & K^0(\mathrm{P}^3\times \mathrm{P}^5)& & 0& & 0\\ & & & & & & & & \\ & & _4_2_2& & & & & & \end{array}$$
(24)
Using the duality <sup>7</sup><sup>7</sup>7I am grateful to Ulrike Tillmann for sketching to me how one could give a rigorous proof between the torsion part of $`K^0`$ and $`K^1`$ for an even-dimensional orientable manifold, we arrive at the following result:
$`K^1(\mathrm{P}^3\times \mathrm{P}^5)`$ $`=`$ $`^2_4_2^2`$
$`K^0(\mathrm{P}^3\times \mathrm{P}^5)`$ $`=`$ $`^2_4_2^2`$ (25)
The order of the torsion subgroups of $`K^0`$ and $`H^{\mathrm{ev}}`$ does not match. Tracing it back through our calculation, we see that this stems from the well-known fact that the order of the torsion of a tensor product is not determined by the orders of the torsion subgroups of the factors. To be precise $`_2_4=_2`$, while $`_2(_2_2)=_2_2`$.
### 4.3 Complete Intersections
For physical reasons it would be nice if the underlying space is Calabi-Yau. Unfortunately hypersurfaces in toric varieties have torsion-free K-theory:
The smooth toric variety (of complex dimension $`m`$) does not have torsion in integer homology. The Lefschetz hyperplane theorem yields torsion free homology of the hypersurface in (real) dimensions $`0`$ to $`m1`$. But Poincaré duality then fixes the torsion part of the whole homology, since $`H_i(X,)_{\mathrm{tors}}H_{dimX1i}(X,)_{\mathrm{tors}}`$. Duality with integer cohomology then gives rise to torsion free cohomology. But torsion in integer cohomology is necessary for K-theory torsion.
## 5 Multiply Connected Spaces
We have seen that integer cohomology provides a necessary although not sufficient tool to determine whether a given manifold has torsion in K-theory. The purpose of this section is to give an easy sufficient criterion. The idea is that line bundles are stably isomorphic if and only if they are isomorphic, so stability is not a relevant concept for one-dimensional vector bundles. Then we just have to construct line bundles where a certain finite sum is (stable) trivial. This happens if the first Chern class $`c_1H^2(X,)`$ is torsion.
### 5.1 Line Bundles
Let us have a closer look to the aforementioned properties of line bundles. A $`n`$-dimensional vector bundle is in general defined via its transition functions on some open cover $`X=_{iI}U_i`$:
$$g_{ij}:U_iU_jU(n,)$$
(26)
For a line bundle, this means
$$g_{ij}:U_iU_jU(1)$$
(27)
Now two line bundles $`L_1,L_2`$ (with transition functions $`g^{(1)},g^{(2)}`$) are stably isomorphic if there exists an $`n`$:
$$L_1^nL_2^n$$
(28)
But the determinant bundle of a line bundle plus a trivial bundle is again the line bundle. Remember that the transition functions $`\stackrel{~}{g}_{ij}^{(k)}`$ of the determinant bundle $`^{n+1}(L_k^n)`$ are the determinants of the transition function matrices of $`L_k^n`$:
$$\stackrel{~}{g}_{ij}^{(k)}=det\left(\begin{array}{cccc}g_{ij}^{(k)}& & & \\ & 1& & \\ & & \mathrm{}& \\ & & & 1\end{array}\right)=g_{ij}^{(k)}k=1,2$$
(29)
Therefore
$$L_1^nL_2^n^{n+1}(L_1^n)^{n+1}(L_2^n)L_1L_2$$
(30)
Of course the “$``$” is trivial.
By a standard argument we identify then the isomorphism classes of transition functions of line bundles with the Cech cohomology group $`H^1(X,C^0(U(1)))`$, where $`C^0(U(1))`$ is the sheaf of $`U(1)`$-valued continuous functions. The long exact sequence associated to the exponential short (sheaf) exact sequence is then
$$\mathrm{}H^1(X,C^0())H^1(X,C^0\left(U(1)\right))\stackrel{c_1}{}H^2(X,)H^2(X,C^0())\mathrm{}$$
(31)
which yields the desired isomorphism since $`C^0()`$ is a fine sheaf, $`H^i(X,C^0())=0i1`$.
### 5.2 Adding line bundles
Now assume $`E`$ is a line bundle on $`X`$ with $`0c_1(E)H^2(X,)`$ pure torsion (according to the previous section then $`[E][1]0K(X)`$). But observe that the group law in $`K(X)`$ is based on the Whitney sum $`EE`$, while the group law in $`H^2(X,)`$ corresponds to the tensor product<sup>8</sup><sup>8</sup>8By the isomorphism in eq. (31) this is the group law in $`H^1(X,C^0(U(1)))`$, which corresponds to multiplying the $`U(1)`$ transition functions $`EE`$. And of course $`[E]K(X)`$ does not generate a torsion subgroup since
$$dim(n[E])=dim\left(\underset{n\mathrm{times}}{\underset{}{[E]+\mathrm{}+[E]}}\right)=n0n\{0\}$$
(32)
However $`[E][1]K(X)`$ is a torsion element ($`1`$ denotes the trivial line bundle). This follows from the Chern isomorphism:
###### Corollary 1
Let $`0xK(X)`$. Then $`x`$ is a torsion element if and only if $`ch(x)=0`$.
###### Proof 1
* $``$”: Since $`ch:K(X)H^{\mathrm{ev}}(X,)`$ is a group homomorphism this is trivial.
* $``$”: Assume that $`xK(X)`$ is free but $`ch(x)=0`$. Thus $`dim(img(ch))<rk(K(X))`$, in contradiction to the Chern isomorphism (eq. (1)).
In our case the Chern character $`ch(E)=e^{c_1(E)}=1+c_1(E)+\mathrm{}=1H^{\mathrm{ev}}(X,)`$ since $`c_1(E)`$ was assumed to be a torsion element in $`H^2(X,)`$ (so its image in $`H^2(X,)`$ vanishes). Therefore $`ch([E][1])=ch(E)ch(1)=0`$ and $`[E][1]`$ generates a nontrivial torsion subgroup.
### 5.3 Multiply connected Calabi-Yau manifolds: Quintics
Consider the Fermat quintic $`Y\mathrm{P}^4`$:
$$\underset{i=1}{\overset{5}{}}z_i^5=0$$
(33)
with the $`_5=G=\{1,g,g^2,g^3,g^4\}`$ symmetry generated by
$$g:[z_1:z_2:z_3:z_4:z_5][z_1:\alpha z_2:\alpha ^2z_3:\alpha ^3z_4:\alpha ^4z_5]\alpha =e^{\frac{2\pi i}{5}}$$
(34)
The group $`G`$ acts freely on $`Y`$: The only fixed point $`[1:0:0:0:0]\mathrm{P}^4`$ of the ambient space is missed by the hypersurface eq. (33).
This means that the quotient $`X=Y/G`$ is a (nonsingular) Calabi-Yau manifold. The quotient is still projective algebraic, but of course not a complete intersection since this would contradict section 4.3; this simply means that it is a hypersurface in some projective space where one cannot eliminate all equations.
Since the quintic $`Y`$ was simply connected (as every complete intersection), we can determine the quotient’s fundamental group from the long exact homotopy sequence (for $`Y`$ as a bundle over $`X`$ with fiber $`G`$):
$$\mathrm{}\underset{=0}{\underset{}{\pi _1(G)}}\underset{=0}{\underset{}{\pi _1(Y)}}\pi _1(X)\underset{=G}{\underset{}{\pi _0(G)}}\underset{=0}{\underset{}{\pi _0(Y)}}\underset{=0}{\underset{}{\pi _0(X)}}$$
(35)
Since $`Y`$ is a complete intersection, $`h^{1,1}(Y)=h^{1,1}(\mathrm{P}^4)=1`$. The quotient $`X`$ is still Kähler (the Kähler class $`\omega =\overline{}\mathrm{log}Z^2`$ is $`G`$-invariant), so that $`h^{1,1}(X)=h^{1,1}(Y)=1`$.
The complex structure deformations $`h^{2,1}(Y)`$ correspond to the monomials modulo $`PGL(4)`$ (the automorphisms of the ambient space) and rescaling of the equation. Here there are $`\left(\genfrac{}{}{0pt}{}{5+51}{5}\right)=126`$ monomials, and $`|PGL(4)|=24`$. Therefore $`h^{2,1}(Y)=126241=101`$. The complex structure deformations of the quotient are the $`G`$-invariant monomials, straightforward counting gives $`26`$. But now by treating every coordinate separately in the $`G`$-action the full $`PGL(4)`$ is broken to the diagonal subgroup (4 parameters). Therefore $`h^{2,1}(X)=2641=21`$. This is confirmed by the Euler number
$$\chi (Y)=2\left(h^{2,1}(Y)h^{1,1}(Y)\right)=200\chi (X)=2\left(h^{2,1}(X)h^{1,1}(X)\right)=40$$
(36)
As we expect for a free $`_5`$ group action, $`\chi (Y)=5\chi (X)`$. The Hodge diamond
$$h^{p,q}(X)=\begin{array}{ccccccc}& & & 1& & & \\ & & 0& & 0& & \\ & 0& & 1& & 0& \\ 1& & 21& & 21& & 1\\ & 0& & 1& & 0& \\ & & 0& & 0& & \\ & & & 1& & & \end{array}$$
(37)
determines the free part of integer cohomology, now we have to find the torsion part. For every manifold $`H^1(X,)`$ is torsion free, since the torsion part is dual to the torsion part in $`H_0(X,)=`$. Furthermore $`H_1(X,)`$ is the abelianization of $`\pi _1(X)=_5`$, which was already abelian. Therefore $`H_1(X,)=_5`$. By the universal coefficient theorem $`H^2(X,)_{\mathrm{tors}}H_1(X,)_{\mathrm{tors}}=_5`$.
The hard part is the torsion in $`H^3`$ (Poincaré duality then determines the rest). We are going to use the following sequence :
$$0\mathrm{\Sigma }_2H_2(X,)H_2(_5)0$$
(38)
where<sup>9</sup><sup>9</sup>9This is corrected version of the sequence in $`\mathrm{\Sigma }_2`$ is the image of $`\pi _2(X)`$ in $`H_2(X,)`$. With other words $`\mathrm{\Sigma }_2`$ are the homology classes that can be represented by $`2`$–spheres.
So we need to determine $`\pi _2(X)`$ first. We know that on the covering space $`\pi _2(Y)=H_2(Y)=`$ (The Hurewicz isomorphism theorem) since $`Y`$ is simply connected. But every map $`f:S^2X`$ can be lifted to $`\stackrel{~}{f}:S^2Y`$ since the $`S^2`$ is simply connected. That is the $`S^2`$ cannot wrap the nontrivial $`S^1X`$. More formally we can use the homotopy long exact sequence:
$$\mathrm{}\underset{=0}{\underset{}{\pi _2(G)}}\pi _2(Y)\pi _2(X)\underset{=0}{\underset{}{\pi _1(G)}}\mathrm{}$$
(39)
to show that $`\pi _2(X)=\pi _2(Y)=`$.
The group homology $`H_2(_5)=0`$, therefore eq. (38) determines an isomorphism $`\mathrm{\Sigma }_2H_2(X,)`$. We know already that the free part $`H_2(X,)_{\mathrm{free}}=`$ from the Hodge diamond. But then the map $`\pi _2(X)\mathrm{\Sigma }_2`$ must have been injective since the domain is $``$ and the image at least $``$. Therefore $`\mathrm{\Sigma }_2=`$ and the torsion part $`H^3(X,)_{\mathrm{tors}}H_2(X,)_{\mathrm{tors}}=0`$.
We have seen that
$$H^i(X,)=\{\begin{array}{cc}& i=6\hfill \\ _5& i=5\hfill \\ & i=4\hfill \\ ^{44}& i=3\hfill \\ _5& i=2\hfill \\ 0& i=1\hfill \\ & i=0\hfill \end{array}$$
(40)
From the Atiyah–Hirzebruch spectral sequence it is obvious that either the $`_5`$ torsion part survives to K-theory or vanishes (there is no subgroup except the trivial group). But according to the previous section there exists a torsion subgroup. Therefore $`K(X)_{\mathrm{tors}}=_5`$. Using Chern isomorphism and duality, this determines K-theory completely:
$$K^i(X)=\{\begin{array}{cc}^{44}_5& i=1\hfill \\ ^4_5& i=0\hfill \end{array}$$
(41)
### 5.4 The Tian–Yau manifold
The first known example was the the Tian–Yau threefold $`X=Y/_3`$, a complete intersection $`Y`$ in $`\mathrm{P}^3\times \mathrm{P}^3`$ with a free $`_3`$ group action.
The Hodge diamond can be found via counting monomials (see for details):
$$h^{p,q}(Y)=\begin{array}{ccccccc}& & & 1& & & \\ & & 0& & 0& & \\ & 0& & 14& & 0& \\ 1& & 23& & 23& & 1\\ & 0& & 14& & 0& \\ & & 0& & 0& & \\ & & & 1& & & \end{array}h^{p,q}(X)=\begin{array}{ccccccc}& & & 1& & & \\ & & 0& & 0& & \\ & 0& & 6& & 0& \\ 1& & 9& & 9& & 1\\ & 0& & 6& & 0& \\ & & 0& & 0& & \\ & & & 1& & & \end{array}$$
(42)
Since the fundamental group of the quotient $`\pi _1(X)=_3`$ we also know the torsion part $`H_1(X,)=_3`$. It remains to determine the torsion part of $`H^3(X,)_{\mathrm{tors}}H^4(X,)_{\mathrm{tors}}H_2(X,)_{\mathrm{tors}}`$.
But in this case the sequence eq. (38) does not suffice. Again the group homology $`H_2(_3)=0`$, but now it is unclear what $`\mathrm{\Sigma }_2`$ is. All we know is $`\pi _2(X)=\pi _2(Y)=^{14}`$, and the free part $`H_2(X,)_{\mathrm{free}}=^6`$. There is no reason for the image $`\pi _2(X)\mathrm{\Sigma }_2`$ not to contain a torsion part.
Therefore
$$H^i(X,)=\{\begin{array}{cc}& i=6\hfill \\ _3& i=5\hfill \\ ^6T& i=4\hfill \\ ^{20}T& i=3\hfill \\ ^6_3& i=2\hfill \\ 0& i=1\hfill \\ & i=0\hfill \end{array}$$
(43)
where $`T`$ is the unknown torsion part. The Atiyah–Hirzebruch spectral sequence not only kills torsion subgroups in cohomology, it also puts them together differently via extensions. So we know that K-theory has torsion, but cannot determine the groups.
## 6 Conclusion
As we have explicitly seen the order of K-theory-torsion and cohomology torsion is in general different. Thus substituting integer cohomology for K-theory not only leads to the wrong charge addition rules, it also does not yield the correct number of charges. Although not being totally independent, one must consider the whole spectral sequence connecting them.
This implies that discrete torsion on the field theory level must be different from the K-theory interpretation of D-brane charges. The most promising idea for a complete treatment is trying to find a pairing (preferably a perfect pairing) between K-theory and something else (maybe again K-theory) to $`U(1)`$ and use this to construct a suitable partition function, as in .
The whole discussion might be even relevant to the real world, since phenomenologically interesting string compactifications need finite non-zero $`H_1(X,)`$ in order to further break the gauge group via Wilson lines. But by the universal coefficient theorem, $`H^2(X,)_{\mathrm{tors}}H_1(X,)_{\mathrm{tors}}`$, so torsion charges appear in all realistic compactifications.
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# 1 Introduction, motivation, and summary
## 1 Introduction, motivation, and summary
Let $``$ denote the class of all probability measures (distributions) on the real line. This paper concerns the transformation $`T`$ defined on $``$ by letting $`T\nu `$ be the distribution of
$$UZ+(1U)Z^{}+g(U),$$
where $`U`$, $`Z`$, and $`Z^{}`$ are independent, with $`Z\nu `$, $`Z^{}\nu `$, and $`U\text{unif}(0,1)`$, and where
(1.1)
$$g(u):=2u\mathrm{ln}u+2(1u)\mathrm{ln}(1u)+1.$$
Of course, $`T`$ can be regarded as a transformation on the class of characteristic functions $`\psi `$ of elements of $``$. With this interpretation, $`T`$ takes the form
$$(T\psi )(t)=𝐄\left[\psi (Ut)\psi ((1U)t)\mathrm{exp}[itg(U)]\right]=_{u=0}^1\psi (ut)\psi ((1u)t)e^{itg(u)}𝑑u,t𝐑.$$
It is well known that (i) among distributions with zero mean and finite variance, $`T`$ has a unique fixed point, call it $`\mu `$; and (ii) if $`C_n`$ denotes the random number of key comparisons required by the algorithm Quicksort to sort a file of $`n`$ records, then the distribution of $`(C_n𝐄C_n)/n`$ converges weakly to $`\mu `$.
There are other fixed points. For example, it has been noted frequently in the literature that the location family generated by $`\mu `$ is a family of fixed points. But there are many more fixed points, as we now describe. Define the Cauchy($`m,\sigma `$) distribution (where $`m𝐑`$ and $`\sigma 0`$) to be the distribution of $`m+\sigma C`$, where $`C`$ has the standard Cauchy distribution with density $`x[\pi (1+x^2)]^1`$, $`x𝐑`$; equivalently, Cauchy($`m,\sigma `$) is the distribution with characteristic function $`e^{imt\sigma |t|}`$. \[In particular, the Cauchy($`m,0`$) distribution is unit mass at $`m`$.\] Now let $``$ denote the class of all fixed points of $`T`$, and let $`𝒞`$ denote the class of convolutions of $`\mu `$ with a Cauchy distribution. Using characteristic functions it is easy to check that $`𝒞`$, and that all of the distributions in $`𝒞`$ are distinct. In this paper we will prove that, conversely, $`𝒞`$, and thereby establish the following main result.
###### Theorem 1.1.
The class $``$ equals $`𝒞`$. That is, a measure $`\nu `$ is a fixed point of the Quicksort transformation $`T`$ if and only if it is the convolution of the limiting Quicksort measure $`\mu `$ with a Cauchy distribution of arbitrary center $`m`$ and scale $`\sigma `$. In particular, $``$ is in one-to-one correspondence with the set $`\{(m,\sigma ):m𝐑,\sigma 0\}`$.
The following corollary is immediate and strengthens Rösler’s characterization of $`\mu `$ as the unique element of $``$ having zero mean and finite variance.
###### Corollary 1.2.
The limiting Quicksort measure $`\mu `$ is the unique fixed point of the Quicksort transformation $`T`$ having finite expectation equal to $`0`$. ∎
The present paper can be motivated in two ways. First, the authors are writing a series of papers refining and extending Rösler’s probabilistic analysis of Quicksort. No closed-form expressions are known for any of the standard functionals (e.g., characteristic function, distribution function, density function) associated with $`\mu `$; information to be obtained about $`\mu `$ must be read from the fixed-point identity it satisfies. We were curious as to what extent additional known information about $`\mu `$, such as the fact that it has everywhere finite moment generating function, must be brought to bear. As one example, it is believed that the continuous Lebesgue density $`f`$ (treated in ) for $`\mu `$ decays at least exponentially quickly to $`0`$ at $`\pm \mathrm{}`$, cf. . But we now know from Theorem 1.1 that there there can be no proof for this conjecture that solely makes use of the information that $`\mu `$.
Second, we view the present paper as a pilot study of fixed points for a class of distributional transformations on the line. In the more general setting, we would be given (the joint distribution of) a sequence $`(A_i:i0)`$ of random variables and would define a transformation $`T`$ on $``$ by letting $`T\nu `$ be the distribution of $`A_0+_{i=1}^{\mathrm{}}A_iZ_i`$, where $`Z_1,Z_2,\mathrm{}`$ are independent random variables with distribution $`\nu `$. \[To ensure well-definedness, one might (for example) require that (almost surely) $`A_i0`$ for only finitely many values of $`i`$.\] For probability measures $`\nu `$ on $`[0,\mathrm{})`$, rather than on $`𝐑`$, and with the additional restrictions that $`A_0=0`$ and $`A_i0`$ for all $`i1`$, such transformations are called generalized smoothing transformations. These have been thoroughly studied by Durrett and Liggett , Guivarc’h , and Liu , and by other authors; consult the three papers we have cited here for further bibliographic references. Generalized smoothing transformations have applications to interacting particle systems, branching processes and branching random walk, random set constructions, and statistical turbulence theory. The arguments used to characterize the set of fixed points for generalized smoothing transformations make heavy use of Laplace transforms; unfortunately, these arguments do not carry over readily to distributions on the line. Other authors (see, e.g., ) have treated fixed points of transformations of measures $`\nu `$ on the whole line as discussed above, but not without finiteness conditions on the moments of $`\nu `$.
We now outline our proof of Theorem 1.1. Let $`\psi `$ be the characteristic function of a given $`\nu `$, and let $`r(t):=\psi (t)1`$, $`t𝐑`$. In Section 2 we establish and solve (in a certain sense) an integral equation satisfied by $`r`$. In Section 3 we then use the method of successive substitutions to derive asymptotic information about $`r(t)`$ as $`t0`$, showing first that $`r(t)=O(t^{2/3})`$, next that $`r(t)=\beta t+o(t)`$ for some $`\beta =\sigma +im𝐂`$ with $`\sigma 0`$, and finally that $`r(t)=\beta t+O(t^2)`$ . In Section 4 we use this information to argue that there exist random variables $`Z_1\nu `$, $`Z_2\mu `$, and $`C\text{Cauchy}(m,\sigma )`$ such that $`Z_1=Z_2+C`$. We finish the proof by showing that one can take $`Z_2`$ and $`C`$ to be independent, whence $`\nu 𝒞`$.
## 2 An integral equation
Let $`\psi `$ denote the characteristic function of a given $`\nu `$. Since $`\psi (t)\overline{\psi (t)}`$, we shall only need to consider $`\psi (t)`$ for $`t0`$. For notational convenience, define
$$r(t):=\psi (t)1,t0.$$
Rearranging the fixed-point integral equation $`(T\psi )(t)\psi (t)`$, we obtain the following result.
###### Lemma 2.1.
The function $`r`$ satisfies the integral equation
$$r(t)=2_{u=0}^1r(ut)𝑑u+b(t),t0,$$
where
(2.1)
$$b(t):=_{u=0}^1r(ut)r((1u)t)𝑑u+it_{u=0}^1[\psi (ut)\psi ((1u)t)1]g(u)𝑑u+a(t)$$
with
(2.2)
$$|a(t):=_{u=0}^1\psi (ut)\psi ((1u)t)[e^{itg(u)}1itg(u)]du|\frac{1}{2}𝐄g^2(U)t^2=(\frac{7}{6}\frac{1}{9}\pi ^2)t^2\text{.}$$
Note that $`r`$ and $`b`$ are continuous on $`[0,\mathrm{})`$, with $`r(0)=0=b(0)`$. Regarding $`b`$ as “known”, the integral equation in Lemma 2.1 is easily “solved” for $`r`$:
###### Proposition 2.2.
For some constant $`c𝐂`$, we have
$$\frac{r(t)}{t}=c2_{v=t}^1\frac{b(v)}{v^2}𝑑v+\frac{b(t)}{t},t>0.$$
###### Proof.
Setting $`h(t):=t[r(t)b(t)]`$, Lemma 2.1 implies
$$h(t)=2_{v=0}^t\left[\frac{h(v)}{v}+b(v)\right]𝑑v,t>0.$$
Thus $`h`$ is continuously differentiable on $`(0,\mathrm{})`$ and satisfies the differential equation
$$h^{}(t)=\frac{2}{t}h(t)+2b(t)$$
there. This is an easy differential equation to solve for $`h`$, and we find that
$$h(t)=ct^22t^2_{v=t}^1\frac{b(v)}{v^2}𝑑v,t>0,$$
for some $`c𝐂`$. After rearrangement, the proposition is proved. ∎
## 3 Behavior of $`r`$ near $`0`$
We now proceed in stages, using Proposition 2.2 as our basic tool, to get ever more information about the behavior of $`r`$ (especially near $`0`$).
###### Lemma 3.1.
Let $`\psi 1+r`$ denote the characteristic function of a given $`\nu `$. Then there exists a constant $`C<\mathrm{}`$ such that
$`|r(t)|Ct^{2/3}`$ for all $`t0`$.
###### Proof.
Let
$$M(t):=\mathrm{max}\{|r(s)|:0st\}2,t0.$$
From (2.1) and (2.2), we see immediately that, for $`0<t1`$,
$$|b(t)|M^2(t)+O(t).$$
Therefore, for $`0<t<1`$, Proposition 2.2 yields
$$|r(t)|M^2(t)+2t_{v=t}^1\frac{M^2(v)}{v^2}𝑑v+ϵ(t)=M^2(t)+2_{u=t}^1M^2(t/u)𝑑u+ϵ(t),$$
where
$$ϵ(t)=O\left(t\mathrm{log}\left(\frac{1}{t}\right)\right)+O(t)=O(t^{2/3})\text{.}$$
Consequently, again for $`0<t<1`$ (but then trivially for all $`t0`$),
$$M(t)M^2(t)+2_{u=0}^1M^2(t/u)𝑑u+O(t^{2/3}).$$
Fix $`0<a<1`$; later in the proof we shall see that $`a=1/8`$ suffices for our purposes. Since $`M(t)0`$ as $`t0`$, we can choose $`t_0>0`$ such that $`M(t_0)a`$. Then, for $`0tt_0`$,
$`M(t)`$ $``$ $`M^2(t)+2{\displaystyle _{u=0}^{t/t_0}}M^2(t/u)𝑑u+2{\displaystyle _{u=t/t_0}^1}M^2(t/u)𝑑u+O(t^{2/3})`$
$``$ $`aM(t)+8{\displaystyle \frac{t}{t_0}}+2a{\displaystyle _{u=t/t_0}^1}M(t/u)𝑑u+O(t^{2/3})`$
and thus
$$M(t)\frac{2a}{1a}_{u=0}^1M(t/u)𝑑u+O(t^{2/3}).$$
Since $`M`$ is bounded, this is trivially true also for $`t>t_0`$. Summarizing, for some constant $`\stackrel{~}{C}<\mathrm{}`$ we have, with $`U\text{unif}(0,1)`$,
(3.1)
$$M(t)\frac{2a}{1a}𝐄M(t/U)+\stackrel{~}{C}t^{2/3},t0.$$
Now fix the value of $`a`$ to be any number in $`(0,1/7)`$, say $`a=1/8`$. Then a straightforward induction \[substituting (3.2) into (3.1) for the induction step\] shows that for any nonnegative integer $`n`$ we have, for all $`t0`$,
(3.2)
$$M(t)\left(\frac{2a}{1a}\right)^n𝐄M\left(\frac{t}{U_1\mathrm{}U_n}\right)+\frac{1a}{17a}\stackrel{~}{C}t^{2/3}.$$
Recalling that $`M`$ is bounded and letting $`n\mathrm{}`$, we obtain the desired conclusion, with $`C:=\frac{1a}{17a}\stackrel{~}{C}`$. ∎
###### Lemma 3.2.
Let $`\psi 1+r`$ denote the characteristic function of a given $`\nu `$, and define $`b`$ by (2.1). Then
$`r(t)=(c2J)t+o(t)`$ as $`t0`$,
where $`J`$ is the absolutely convergent integral
(3.3)
$$J:=_{v=0}^1\frac{b(v)}{v^2}𝑑v.$$
###### Proof.
Combining (2.1)–(2.2) and Lemma 3.1, we obtain
$$|b(t)|O(t^{4/3})+O(t^{1+(2/3)})+O(t^2)=O(t^{4/3}).$$
Thus the integral $`J`$ converges absolutely, and from Proposition 2.2 we obtain the desired conclusion about $`r`$. ∎
Lemma 3.2 is all we will need in the next section, but the following refinement follows readily and takes us as far as we can go with the method of successive substitutions.
###### Corollary 3.3.
Let $`\psi `$ denote the characteristic function of a given $`\nu `$. Then there exists a constant $`\beta =im\sigma 𝐂`$ with $`\sigma 0`$ such that
$`\psi (t)=1+\beta t+O(t^2)`$ as $`t0`$.
###### Proof.
Combining (2.1)–(2.2) and Lemma 3.2, we readily obtain $`b(t)=O(t^2)`$. Therefore, by Proposition 2.2,
$$\psi (t)1=r(t)=(c2J)t+2t_{v=0}^t\frac{b(v)}{v^2}𝑑v+b(t)=\beta t+O(t^2),$$
with $`\beta =im\sigma :=c2J`$. Since $`|\psi (t)|1`$ for all $`t`$, we must have $`\sigma 0`$. ∎
## 4 Proof of the main theorem
### 4.1 Further preliminaries
In Sections 4.14.2 we complete the proof of our main Theorem 1.1. To do this, we begin with a key result that any characteristic function with expansion as in Corollary 3.3 \[more generally, we allow the remainder term there to be simply $`o(t)`$\] is in the domain of attraction of (iterates of) the “homogeneous” analogue $`T_0`$ of $`T`$. (Here $``$ denotes weak convergence of probability measures.)
###### Theorem 4.1.
Let $`\psi `$ be any characteristic function satisfying
(4.1) $`\psi (t)=1+\beta t+o(t)=1+imt\sigma t+o(t)`$ as $`t0`$
for some $`\beta =im\sigma `$ $`𝐂`$, with $`m𝐑`$ and $`\sigma 0`$. Let $`\nu `$ be the corresponding probability measure. Then
$$T_0^n\nu \text{Cauchy}(m,\sigma ),$$
where $`T_0`$ is the homogeneous analogue of the Quicksort transformation $`T`$ mapping distributions as follows (in obvious notation):
(4.2)
$$T_0:ZUZ+(1U)Z^{}.$$
###### Proof.
Let $`Z_1,Z_2,\mathrm{};U_1,U_2,\mathrm{}`$ be independent random variables, with every $`Z_i\nu `$ and every $`U_j\text{unif}(0,1)`$. Then, using the definition of $`T_0`$ repeatedly,
$$W_n:=\underset{i=1}{\overset{2^n}{}}V_i^{(n)}Z_iT_0^n\nu ,n0,$$
where we define the random variables $`V_i^{(n)}`$ as follows. Using $`U_1`$ in the obvious fashion, split the unit interval into intervals of lengths $`U_1`$ and $`1U_1`$. Now using $`U_2`$ and $`U_3`$, split the first interval into subintervals of lengths $`U_1U_2`$ and $`U_1(1U_2)`$ and the second interval into subintervals of lengths $`(1U_1)U_3`$ and $`(1U_1)(1U_3)`$. Continue in this way (using $`U_1,\mathrm{},U_{2^n1}`$) until the unit interval has been divided overall into $`2^n`$ subintervals. Call their lengths, from left to right, $`V_1^{(n)},\mathrm{},V_{2^n}^{(n)}`$.
Let $`L_n:=\mathrm{max}(V_1^{(n)},\mathrm{},V_{2^n}^{(n)})`$. We show that $`L_n`$ converges in probability to $`0`$ as $`n\mathrm{}`$. Luckily, the complicated dependence structure of the variables $`V_i^{(n)}`$ does not come into play; the only observation we need is that that each $`V_i^{(n)}`$ marginally has the same distribution as $`U_1\mathrm{}U_n`$. Indeed, abbreviate $`V_1^{(n)}`$ as $`V_n`$; briefly put, we derive a Chernoff’s bound for $`\mathrm{ln}(1/V_n)`$ and then simply use subadditivity. To spell things out, let $`x>0`$ be fixed and let $`t0`$. Then
$$𝐏(V_ne^x)e^{tx}𝐄V_n^t=e^{tx}\underset{j=1}{\overset{n}{}}𝐄U_j^t=e^{tx}(1+t)^n=\mathrm{exp}[(n\mathrm{ln}(1+t)xt)].$$
Choosing the optimal $`t=\frac{n}{x}1`$ (valid for $`nx`$), this yields
$$𝐏(V_ne^x)\mathrm{exp}[(n\mathrm{ln}(n/x)n+x)]=\mathrm{exp}[(n(\mathrm{ln}n\mathrm{ln}(ex))+x)]$$
and thus
$$𝐏(L_ne^x)2^n\mathrm{exp}[(n(\mathrm{ln}n\mathrm{ln}(ex))+x)]=\mathrm{exp}[(n(\mathrm{ln}n\mathrm{ln}(2ex))+x)]0$$
as $`n\mathrm{}`$.
Since $`L_n`$ converges in probability to $`0`$, we can therefore choose $`ϵ_n0`$ so that $`𝐏(L_n>ϵ_n)0`$. To prove the theorem, it then suffices to prove
$$\stackrel{~}{W}_n:=\mathrm{𝟏}(L_nϵ_n)W_n\text{Cauchy}(m,\sigma ).$$
For this, we note that the characteristic function $`\varphi _n`$ of $`\stackrel{~}{W}_n`$ is given for $`t𝐑`$ by
(4.3)
$$\varphi _n(t)=𝐏(L_n>ϵ_n)+𝐄\left[\mathrm{𝟏}(L_nϵ_n)\underset{i=1}{\overset{2^n}{}}\psi (V_i^{(n)}t)\right].$$
We will show that $`\varphi _n(t)`$ converges to $`e^{\beta t}=e^{imt\sigma t}`$ for each fixed $`t0`$, and \[since, further, $`\varphi _n(t)\overline{\varphi _n(t)}`$\] this will complete the proof of the lemma.
Indeed, we need only consider the second term in (4.3). For that, the calculus estimates outlined in the proof of the lemma preceding Theorem 7.1.2 in demonstrate that, when $`L_nϵ_n`$,
$$\underset{i=1}{\overset{2^n}{}}\psi (V_i^{(n)}t)=(1+D_n)e^{\beta t}$$
for complex random variables $`D_n`$ (depending on our fixed choice of $`t0`$) satisfying $`|D_n|\delta _n`$ for a deterministic sequence $`\delta _n[\delta (ϵ_nt)]0`$ \[with $`\delta (s)0`$ as $`s0`$\]. \[Leaving out the error estimates, the argument is
$`\mathrm{log}\left[{\displaystyle \underset{i=1}{\overset{2^n}{}}}\psi (V_i^{(n)}t)\right]`$ $``$ $`{\displaystyle \underset{i=1}{\overset{2^n}{}}}\left(\psi (V_i^{(n)}t)1\right)`$
$``$ $`{\displaystyle \underset{i=1}{\overset{2^n}{}}}\beta V_i^{(n)}t=\beta t.]`$
It now follows easily that $`\varphi _n(t)e^{\beta t}`$, as desired. ∎
Both the next lemma and its immediate corollary (Lemma 4.3) will be used in our proof of Theorem 1.1.
###### Lemma 4.2.
Let $`\nu _i`$, $`i=1,2`$. Suppose that $`(Z_1,Z_2)`$ is a coupling of $`\nu _1`$ and $`\nu _2`$ such that the characteristic function of $`Z_1Z_2`$ satisfies (4.1). Then there exists a coupling $`(\stackrel{~}{Z}_1,\stackrel{~}{Z}_2)`$ of $`\nu _1`$ and $`\nu _2`$ such that $`\stackrel{~}{Z}_1\stackrel{~}{Z}_2`$ Cauchy$`(m,\sigma )`$.
###### Proof.
Extend $`T`$ to a transformation $`T_2`$ on the class $`_2`$ of probability measures on $`𝐑^2`$ by mapping the distribution $`\xi _2`$ of $`(X,Y)`$ to the distribution $`T_2\xi `$ of
$$(UX+(1U)X^{}+g(U),UY+(1U)Y^{}+g(U)),$$
where $`U`$, $`(X,Y)`$, and $`(X^{},Y^{})`$ are independent, with $`(X,Y)\xi `$, $`(X^{},Y^{})\xi `$, and $`U\text{unif}(0,1)`$, and where $`g`$ is given by (1.1). (Note that we use the same uniform $`U`$ for the $`Y`$s as for the $`X`$s!) Of course, $`T_2`$ maps the marginal distributions $`\xi _1()=\xi (\times 𝐑)`$ of $`X`$ and $`\xi _2()=\xi (𝐑\times )`$ of $`Y`$ into $`T\xi _1`$ and $`T\xi _2`$, respectively; more importantly for our purposes, it maps the distribution, call it $`\widehat{\xi }`$, of $`XY`$ into the distribution $`T_0\widehat{\xi }`$, with $`T_0`$ defined at (4.2).
Now let $`\nu _2`$ have marginals $`\nu _i`$, $`i=1,2`$. Then ($`T_2^n\nu )_{n1}`$ has constant marginals $`(\nu _1,\nu _2)`$ as $`n`$ varies and so is a tight sequence. We then can find a weakly convergent subsequence, say,
$$T_2^{n_k}\nu \nu ^{\mathrm{}}_2;$$
of course, the limit $`\nu ^{\mathrm{}}`$ again has marginals $`\nu _i`$, $`i=1,2`$. Moreover,
$$T_0^{n_k}\widehat{\nu }=\widehat{T_2^{n_k}\nu }\widehat{\nu ^{\mathrm{}}}.$$
But, by supposition, the characteristic function of $`\widehat{\nu }`$ satisfies (4.1), so Theorem 4.1 implies that $`\widehat{\nu ^{\mathrm{}}}`$ is Cauchy$`(m,\sigma )`$. Thus $`\nu ^{\mathrm{}}_2`$ supplies the desired coupling. ∎
###### Lemma 4.3.
Let $`\nu _i`$, $`i=1,2`$. Suppose that $`(Z_1,Z_2)`$ is a coupling of $`\nu _1`$ and $`\nu _2`$ such that $`Z_1Z_2`$ has zero mean and finite variance. Then $`\nu _1=\nu _2`$. ∎
### 4.2 The proof
We now complete the proof of Theorem 1.1.
###### Proof.
As discussed in Section 1, it is simple to check that $`𝒞`$ (and that the elements of $`𝒞`$ are all distinct).
Conversely, given $`\nu `$, let $`Z_1\nu _1:=\nu `$ and $`Z_2\nu _2:=\mu `$ be independent random variables (on some probability space); recall that $`\mu `$ is the limiting Quicksort measure, with zero mean and finite variance. Write $`\psi _i`$, $`i=1,2`$, for the characteristic functions corresponding respectively to $`\nu _i`$, $`i=1,2`$. By Lemma 3.2 (or see Corollary 3.3), $`\psi _1`$ satisfies (4.1) \[for some $`(m,\sigma )`$\]. Of course, $`\psi _2`$ satisfies (4.1) with $`\beta `$ taken to be $`0`$, so the characteristic function $`t\psi _1(t)\psi _2(t)`$ of $`Z_1Z_2`$ satisfies (4.1) for the same $`(m,\sigma )`$ as for $`\psi _1`$. Applying Lemma 4.2, there exists a coupling $`(\stackrel{~}{Z}_1,\stackrel{~}{Z}_2)`$ of $`\nu _1`$ and $`\nu _2`$ such that $`C:=\stackrel{~}{Z}_1\stackrel{~}{Z}_2`$ Cauchy$`(m,\sigma )`$. Without loss of generality (by building a suitable product space), we may assume the existence of a random variable $`Y\mu `$ on the same probability space as $`\stackrel{~}{Z}_1`$ and $`\stackrel{~}{Z}_2`$ such that $`Y`$ and $`C`$ are independent.
We know that the distribution $`\nu _1`$ of $`\stackrel{~}{Z}_1=\stackrel{~}{Z}_2+C`$ is a fixed point of $`T`$. But so is the distribution $`\nu _1^{}𝒞`$ of $`Z:=Y+C`$. By Lemma 4.3 applied to $`(\stackrel{~}{Z}_1,Z)`$, $`\nu =\nu _1=\nu _1^{}𝒞`$, as desired. ∎
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# Aspects of the stochastic Burgers equation and their connection with turbulence.
Summary.
We present results for the 1 dimensional stochastically forced Burgers equation when the spatial range of the forcing varies. As the range of forcing moves from small scales to large scales, the system goes from a chaotic, structureless state to a structured state dominated by shocks. This transition takes place through an intermediate region where the system exhibits rich multifractal behavior. This is mainly the region of interest to us. We only mention in passing the hydrodynamic limit of forcing confined to large scales, where much work has taken place since that of Polyakov.
In order to make the general framework clear, we give an introduction to aspects of isotropic, homogeneous turbulence, a description of Kolmogorov scaling, and, with the help of a simple model, an introduction to the language of multifractality which is used to discuss intermittency corrections to scaling.
We continue with a general discussion of the Burgers equation and forcing, and some aspects of three dimensional turbulence where - because of the mathematical analogy between equations derived from the Navier-Stokes and Burgers equations - one can gain insight from the study of the simpler stochastic Burgers equation. These aspects concern the connection of dissipation rate intermittency exponents with those characterizing the structure functions of the velocity field, and the dynamical behavior, characterized by different time constants, of velocity structure functions. We also show how the exponents characterizing the multifractal behavior of velocity structure functions in the above mentioned transition region can effectively be calculated in the case of the stochastic Burgers equation.
Table of contents.
I. Introduction.
I.1. Kolmogorov scaling.
I.2. A simple model.
I.3. The language of multifractality.
II. The stochastic Burgers equation.
II.1. Shock structure and extreme multifractality.
II.2. Stochastic forcing.
III. Three dimensional turbulence and the stochastic Burgers equation.
III.1. Multifractal exponents.
III.2. Dissipation rate correlation and intermittency.
III.3. Dynamic behavior.
IV. Remarks on intermittency.
I. Introduction.
We study some aspects of statistically stationary, homogeneous and isotropic fully developed turbulence. This is the typical framework in which such studies are done. The quantities of interest are the equal time spatial correlations of the velocity field $`\stackrel{}{u}(\stackrel{}{x},t)`$, the so-called structure functions. The longitudinal structure functions, which are the ones usually discussed, are defined by
$$S_p(r)=<[(\stackrel{}{u}(\stackrel{}{x}+\stackrel{}{r},t)\stackrel{}{u}(\stackrel{}{x},t)).\stackrel{}{n}]^p>$$
(1)
where $`\stackrel{}{n}`$ is the unit vector in the direction $`\stackrel{}{r}`$. Some components of the velocity field difference can be projected onto the direction transverse to $`\stackrel{}{n}`$, and thus there are other correlations of $`pth`$ order, which involve longitudinal and an (even) number of transverse projections.
The velocity satisfies the incompressible Navier-Stokes equation
$$_t\stackrel{}{u}+\stackrel{}{u}.\stackrel{}{}\stackrel{}{u}=\stackrel{}{}p+\nu \mathrm{}\stackrel{}{u}+\stackrel{}{f}$$
(2)
with
$$\stackrel{}{}.\stackrel{}{u}=0$$
(3)
Here $`p`$ is the pressure divided by the constant mass density, $`\nu `$ the kinematic viscosity. We have added $`\stackrel{}{f}=\stackrel{}{f}(\stackrel{}{x},t)`$, an external stochastic force which acts on large scales, and maintains a turbulent steady state. The average in (1) then includes as well an average over time.
In the usual picture of turbulence (see I.1.), when the distance $`r=|\stackrel{}{r}|`$ in (1) is small compared to large scales $`L`$ of the order of the system size, and large compared to the scales where dissipation takes place, the structure functions are expected to behave as
$$S_p(r)(r/L)^{\zeta _p}$$
(4)
An important aspect of solving the problem of statistical isotropic, homogeneous turbulence is deriving the values of the exponents $`\zeta _p`$ in (4) from the Navier-Stokes equation. This has not been done except for $`\zeta _3`$, the value of which is fixed by the Von Karman-Howarth relation. It turns out however that the experimentally measured $`\zeta _p`$’s (up to $`p=10`$ or so) are not too different from their scaling values as they arise in the picture of fully developed turbulence proposed by Kolmogorov. This is the reason a large number of phenomenological models exist, which by breaking scale invariance slightly, give improved fits to the data. The usual language in which to express deviations from scaling is that of multiscaling or multifractality.
We will therefore discuss first in this introductory section Kolmogorov scaling, then a simple model, which allows one to introduce non-scaling elements, and provides a simple introduction to the language of multiscaling which we present next. A general reference for these subjects is the book of Frisch.
In the second section we discuss the stochastic Burgers equation, its shock structure and the associated extreme multifractality, and its behavior when the spatial range of the random forcing varies from small to large scales. In section III we take up the point about statistical aspects of the stochastic Burgers equation and their connection with three dimensional, forced, isotropic and homogeneous turbulence. First we show how the problem of multifractality can be solved for the stochastic Burgers equation. Then we discuss the relation between intermittency in the energy dissipation to intermittency in the velocity field, and end up by making a number of observations concerning the dynamical behavior of structure functions. General remarks about intermittency in fully developed turbulence and for the stochastic Burgers equation are made in section IV.
This report is based on a number of results or points made in references.
I.1. Kolmogorov scaling.
The picture is that of an energy cascade from the large scale $`L`$ where the energy is put into the system, to the dissipation scale $`\delta `$ where it is dissipated. On intermediate scales $`\delta rL`$, which make up the so-called inertial range, the only quantity which matters is $`ϵ`$, the mean energy dissipation rate per unit mass, considered to be independent of scale. $`ϵ`$ has the dimension of velocity squared divided by time, or velocity cubed divided by distance.
The dissipation scale $`\delta `$ can only depend on $`\nu `$ and $`ϵ`$, and thus for dimensional reasons $`\delta (\nu ^3/ϵ)^{\frac{1}{4}}(1/Re)^{\frac{3}{4}}L`$, where after replacing $`ϵ`$ in terms of a characteristic velocity $`U`$ and the large scale $`L`$, we are able to introduce the Reynolds number $`Re=UL/\nu `$. In the limit of small viscosity or large Reynolds number there is thus a definite separation of scales between $`\delta `$ and $`L`$.
In the inertial region, dimensions are determined by $`ϵ`$ alone, and therefore one predicts on dimensional grounds, that $`S_p(r)`$ which has the dimension of velocity to the $`pth`$ power behaves as
$$S_p(r)ϵ^{\frac{p}{3}}r^{\frac{p}{3}}$$
(5)
This is Kolomogorov scaling. The scaling values of the exponents in (4) are then
$$\zeta _p=p/3$$
(6)
This gives $`\zeta _2=2/3`$, which by Fourier transform is equivalent to the experimentally observed $`5/3`$ behavior of the energy spectrum, namely $`E(k)=k^2<\stackrel{}{u}(\stackrel{}{k}).\stackrel{}{u}(\stackrel{}{k})>ϵ^{\frac{2}{3}}k^{\frac{5}{3}}`$. One also obtains $`\zeta _3=1`$, which is the value fixed by the Von Karman-Howarth relation. The other general result is that $`\zeta _p`$ is a convex function of $`p`$. Measurements of the structure functions show however that Kolmogorov scaling does not hold: the measured $`\zeta _p`$’s for $`p>3`$ lie below the scaling values. For instance $`\zeta _6=1.80\pm 0.05`$ rather than the scaling value of $`2`$, obtained from (6). This effect is called intermittency or multifractality, and can be related heuristically to the non-space filling property of the eddies which make up the energy cascade, and therefore to their fractal dimension. A simple model will serve to illustrate these points.
I.2. A simple model.
Among models which describe the energy cascade, the so-called $`\beta `$-model is instructive. Imagine, as the energy cascades down to smaller scales from the large scale $`L`$, that at scales $`r=\alpha ^nL`$ in the inertial range, the eddies at this scale, which themselves have a typical size of $`r`$, occupy only a fraction $`\beta `$ of the available space, such that $`p_r=\beta ^n`$, where $`p_r`$ can be interpreted as the probability of finding an eddy of size $`r`$ at scale $`r`$. Eliminating the ”generation” number $`n`$ between the expressions for $`r`$ and $`p_r`$, on finds
$$p_r=(r/L)^{3D}$$
(7)
where $`3D=\mathrm{ln}\beta /\mathrm{ln}\alpha `$. If the eddies are space filling, then $`\beta =1`$, and therefore $`D=3`$. The value of $`3`$ corresponds to the fact that we pretend our discussion is for eddies in 3 dimensions. The argument itself is clearly independent of space dimension. One now interprets $`D`$ as the fractal dimension of the space on which the eddies exist, assuming that $`D`$ is smaller than $`3`$.
What are the structure functions in this model?
The typical energy of an eddy of size $`r`$ is $`E_r\delta v_{r}^{}{}_{}{}^{2}p_r`$, and therefore the average energy dissipation rate (per unit mass) at scale $`r`$, with a typical time scale $`t_r=r/\delta v_r`$, is
$$ϵ_r\frac{\delta v_{r}^{}{}_{}{}^{3}}{L}(r/L)^{3D1}$$
(8)
Here $`\delta v_r`$ is the velocity variation across the eddy. The value of $`ϵ_r`$ is independent of $`r`$ if homogeneity holds (existence of an inertial scale), and therefore one has for the velocity
$$\delta v_r(ϵL)^{1/3}(r/L)^{\frac{1}{3}(3D)/3}$$
(9)
from which follows for the structure function
$$S_p(r)=<\delta v_{r}^{}{}_{}{}^{p}>=\delta v_r^pp_r(ϵL)^{p/3}(r/L)^{p/3+(3D)(1p/3)}$$
(10)
One thus finds for the exponents $`\zeta _p`$ of the structure functions, a convex function of $`p`$, namely
$$\zeta _p=p/3+(3D)(1p/3)$$
(11)
which satisfies the condition (Von Karman-Howarth relation) $`\zeta _3=1`$. The scaling violating part in $`\zeta _p`$ is given by $`(3D)(1p/3)`$. For instance $`\zeta _6=2(3D)`$, which, by comparison with the experimental result $`\zeta _6=20.2`$, leads to a fractal dimension $`D=2.8`$. Note that the velocity variation at $`r`$ ($`\delta v_rr^h`$) is itself characterized by an exponent $`h=1/3(3D)/3`$. For $`D=3`$, when the eddies fill all space at any inertial scale, one has the scaling (Kolmogorov) result $`h=1/3`$ and $`\zeta _p=p/3`$.
In the simple model we have considered, the structure functions and the variations of the velocity field are characterized by a single $`h`$ and $`D`$. However here, as opposed to the Kolmogorov scaling behavior, the eddies are not space filling, but are characterized by a fractal dimension $`D`$.
Simple fractal models such as the one we have described are not believed to give the whole picture required to describe fully developed turbulence. Experimental data suggest that $`\zeta _p`$ depends non linearly on $`p`$ in contrast to equation (11). It is believed that one needs to consider a more general picture, with a range of possible $`h`$’s and of corresponding fractal dimensions $`D(h)`$ (see section IV.). This picture, or the language in which it is formulated, is that of multifractality, which we discuss next.
I.3. The language of multifractality.
Assume now that $`h`$ can take values in an interval $`(h_{min},h_{max})`$, and that to each $`h`$ there corresponds a set in three dimensional space of fractal dimension $`D(h)`$, in such a way that across any distance $`r`$ ( $`r`$ belongs to the inertial range) in the vicinity of that set, one has
$$\delta v_r(r/L)^h$$
(12)
and
$$p_r(r/L)^{3D(h)}$$
(13)
where $`p_r`$ is the probability for being within a distance of the set of fractal dimension $`D(h)`$, and $`\delta v_r`$ is the velocity variation. As a consequence one has the following expression for $`S_p(r)`$ for a given set with scaling dimension $`h`$
$$S_p(r)<\delta v_{r}^{}{}_{}{}^{p}>(r/L)^{ph+3D(h)}$$
(14)
All $`h`$ can contribute to the right-hand side, but since $`r/L1`$, the dominant exponent $`\zeta _p`$ is given by
$$\zeta _p=\underset{h}{\mathrm{min}}(ph+3D(h))$$
(15)
This exponent $`\zeta _p`$ is the dominant one in the expression of the structure factors (cf. equation (4)).
Remarks:
\- the scaling result corresponds to $`h=1/3`$ and $`D(1/3)=3`$.
\- the argument is the same in 1 or 2 dimensions with the replacement of the number $`3`$ in $`3D(h)`$ by respectively $`1`$ and $`2`$.
\- the quantity $`3D(h)`$ is positive or zero, since $`D(h)`$ cannot exceed the dimension of the embedding space. It is generally assumed that $`h_{min}0`$. In the case of the Burgers equation where exponents can be calculated, we find (cf. section III.1.) that the $`h`$’s corresponding to higher order structure functions reach the value 0 when the stochastic forcing has moved to sufficiently large scales, and stay at the value 0 when the scale of the forcing increases further.
II. The stochastic Burgers equation.
This is a 1 dimensional version of the Navier-Stokes equation, a version without incompressibility and pressure, which describes the evolution of the compressible field $`u(x,t)`$, by
$$_tu+u\frac{u}{x}=\nu \frac{^2u}{x^2}+f$$
(16)
where $`f=f(x,t)`$ is a stochastic forcing.
We will discuss later the forcing and its influence on the dynamics of the field. For the moment, we will ignore it, and summarize some results concerning the plain Burgers equation.
II.1. Shock structure and extreme multifractality.
If one starts from an initial sinusoidal velocity profile of large wavelength, then under the influence of the nonlinear term in the equation, the sinusoid will for sufficiently small viscosity, steepen into a series of shocks. After some time the shocks will fade away, their energy being dissipated by the viscous term. This viscous term plays a role mainly at the position of the shocks, where it is counterbalanced by the nonlinear term. The equality of these two terms leads to
$$\nu =\mathrm{}u.\delta $$
(17)
where $`\mathrm{}u`$ is the velocity jump across the shock, and $`\delta `$ is the shock thickness. There are thus two scales here: a large scale $`L`$ corresponding to some average distance between shocks, and a dissipation scale $`\delta \nu `$, very much smaller than $`L`$ when $`\nu `$ goes to zero. Distances away from both extremes make up the inertial range.
In terms of multifractal language, the Burgers equation ( one averages, in the limit $`\nu 0`$, over an ensemble of initial states, or considers stochastic forcing on large scales) shows extreme multifractality, a situation called bifractality in the literature. The behavior of $`u`$ is essentially linear between shocks ($`ux`$), and thus here $`h=1,D(1)=1`$. At the shocks themselves $`h=0,D(0)=0`$, since the shocks are discontinuities of the velocity field occurring at a point (in the $`\nu 0`$ limit). The velocity variation across the shock is independent of distance, and the probability of being within a distance $`r`$ is linear in $`r`$ (cf. equations (12) and (13) for the case of 1 dimension).
There are thus two possible values for the exponent $`ph+1D(h)`$ (cf. section I.3.), namely $`p`$ or $`1`$, and therefore the dominant exponent $`\zeta _p`$ (equation (15)) characterizing the behavior of the structure functions in the inertial scale, is such that
$`\zeta _p=1,`$ $`p1`$ (18)
This is an extreme case of multifractality ( all exponents have the same value for integer $`p`$ greater than 1), very much different from the case of three dimensional homogeneous, isotropic turbulence where the experimentally determined exponents remain relatively close to the scaling ones, which increase linearly with $`p`$ (see equation (6)).
However - as we have discovered - there is a whole range of multifractal behavior as the spatial extent of the stochastic force in the Burgers equation varies, and the situation is much more interesting.
II.2. Stochastic forcing.
For the stochastic forcing in (16) we take a Gaussian, such that in $`k`$ space
$$<f(k,t)>=0$$
$$<f(k,t)f(k^{},t^{})>=2D_0|k|^\beta \delta _{k,k^{}}\delta (tt^{})$$
(19)
The exponent $`\beta `$ determines over which scales the forcing acts. For $`\beta >0`$ it acts effectively on small scales, whereas as $`\beta `$ becomes negative, larger and larger scales matter. The limit relevant to forcing in three dimensional turbulence is that of large scales, of the order of the system size $`L`$.
The range of values of $`\beta `$ goes from $`\beta =2`$, which corresponds to thermal noise, to $`\beta =3/2`$. For values smaller than the latter, the statistics of the velocity field is independent of $`\beta `$, unchanged from its behavior at $`\beta =3/2`$. At $`\beta =3/2`$ the system behaves as the steady state of the plain Burgers equation: it exhibits the extreme multifractal behavior discussed in II.1., characteristic of a shock dominated velocity field. For $`\beta >0`$ however, the presence of noise on small scales prevents the shocks from developing, and therefore the behavior appears chaotic, i.e. random and structureless. Thus as $`\beta `$ moves from positive to large negative values, the velocity field goes from a chaotic to a shock dominated state, through an intermediate region ($`3/2<\beta <0`$), where for $`1<\beta <0`$ it displays complex dynamics of appearing, interacting and disappearing shocks. This region is one of rich multifractal behavior, and is the principal object of our study. It is through this region that one approaches the hydrodynamic limit of large scale forcing from a purely chaotic state.
To be complete, we mention that for positive values of $`\beta `$ one can use a renormalization group approach. As soon as $`\beta `$ becomes negative, all sorts of non-linear terms become important in the equations, and the perturbative renormalization group approach breaks down. This approach has been usually applied to the equivalent KPZ (Kardar-Parisi-Zhang) equation for fluctuations of an interface height $`h(x,t)`$, related to $`u`$ by $`u=_xh`$. With a noise of the form considered, the renormalization group has also been applied to the Navier-Stokes equation.
For $`\beta `$ positive , close to zero, the scaling analysis leads to the following result for the exponents $`z`$ and $`\zeta _2`$, which appear in the scaling form assumed for $`S_2(r,\tau )=<(u(x+r,t+\tau )u(x,t))^2>`$, namely $`S_2(r,\tau )=r^{\zeta _2}g(\tau /r^z)`$:
$$z+\zeta _2/2=1$$
(20)
and
$$\zeta _2z=1\beta $$
(21)
The first relation is a consequence of Galilean invariance, the second of the fact that the coefficient $`D_0`$ of noise fluctuations is not rescaled because of the non-analytic form of the noise. One obtains from (20) and (21) that $`\zeta _2=2\beta /3`$ and $`z=1+\beta /3`$.
We will from now on consider the region of negative $`\beta `$, which is so to speak the gateway to hydrodynamic behavior.
III. Three dimensional turbulence and the stochastic Burgers equation.
We believe that because of the mathematical similarity of the Navier-Stokes equation with forcing, and the stochastic Burgers equation, the latter can be used as a key to the understanding of a number of issues in the statistical behavior of isotropic, homogeneous turbulence. In the work we have been doing, we highlight this similarity on a number of occasions, in different situations. To give a simple example here, we compare the Von Karman-Howarth relation for $`S_3`$ for both equations.
For the Navier-Stokes equation with forcing $`\stackrel{}{f}`$, this relation takes the following form for the (equal time) 3rd order structure function $`S_{3j}=<(\stackrel{}{u}_1\stackrel{}{u}_2)^2(u_{1j}u_{2j})>`$, where ”1” refers to the point $`\stackrel{}{x}+\stackrel{}{r}`$, ”2” to the point $`\stackrel{}{x}`$, and ”j” denotes the j-th component of $`\stackrel{}{u}`$
$$\frac{1}{2}_{r_j}S_{3j}(r)=\nu \mathrm{}S_2(r)2<ϵ>+<(\stackrel{}{u}_1\stackrel{}{u}_2).(\stackrel{}{f}_1\stackrel{}{f}_2)>$$
(22)
where $`S_2(r)=<(\stackrel{}{u}_1\stackrel{}{u}_2)^2>`$, while for the stochastic Burgers equation, where $`S_3(r)=<(u_1u_2)^3>`$, it reads
$$\frac{1}{6}dS_3(r)/dr=\nu d^2S_2/dr^22<ϵ>+<(u_1u_2)(f_1f_2)>$$
(23)
The structural similarity of the two equations is clear.
One can derive the above two Von Karman-Howarth relations in a straightforward way from the space and time dependent $`S_2`$, using the homogeneity in time of expectation values. More precisely, one writes that $`S_2(r,\tau )/t_1+S_2(r,\tau )/t_2=0`$, where $`r=x_1x_2,\tau =t_1t_2`$. This derivation highlights the fact, which we have several times pointed out in our work, that it is often useful for deriving equal time correlations to pass through time dependent calculations. Many identities can be obtained this way.
The two equations (22) and (23) are very similar. The 3 dimensional result contains Kolmogorov’s ”4/5th” law for the longitudinal structure function. In both cases $`<ϵ>`$ represents the energy dissipation rate. Since $`r`$ belongs to the inertial scale the term multiplied by $`\nu `$ is negligible in both equations in the zero viscosity limit. The noise dependent term can be evaluated in the equal time limit with the help of the Novikov-Donsker formalism. When the noise is cut-off at large scales (the hydrodynamic limit) this term leads to a subdominant correction of order $`(r/L)^2`$. We will discuss later, for the stochastic Burgers equation, the general case when the noise ranges over small scales as well.
Though this comparison of the Von Karman-Howarth relations is based on a simple case, we have found that the same similarity term by term, with an obvious display of the 3 dimensional space indices, holds for any other equation we have derived involving velocity or dissipation rate correlations, with the exclusion of course of terms involving pressure.
We will discuss in the following three main points:
(i) first, we are going to face for the stochastic Burgers equation the problem of turbulence, namely calculate, for small $`p`$, in the multifractal region ($`1<\beta <0`$) the exponents $`\zeta _p`$ characterizing the statistical behavior of velocity structure functions,
(ii) second, we are going to give the general equation satisfied by the equal time correlation of the dissipation rate, and connect its intermittent behavior, which exhibits a hierarchy of exponents, to the intermittent behavior of the velocity structure functions,
(ii) third, we investigate the dynamics of the second order structure function, and show how - even in the absence of any average flow - $`S_2`$ satisfies a wave equation with characteristic velocity $`\sqrt{<u^2>}`$. These dynamic considerations enable us to disentangle, in our Eulerian framework, the intrinsic dynamical and the kinetic, ballistic characteristic times which describe the time evolution of flow structures.
III.1. Multifractal exponents.
We are interested in the region where $`1<\beta <0`$. Here also exists the possibility of scaling behavior, in the same way as there is Kolmogorov scaling for three dimensional turbulence, where the dimension of $`<ϵ>`$ or equivalently $`D_0`$, determines the dependence on distance of the $`S_p`$’s in the inertial range. One thus has
$$S_p(r)(D_0/L)^{p/3}r^{p\beta /3}$$
(24)
which corresponds to $`\zeta _p=p\beta /3`$ and $`h=\beta /3`$. This is the value of $`h`$ in the scaling regime. (Notice that at $`\beta =1`$ the exponents are the same as those of Kolmogorov scaling, equation (6).)
This scaling regime is however dominant only in the region of $`\beta `$ negative close to zero, and gives way to multifractal behavior as $`\beta `$ goes towards $`1`$. We are going to study this behavior directly on equations for the structure functions derived from the stochastic Burgers equation. We proceed systematically discussing first $`S_2`$ and $`S_3`$, and then $`S_4,S_5`$ and general $`S_p`$.
(i) $`S_2`$ and $`S_3`$.
One cannot derive directly from the stochastic Burgers ( or from the Navier-Stokes equation in three dimensions) a closed equation for the equal time structure function $`S_2`$. We therefore check numerically that $`S_2(r)`$ behaves in the following way
$$S_2(r)(r/L)^{2\beta /3}$$
(25)
for all $`3/2<\beta <0`$. Precise numerical results, and therefore a precise value of the exponent, can be obtained from evaluating the energy spectrum ($`E(k)|k|^{1+2\beta /3}`$), related to $`S_2`$ by Fourier transform, rather than from $`S_2`$ itself ( Figure 1). $`S_2(r)`$ thus scales, in the sense that $`\zeta _2=2\beta /3`$ has its scaling value (cf. equation (24)).
As to $`S_3(r)`$, it is determined from the Von Karman-Howarth relation, equation (23). In this equation the noise term takes in the equal time limit (Novikov-Donsker formalism) the form
$$<(u_1u_2)(f_1f_2)>=2(1/L^2)\underset{k}{}D_0|k|^\beta (1coskr)$$
(26)
The term proportional to ”1” in $`(1coskr)`$ cancels the $`2ϵ`$ in equation (23) because $`(1/L^2)_kD_0|k|^\beta `$ is the total rate of energy input. One thus obtains from equation (23) (in the $`\nu 0`$ limit)
$$\frac{1}{6}dS_3/dr=2(1/L^2)\underset{k}{}D_0|k|^\beta coskr$$
(27)
The ”coskr” term leads by rescaling to the following result
$$S_3(r)r^\beta $$
(28)
for $`1<\beta <0`$, in the case where the noise does not have a cut-off at scales of order $`L`$. ( At $`\beta =1`$ there is an additional logarithm, $`S_3rlogr`$.)
The exponents characterizing the inertial range behavior of $`S_2`$ and $`S_3`$ have therefore their scaling values throughout the domain $`1<\beta <0`$. For $`S_2`$ the result is based on simulations, for $`S_3`$ the expression of the exponent is obtained from the Von Karman-Howarth relation.
(ii) $`S_4,S_5`$ and general $`S_p`$.
For $`p4`$ scaling no longer holds through the entire $`1<\beta <0`$ range. The following are the equations we obtain from the stochastic Burgers equation after isolating the terms which in the inertial range go to zero when the viscosity does, and simplifying the noise terms
$$\frac{1}{6}dS_4(r)/dr=\frac{2}{3}\nu d^2S_3/dr^22<(ϵ_1+ϵ_2)(u_1u_2)>$$
(29)
$`{\displaystyle \frac{1}{40}}dS_5(r)/dr`$ $`=`$ $`{\displaystyle \frac{1}{12}}\nu d^2S_4/dr^2{\displaystyle \frac{1}{2L^2}}{\displaystyle \underset{k}{}}D_0|k|^\beta cos(kr)<(u_1u_1)^2>`$ (30)
$`{\displaystyle \frac{1}{2}}[<(ϵ_1+ϵ_2)(u_1u_2)^2>`$
$`<(ϵ_1+ϵ_2)><(u_1u_2)^2>]`$
$`dS_p(r)/dr`$ $``$ $`{\displaystyle \frac{1}{2L^2}}{\displaystyle \underset{k}{}}D_0|k|^\beta cos(kr)<(u_1u_2)^{p3}>`$ (31)
$`+\mathrm{}<(ϵ_1+ϵ_2)(u_1u_2)^{p3}>`$
The right-hand sides of equations (29) and (30) contain terms ( not written for equation (31)) which go to zero in the small viscosity limit, a noise dependent term and a dissipation rate dependent term. The noise term has the general form
$$\underset{k}{}D_0|k|^\beta cos(kr)<(u_1u_2)^{p3}>\frac{dS_3(r)}{dr}S_{p3}(r)$$
(32)
since $`dS_3(r)/dr_kD_0|k|^\beta coskr`$ (cf. equation (27))
Therefore scaling behavior in $`S_p`$ is present, whether dominant or subdominant, whenever there is scaling behavior in $`S_{p3}`$. Thus the presence of a scaling term in $`S_2,S_3`$ and $`S_4`$ guarantees the presence of one in any $`S_p`$ for $`p4`$. We have already pointed out that both $`S_2`$ and $`S_3`$ scale through the domain $`1<\beta <0`$. The case of $`S_4`$ is trickier because of the absence of an explicit noise term in equation (29). We discuss it below. First we turn to extracting the multifractal behavior of $`S_4`$ and higher order structure functions. This behavior becomes relevant when the associated exponents are smaller than the scaling ones, and therefore the corresponding non-scaling term dominates over the scaling one, since $`r/L1`$.
We first note that in $`k`$-space both $`S_3`$ and $`S_4`$ depend on $`<u(k_1)u(k_2)u(k_3)>,k_1+k_2+k_3=0`$, the first one through its definition, the second one through the $`ϵ`$ dependent term in (29). We thus make the following general ansatz
$$Im<u(k_1)u(k_2)u(k_3)>\frac{|k_1|^{\mu _1}|k_2|^{\mu _2}|k_3|^{\mu _3}}{k_1k_2k_3}+permutations$$
(33)
The constraint that $`S_3(r)r^\beta `$(cf. equation (28)) leads to $`\mu _1+\mu _2+\mu _3=1+\beta `$. We can show that the lowest exponent is obtained when $`\mu _1=\mu _2=\mu _3=\mu /3=(1+\beta )/3`$. Putting the ansatz into the 2nd term of (29) leads to
$$dS_4/dr\nu _{\mathrm{}}^{\mathrm{}}𝑑\alpha 𝑑k_1𝑑k_2𝑑k_3\mathrm{sin}(k_1r)\frac{|k_1k_2k_3|^\mu }{k_1}\mathrm{exp}i\alpha (k_1+k_2+k_3)$$
(34)
Performing the $`k`$ integrals with a cutoff $`\delta `$ and then integrating over $`\alpha `$, with $`0<\mu _1<1`$, one obtains
$$dS_4/dr\nu (2\pi /\delta )^{\mu _2+\mu _3}(1/\delta )r^{\mu _1}$$
(35)
and thus, with $`\mu _1=\mu _2=\mu _3`$,
$$S_4(r)\frac{\nu }{\delta ^{1+2\mu /3}}r^{1\mu /3}$$
(36)
It is important to note here that the non-scaling behavior arises from the term in the equation which involves $`ϵ`$. The expression for $`S_4`$ contains two results:
(i) the fact that in the limit $`\nu 0`$ ,
$$\nu \delta ^{1+2(1+\beta )/3}$$
(37)
whereas in the scaling limit $`\nu \delta ^{1\beta /3}`$. (By writing that at the dissipation scale $`\delta `$, the characteristic eddy time $`t_\delta \delta /\delta ^h`$ is of order of the dissipation time $`\delta ^2/\nu `$, one finds $`\nu \delta ^{1+h}`$) One thus has a new dissipation scale in $`S_4`$, namely $`\delta \nu ^{\frac{1}{1+h_4}}`$. This dissipation scale depends on the corresponding multifractal exponent $`h_4=2(1+\beta )/3`$. For the dominant term this multifractal exponent has to be construed as the one which minimizes $`\zeta _p`$ (cf. (15)).
(ii) second it gives the non-scaling exponent $`\zeta _4=(2\beta )/3`$, which being smaller than the scaling exponent $`\zeta _4=4\beta /3`$ in the region $`1<\beta <2/3`$, dominates over the scaling term in this region.
We now have to get back to the question how scaling behavior arises in $`S_4`$. One can show that it arises through the $`\nu dS_2/dr`$ contribution in $`S_3`$ present in the Von Karman-Howarth relation (cf. equation (23)). It corresponds to $`\mu _1+\mu _2+\mu _3=2+2\beta /3`$ in the ansatz for $`S_3`$ (see above) with however $`\mu _1\mu _2=\mu _3`$.
One can now proceed along the same lines to find the behavior of $`S_5(r)`$, taking as a starting point an ansatz similar to the one used for $`S_4`$, but now for $`<u(k_1)u(k_2)u(k_3)u(k_4)>,k_1+k_2+k_3+k_4=0`$. There are now four $`\mu `$’s, the sum of which is constrained by the known behavior of $`S_4`$ in two different regions $`1<\beta <2/3`$ and $`2/3<\beta <0`$. We know already that in $`S_5`$ because of the presence in equation (30) of the noise term, a scaling contribution will be present. The question that is to be settled through making the ansatz on the 4-point function, is whether there are regions in which the scaling term is subdominant, as happens for $`S_4`$. The answer is yes, and one finds that there are three different regions:
(i) $`1/2<\beta <0`$, where scaling behavior dominates, and thus $`\zeta _5=5\beta /3`$,
(ii) $`2/3<\beta <1/2`$, where $`S_5`$ does not scale, $`\zeta _5=(34\beta )/6`$, and this exponent is smaller than the scaling one and therefore the corresponding term dominates in $`S_5(r)`$,
(iii) $`1<\beta <2/3`$, where $`S_5`$ has still another multifractal exponent, $`\zeta _5=(5\beta )/6`$, which gives the dominant behavior in this region of $`\beta `$. The three exponents connect smoothly at the end points of each interval. In each interval all three terms are present, but the term with the smallest exponent dominates. The first four $`\zeta _p`$’s are shown in Figure 2.
The following general scenario emerges from these results: as $`p`$ increases, simple scaling with $`\zeta _p=p\beta /3`$ occurs over a progressively diminishing range of values for $`\beta `$ close to zero (and negative). Over most of the considered domain therefore, multiscaling occurs as soon as $`p4`$, with the $`\zeta _p`$’s continuous and piecewise linear, the number of linear segments increasing as $`p`$ gets larger. As $`\beta 1`$ all the $`\zeta _p`$’s for $`p3`$ go towards $`1`$. This extreme multifractal regime is a manifestation of the increasingly important role played by shocks as the noise acts on larger and larger scales.
Several remarks are in order here:
(i) if one extracts a fractal scaling exponent for velocity variations from the calculations, as we have done above for $`S_4`$ (equations (12) and (37)), one finds a different value for $`h_5`$ in each of the three regions of $`\beta `$, where different $`\zeta _5`$’s dominate, namely $`h_5=\beta /3`$ for $`1/2<\beta <0`$, $`h_5=1/2+2\beta /3`$ for $`2/3<\beta <1/2`$, and $`h_5=(1+\beta )/6`$ for $`1<\beta <2/3`$. Thus $`h_5`$ is continuous and piecewise linear, and goes to zero as $`\beta 1`$, which is a reflection of the increasing dominance of shocks. The same is true for all $`h_p`$’s with $`p4`$.
(ii) one can also calculate continuous and piecewise linear fractal dimensions $`D(h_p)`$ with the help of equation (15), assuming that the corresponding $`h_p`$ minimizes the right hand side, and using the values of $`h_p`$ and $`\zeta _p`$ which result from the ”ansatz” calculation. One finds that all fractal dimensions tend towards zero as $`\beta 1`$, which again is consistent with the dominance of shock structure.
(ii) we cannot show in general that our calculation based on an ansatz in $`k`$-space, and the assumption of the equality of $`\mu `$’s in $`S_5`$ (cf. equations (33) and (34)) leads to the ”true” dominant behavior in each domain. It is possible that our continuous, piecewise linear $`\zeta _p`$’s, are only an approximation to the ”true” function $`\zeta _p(\beta )`$.
III.2. Dissipation rate correlation and intermittency.
By studying the full equation satisfied by the dissipation rate correlation
$$<ϵ(x+r)ϵ(x)><ϵ>^2(r/L)^\mu $$
(38)
we are able to find expressions for the intermittency exponent $`\mu `$ in terms of static and dynamic exponents of velocity field correlations. Here $`ϵ(x)=\nu (u/x)^2`$ for the Burgers equation and $`ϵ(\stackrel{}{x})=\frac{\nu }{2}(_iu_j+_ju_i)^2`$ for the Navier-Stokes equation are the local dissipation rates. In our previous discussion, we have taken the energy dissipation rate $`ϵ`$ to be a constant, and this is all that is required to obtain Kolmogorov scaling of the structure functions. In this section $`ϵ(\stackrel{}{x})`$ is considered to be a fluctuating quantity which has non trivial correlations, as experiment shows. One still has $`<ϵ(\stackrel{}{x})>=ϵ=constant`$ because of homogeneity.
The following two relations have been proposed for the intermittency exponent $`\mu `$:
$$\mu _1=2\zeta _6$$
(39)
and
$$\mu _2=2\zeta _2\zeta _4$$
(40)
The first one, the most discussed, because experimentally the value of $`\zeta _61.8`$ agrees with that of $`\mu 0.25`$, is essentially obtained by a scaling argument, which uses the dimension of $`ϵ`$, namely $`V^3/L`$, to set $`<ϵ(x+r)ϵ(x)>S_6(r)/r^2(r/L)^{\zeta _62}`$.
The advantage of our approach lies in the fact that relations between $`\mu `$ and structure function exponents $`\zeta _p`$ are derived directly, and simultaneously, from the equation satisfied by the dissipation rate correlation. This equation can be derived from the stochastic Burgers or the Navier-Stokes equation by considering correlations in both space $`r`$ and time $`\tau `$, and then passing to the $`\tau 0`$ limit. In this limit the noise term can be expressed using the Novikov-Donsker formalism. One finds in this way, with $`ϵ_1=ϵ(x+r,t+\tau ),ϵ_2=ϵ(x,t)`$
$`<ϵ_1ϵ_2>`$ $`=`$ $`{\displaystyle \frac{1}{4}}_\tau <(ϵ_1ϵ_2)(u_1u_2)^2>{\displaystyle \frac{1}{6}}_r<(ϵ_1+ϵ_2)(u_1u_2)^3>`$ (41)
$`{\displaystyle \frac{1}{4}}_r<(u_1u_2)^2(ϵ_2u_2ϵ_1u_1)>+{\displaystyle \frac{\nu }{4}}_r^2<(ϵ_1+ϵ_2)(u_1u_2)^2>`$
$`+{\displaystyle \frac{1}{2}}<(u_1u_2)(ϵ_2f_1ϵ_1f_2)>`$
The 3rd term on the right-hand side ensures Galilean invariance together with the first term (the left-hand side is Galilean invariant). The viscosity dependent term, which is connected to $`d^3S_5/dr^3`$ (cf. equation (30)), goes to zero for inertial $`r`$ in the zero viscosity limit.
In order to show again the mathematical similarity of expressions derived from the Burgers and Navier-Stokes equations, we show the equivalent expression in three dimensions derived from equation (2):
$`<ϵ_1ϵ_2>`$ $`=`$ $`{\displaystyle \frac{1}{4}}_\tau <(ϵ_1ϵ_2)(\stackrel{}{u}_1\stackrel{}{u}_2)^2>{\displaystyle \frac{1}{4}}_{r_j}<(ϵ_1+ϵ_2)(u_{1j}u_{2j})(\stackrel{}{u}_1\stackrel{}{u}_2)^2>`$ (42)
$`+{\displaystyle \frac{1}{4}}_{r_j}<(\stackrel{}{u}_1\stackrel{}{u}_2)^2(ϵ_1u_{1j}ϵ_2u_{2j}))`$
$`+{\displaystyle \frac{\nu }{4}}_{r_j}^2<(ϵ_1+ϵ_2)(\stackrel{}{u}_1\stackrel{}{u}_2)^2>+{\displaystyle \frac{\nu }{2}}_{r_i}_{r_j}<ϵ_1u_{2i}u_{2j}+ϵ_2u_{1i}u_{1j}>`$
$`{\displaystyle \frac{1}{2}}_{r_i}<(u_{1i}u_{2i})(ϵ_2p_1+ϵ_1p_2)>`$
$`+{\displaystyle \frac{1}{2}}<(u_{1i}u_{2i})(ϵ_2f_{1i}ϵ_1f_{2i})>`$
Apart from the pressure term and a more complicated viscosity term due to the difference in structure of the definitions of $`ϵ`$ in the Burgers and Navier-Stokes case (see the beginning of this section), the two equations correspond to each other term by term, with an obvious generalization of space indices when going from one to three dimensions.
Now going back to equation (31) with $`p=6`$, one sees that the expression $`<(ϵ_1+ϵ_2)(u_1u_2)^3>`$, which occurs in (41), is precisely the term in $`dS_6/dr`$ which, as argued in section III.1., leads to intermittency. Therefore from (41), $`<ϵ_1ϵ_2>`$ (in the $`\tau 0`$ limit) contains the intermittent behavior $`(r/L)^{\mu _1}`$, with
$$\mu _1=2\zeta _6$$
(43)
as given in equation (39).
As to the first term on the right hand side of (41), one can show that the expression $`<(ϵ_1ϵ_2)(u_1u_2)^2>`$ appears in $`S_4/\tau `$, where it is the only one involving the dissipation rate, and therefore leads to intermittency. There is thus a contribution here to the intermittent behavior of $`<ϵ_1ϵ_2>`$ of exponent
$$\mu _2=z_{4,2}\zeta _4$$
(44)
where $`z_{4,2}`$ characterizes the behavior of the second order partial derivative of $`S_4`$ in time, in the limit $`\tau 0`$. The origin of $`\mu _2`$ is thus dynamical. If simple scaling in time holds, then $`z_{4,2}=2z`$, where $`z=1h`$, with the value of $`h`$ equal to its scaling value. $`z`$ here is the dynamical exponent, not the ”frozen turbulence” exponent of value 1, which characterizes the advection of small structures by large ones. Our preceding result and remarks apply as well to Navier-Stokes turbulence. In this latter case $`z=2/3`$, which is numerically equal to $`\zeta _2`$ (we are going to show in III.3. that this result is general and exact). Substituting $`\zeta _2`$ for $`z`$ (recall that in the scaling limit $`z_{4,2}=2z`$) in (44) leads to the result given in equation (40), which thus appears as a static approximation to what our derivation shows to be the dynamical intermittency exponent given by equation (44).
For the Burgers equation the two intermittency exponents of equations (43) and (44) are the two main ones. For the Navier-Stokes equation we can only assert that these same two occur as well, because our discussion does not take into account the pressure term in equation (42).
III.3. Dynamic behavior.
Except for the discussion of $`\mu _2`$ in the preceding section, our concern up to now has been with the equal time structure functions. We now address the problem of their dynamical behavior. We are interested in relationships between dynamic and static exponents, and also in shedding light on Taylor’s frozen turbulence hypothesis in the case when there is no average flow field. In particular we wish to understand how it happens that the square root of the rms fluctuations of the velocity field replaces the average velocity when the latter is zero, thus allowing ballistic behavior with $`z=1`$ ($`z`$ is defined by $`\tau r^z`$). The objects of our study are now the space and time dependent structure functions
$$S_p(r,\tau )=<(u_1u_2)^p>$$
(45)
where $`u_1=u(x+r,t+\tau ),u_2=u(x,t)`$. The generalization to the three dimensional case is straightforward.
We will concentrate on $`S_2`$. One can derive the following equation from the stochastic Burgers equation
$$S_2(r,\tau )/\tau =\frac{1}{2}T_3/r+<u_1f_2><u_2f_1>$$
(46)
where
$$T_3(r,\tau )=<(u_1+u_2)(u_1u_2)^2>$$
(47)
which apart from additive constants is the same as $`<u_1^2u_2+u_1u_2^2>`$. The term on the left-hand side and the first term on the right-hand side form a Galilean invariant pair. In the $`\tau 0`$ limit $`T_3`$ does not contribute because of symmetry reasons. In this limit there is a discontinuity in the noise term because $`<u_1f_2>`$ contributes for $`\tau >0`$, and $`<u_2f_1>`$ for $`\tau <0`$. One thus has, using equations (23) and (27),
$$S_2(r,\tau =0^+)/\tau =(1/L^2)\underset{k}{}D_0|k|^\beta coskr=\frac{1}{12}dS_3/dr$$
(48)
Assuming simple dynamic scaling for the first time derivative of $`S_2`$ (in the $`\tau 0`$ limit), with $`\tau r^z`$, equation (48) leads to the following relation
$$\zeta _2z=\zeta _31$$
(49)
This equation is the same as equation (21). However here it follows from an exact equation, whereas before it was obtained from a renormalisation analysis. Moreover $`z`$ here is precisely defined as the exponent which characterizes the behavior of the first order partial derivative of $`S_2`$ in time in the limit $`\tau 0`$.
Since $`\zeta _3`$ is known from the Von Karman-Howarth relation (equations (22) or (23)), this equation relates the temporal and spatial exponents which characterize the behavior of the 2nd order velocity structure function. Since $`\zeta _3`$ has its scaling value set by the Von Karman-Howarth relation, any scaling violations in $`\zeta _2`$ has to be compensated by an equal one in $`z`$. Introducing the value of $`\zeta _3`$, one thus has in the case of the Burgers equation
$$\zeta _2z=\beta 1$$
(50)
and in the case of Navier-Stokes
$$\zeta _2z=0$$
(51)
Thus $`\zeta _2`$ and $`z`$ are not independent, the knowledge of one determines the other. This is the first constraint we have found for $`\zeta _2`$, for which none can be found when one limits one’s investigations to static quantities only. In particular, in the Navier-Stokes case $`\zeta _2=2/3=z`$, whereas in the Burgers case one obtains $`z=1+\beta /3`$. The latter results are consistent with the simple Kolmogorov type scaling argument which entails $`z=1h`$.
As $`\tau 0`$ what matters is clearly this dynamical $`z`$, the one appropriate for a Galilean invariant situation. However as soon as $`\tau `$ departs from zero, the ballistic behavior with $`z=1`$ asserts itself. We have checked this numerically for $`S_2(r=0,\tau )`$ and $`S_4(r=0,\tau )`$, for which, if dynamical scaling holds and for example $`S_2(r,\tau )=r^{\zeta _2}g(\tau /r^z)`$, time dependence is of the form $`\tau ^{\zeta _2/z}`$, and similarly for $`S_4`$. Numerically one is able to distinguish satisfactorily between the dynamic and ballistic values of $`z`$ (Figure 3). One thus verifies that as soon as $`\tau `$ is positive, ballistic behavior with $`z=1`$ occurs.
The question now arises in which way ballistic behavior emerges, and with it the use of Taylor’s frozen turbulence hypothesis, in the case when there is no average flow, i.e. $`<u(x,t)>=0`$.
In reference we have shown that if one differentiates relative to $`\tau `$ equation (46), one is lead to the following equation
$$^2S_2(r,\tau )/\tau ^2=<u^2>^2S_2/r^2+\mathrm{}\mathrm{}$$
(52)
The term on the right-hand side is a result of the fact that
$$T_3/\tau <u^2>S_2/r$$
(53)
after use of the assumption that in the $`\nu 0`$ limit the term $`<(u_1u_2)^2(u_{1}^{}{}_{}{}^{2}+u_{2}^{}{}_{}{}^{2})>2<u^2>S_2`$. The latter assumption arises from the observation already made by Tennekes that large scales eddies advect inertial scale information past an Eulerian observer. Here we show that this assumption is encapsulated in the fact that $`S_2(r,\tau )`$ satisfies precisely a wave type equation of characteristic velocity given by the rms fluctuations of the velocity field. One expects this behavior to occur over time scales large compared to the dissipation time and small compared to the turnover time of the large scale structures in the system. A detailed discussion of the other terms occurring in the equation can be found in reference .
IV. Remarks on intermittency.
Before embarking on these remarks one should point out that the nature of turbulence is different for the Burgers and Navier-Stokes equations: for example vortex stretching is believed to be an important ingredient in three dimensional developed turbulence.
Intermittency - the non-scaling behavior of the structure functions in the inertial range - is a hallmark of three dimensional turbulence. The language of multifractality is a convenient way to describe it. What is the origin of intermittency in the statistical behavior of turbulence? The answer is not clear, though intermittency has been connected to the presence of vortex filaments in the flow. In one experiment, where the size of the filaments is several times the dissipation scale, they are associated with events in the velocity field where the velocity derivative has large jumps. This is of course what happens across shocks, which play the role of coherent structures in the one dimensional stochastic Burgers equation. Here one has a clear connection between intermittency and the presence of shocks, though we are unable to give a numerical measure of the number and sizes of shocks. Typically the velocity variation across a shock occurs on length scales of the order of the dissipation scale. For $`\beta `$ negative close to zero, shocks are barely apparent in the velocity profile, and the structure functions show scaling behavior. As $`\beta `$ approaches $`1`$ the shocks play a larger and larger role, and intermittency, the difference between the actual values of the $`\zeta _p`$’s and their scaling values, increases correspondingly (for $`p4`$). At $`\beta 3/2`$ the shocks are present in full, dominating the velocity profile, and intermittency is extreme: all $`\zeta _p`$’s are equal to $`1`$. There is thus an obvious link between the dynamics of shocks - the small scale coherent structures - and intermittency.
We provide two other insights:
\- we connect - not by a self-similarity argument, but from the exact equation - the values of the exponents measuring intermittency in the energy dissipation rate to those measuring intermittency in the velocity structure functions (see III.2.),
\- we show that in the equations for the velocity structure functions the terms responsible for intermittent behavior are those which contain the energy dissipation rate. Intermittent behavior at the inertial scale is thus a consequence of dynamics which occurs at dissipation scales (see III.1.).
For the stochastic Burgers equation we are of course able to provide an extra bonus: namely, with the help of an ansatz, we are able to calculate from the basic equations the low order structure function exponents as $`\beta `$ varies. Such a calculation remains the ”holy grail” for statistical three dimensional turbulence.
Acknowledgments. This review paper was written while F. H. was on sabbatical at the Courant Institute of Mathematical Sciences at New York University. He wishes to thank Dave McLaughlin and Mike Shelley for their hospitality and support.
We are also grateful to Mark Nelkin for a number of comments and suggestions he made concerning the manuscript. We thank the Ohio Supercomputer Center for continuing support.
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# On the Moduli Space of the Localized 1-5 System
## 1 Introduction
Brane moduli spaces, in particular those of black hole configurations, have been investigated extensively over the last few years . One motivation is that in various limits such moduli spaces should be related to the ADS/CFT conjecture . In particular, recent work has focused on cases connected to the elusive relation between $`1+1`$ ADS space and a $`0+1`$ CFT.
A closely related motivation of is the attempt to understand the internal states of black holes in terms of the multi-black hole moduli space. In contrast to the description used in to account for the entropy of such black holes, this new description could be valid in a regime of couplings more properly associated with large classical black holes. Such a description could then lead to better insight into the nature of black hole internal states.
Since in M-theory the black holes themselves are believed to be marginally bound states of various types of branes, one might expect the full moduli space associated with the constituent branes to be relevant to this problem. After all, if states in the near horizon limit of the multi-black hole moduli space are to be interpreted as internal states of a black hole, then the same interpretation naturally applies to all states in the near-horizon limit of the moduli space describing the interactions of the various component branes. As such, it may be important for this program to understand any new features inherent in moduli spaces of localized branes.
We should stress here that by the term “localized” brane we refer to a p-brane whose classical supergravity description is in terms of some p+1 dimensional hypersurface, as opposed to a smeared “fluid” of such objects. It is this set of classical configurations that one would use to define a moduli space on which one could then consider quantum mechanical wavefunctions. While such wavefunctions will of course spread out over the moduli space, this is a different sort of localization/delocalization than we will address in this work.
Now, it is true that when a localized 1-brane approaches a 5-brane there is a certain ‘spontaneous delocalization’ of the 1-brane charges so that the resulting object resembles the black holes studied in . Thus, one might expect a similar result for the moduli space dynamics. However, this is not what we find. Instead, we find the moduli space for a localized 1-brane with longitudinal momentum scattering off such a 5-brane to be substantially different from that of a delocalized 1-brane in the near-horizon limit.
With the exception of the test-brane calculations of , 1/5-brane moduli space calculations in supergravity proceed by dimensionally reducing 10 dimensional solutions to 5 or 4 dimensions so that the branes appear as point particles. The effective action is obtained in the small velocity approximation as shown by Ferrell and Eardley for Reissner-Nordstrom black holes , building on previous work . The moduli space metric is then obtained from the kinetic terms, as first demonstrated for BPS monopoles . In several cases, the moduli spaces have been related to the target space of 1-dimensional supersymmetric sigma models and a connection made between the number of supersymmetries of the effective theory and the complex structure on the moduli space .
In order to get an effective theory of point particles, the brane configuration in 10-dimensions must have appropriate isometries along all possible brane directions. If however, one begins with both 5-branes and localized 1-branes, the 1-branes break translational symmetry along the 5-branes and dimensional reduction in these directions is not possible. The effective theory is necessarily one of extended objects. Localized brane configurations have received considerable interest in recent years, but a study of the scattering of such objects has not yet been carried out. It is the purpose of this work to extend the calculations of to a particular localized system.
In this paper, we calculate the moduli space metric for the system of containing Neveu-Schwarz 5-branes and localized fundamental strings (F1-branes) in the near horizon limit. By S-duality, the moduli space of the localized D1/D5-brane system with momentum will be identical. The solution has a 5-brane wrapped on a $`T^5`$ and a separated 1-brane wrapped on one of the $`T^5`$ cycles. Thus, unlike the delocalized case, there are four extra moduli in the problem labeling the location of the string in the remaining $`T^4`$ directions. This means that there is a single spatial isometry along the cycle on which the string is wrapped and along which a dimensional reduction can be performed. We also include a third charge corresponding to momentum directed along the string. We then calculate the effective action in the low velocity limit following Ferrell and Eardley .
The procedure for computing the moduli space involves first replacing the branes with a smooth ‘dust’ source and then taking the distributional limit where the dust describes a set of branes. We derive the effective action for smooth dust sources in section 2. One of the nice features of this approach is that it is insensitive to the details of the 1-brane singularity. Any distribution of 1-brane charge localized in a region much smaller than the length scales $`L`$ and $`r_5`$ associated with the size of the 4-torus and the 5-brane charge, respectively, will produce much the same results. For large $`L`$, $`r_5`$ the curvatures and dilaton remain small at the 1-brane source, and ten-dimensional supergravity is an adequate description of the system.
The action derived in section 2 can be readily generalized to include many independent dust distributions. In section 3, we consider a special two body case describing a single localized stack of 1-branes and a single stack of 5-branes carrying delocalized 1-brane and momentum charges. We obtain an explicit expression for the moduli space metric in the limit in which the localized branes approach the 5-brane horizon. As with the string metric for a single 5-brane, the metric has a warped product structure, with the transverse radial directions warping the internal $`T^4`$ directions. As mentioned above, the metric in the transverse directions differs from that of the delocalized case. In particular, the transverse metric for relative motion in isotropic coordinates is no longer simply a conformal factor times the standard Euclidean metric. In addition to the terms familiar from the delocalized case, we identify a new one which depends on the ratio $`r_5/L`$. This term is sufficiently large for large $`L`$ to make a finite contribution in the limit $`L\mathrm{}.`$ We close with some discussion in section 4.
## 2 The Effective Action
In the string frame, the ten dimensional action of type IIB supergravity contains the terms
$$S_{10}=\frac{1}{16\pi G_{10}}d^{10}x\sqrt{G}e^{\widehat{\varphi }}[^GR+(\widehat{\varphi })^2\frac{1}{12}\widehat{}^2],$$
(1)
where $`G_{MN}`$ is the 10 dimensional string metric, $`\widehat{\varphi }`$ is the dilaton and $`\widehat{}`$ is the 3-form field associated with an antisymmetric Neveu-Schwarz 2-form field, $`\widehat{B}`$. The symbol $`{}_{}{}^{G}R`$ refers to the Ricci scalar of $`G_{MN}`$ as opposed to that of other metrics that will appear later. A stationary point of (1) represents a solution of type IIB supergravity with all fermions and Ramond-Ramond fields set to zero. The hats on fields serve to simplify the notation later in the paper, after we dimensionally reduce the solution along the single translational symmetry. The overall normalization of the action will not be needed for our purposes.
We are interested in the case of a separated F1-NS5 brane solution, where the one brane is localized in the transverse 5-brane directions. Such a solution was found in , and belongs to the class of chiral null models of . Five of the spatial directions are compactified on a $`T^5`$ on which the 5-brane is wrapped. The 1-brane is wrapped along a single cycle of the $`T^5`$. We employ the coordinates $`(t,z,x^i,y^a)`$, where $`z`$ is the direction along which the 1-brane is wrapped, $`x^i`$ are the 4 spatial directions transverse to the 5-brane, and $`y^a`$ are the remaining 4 directions transverse to the 1-brane along the 5-brane. For simplicity we take the $`T^5`$ to be an orthogonal torus with $`z,y^a`$ labeling the orthogonal directions and with the corresponding cycles having length $`L_z,L`$. In these coordinates the non-vanishing components of the solution are,
$`G_{MN}dX^MdX^N`$ $`=`$ $`H_1^1du(dv+Kdu)+H_5dx^idx^i+dy^ady^a`$
$`e^{\widehat{\varphi }}`$ $`=`$ $`{\displaystyle \frac{H_1}{H_5}}`$
$`\widehat{}_{ijk}`$ $`=`$ $`ϵ_{ijkl}_lH_5\widehat{B}_{uv}=G_{uv}`$ (2)
where, $`H_5(x^i),H_1(x^i,y^a),K(x^i,y^a)`$ are functions associated respectively with the 5-brane charge, the 1-brane charge, and momentum in the z-direction. From , it follows that they satisfy the coupled equations:
$`_i^2H_5`$ $`=`$ $`c_5\rho _5`$
$`_i^2H_1+H_5_a^2H_1`$ $`=`$ $`c_1\rho _1`$
$`_i^2K+H_5_a^2K`$ $`=`$ $`c_K\rho _K,`$ (3)
where $`\rho _5,\rho _1,\rho _K`$ are the brane charge densities for the 5-brane charge, 1-brane charge and momentum respectively. The solution for a 1-brane separated in the transverse directions from a 5-brane was found in (see Eq. (3.1)). Being a chiral null solution, it preserves $`1/8`$ of the supersymmetries, or 4 supercharges. Unlike the solutions considered earlier , neither $`H_1`$ nor $`K`$ are harmonic functions, which is the crucial point of departure for this analysis.
Our choice of convention will be to take $`x^i`$ along the $`(1,2,3,4)`$ directions and the $`y^a`$ along the $`(5,6,7,8)`$ directions. In what follows, we use the collective spatial label $`\alpha =(i,a)`$.
The isometry along the $`z`$ direction makes it possible to dimensionally reduce this solution to 8+1 spacetime dimensions. Proceeding as in , we find the 8+1-d non-vanishing fields $`g_{\mu \nu },A_\mu ^K,A_\mu ^1,\varphi ,\psi ,_{\mu _1\mu _2\mu _3}`$,
$`ds_{8+1}^2`$ $`=`$ $`g_{\mu \nu }dx^\mu dx^\nu `$
$`=`$ $`(H_1H_K)^1dt^2+H_5dx^2+dy^2`$
$`A_0^K`$ $`=`$ $`G_{0z}={\displaystyle \frac{H_K1}{H_K}}`$
$`A_0^1`$ $`=`$ $`\widehat{B}_{0z}={\displaystyle \frac{1}{H_1}}`$
$`\varphi `$ $`=`$ $`\widehat{\varphi }{\displaystyle \frac{1}{2}}\mathrm{ln}G_{zz}=\mathrm{ln}{\displaystyle \frac{H_5}{\sqrt{H_1H_k}}}`$
$`\psi `$ $`=`$ $`\mathrm{ln}G_{zz}=\mathrm{ln}{\displaystyle \frac{H_K}{H_1}}`$
$`_{ijk}`$ $`=`$ $`\widehat{}_{ijk}=ϵ_{ijkl}_lH_5,`$ (4)
where $`H_K=1+K`$, $`dx^2=_{i=1}^4dx^idx^i`$, and $`dy^2=_{a=1}^4dy^ady^a`$. The notation reflects the fact that the potential $`A_\mu ^1`$ couples to the 1-brane charge while $`A_\mu ^K`$ couples to momentum.
As we can see from (2), the 5-brane couples magnetically to the field strength $``$. However, in order to explicitly couple the potentials to sources, it is useful to work in a formalism where the charges all couple electrically to the gauge fields. Now that we have reduced the system to 8+1 dimensions, the 3-form field strength produced by the 5-brane does not couple directly to any other charges. Thus, we are free to consider a dual 6-form field strength and the associated potential. One can check that, in order for the dual 6-form field strength to be an exact form, one must take the dual using the auxiliary metric $`g_{\mu \nu }^{\mathrm{aux}}=(H_1H_K)^{\frac{1}{3}}H_5^{\frac{2}{3}}g_{\mu \nu }`$,
$$_{\mu _1\mu _2\mu _3\mu _4\mu _5\mu _6}=\frac{1}{\sqrt{g^{\mathrm{aux}}}}ϵ_{\mu _1\mu _2\mu _3\mu _4\mu _5\mu _6}^{(\mathrm{aux})\nu _1\nu _2\nu _3}_{\nu _1\nu _2\nu _3}.$$
(5)
For the solution (2), the associated potential $`A_{\mu _1\mu _2\mu _3\mu _4\mu _5}^5`$ for $``$ (i.e., satisfying $`dA^5=`$) then has a single nonzero component,
$$A_{05678}^5=H_5^1.$$
(6)
Finally, since our 5-brane will always remain parallel to the 5,6,7,8 directions (and to the 9 direction, which is hidden in the 8+1 formalism), it is convenient to introduce the notation $`A_\mu ^5=A_{\mu 5678}^5`$, so that $`A^5`$ can be described as a vector potential in parallel with $`A^1`$ and $`A^K`$. Note, however, that the 5,6,7,8 components of $`A_\mu `$ will always vanish. For the static solution we have
$$A_0^5=H_5^1.$$
(7)
The dimensionally reduced 8+1 dimensional action is,
$`S`$ $`=`$ $`{\displaystyle \frac{1}{16\pi G_9}}{\displaystyle }d^9x\sqrt{g}e^\varphi [R+^\mu \varphi _\mu \varphi {\displaystyle \frac{1}{4}}^\mu \psi _\mu \psi {\displaystyle \frac{1}{4}}e^\psi ^{\mu \nu }_{\mu \nu }`$ (8)
$`{\displaystyle \frac{1}{4}}e^\psi ^{\mu \nu }_{\mu \nu }{\displaystyle \frac{1}{2\times 6!}}e^{2\varphi }^2]+{\displaystyle \frac{1}{4}}{\displaystyle }A^K+S_{\mathrm{source}},`$
where $`,`$ are the field strengths associated with the fields, $`A_\mu ^1`$ and $`A_\mu ^K`$ respectively and $`G_9=G_{10}/L_z`$. The source terms are $`S_{source}=S_{matter}+S_{current}`$. The first of these contains the kinetic terms for the branes:
$`S_{matter}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi G_9}}{\displaystyle }dtd^4xd^4y\{c_1\rho _1\sqrt{{\displaystyle \frac{d\tau _1^2}{dt^2}}}e^{\frac{\psi }{2}}c_K\rho _K\sqrt{{\displaystyle \frac{d\tau _K^2}{dt^2}}}e^{\frac{\psi }{2}}`$ (9)
$`c_5\rho _5\sqrt{{\displaystyle \frac{d\tau _5^2}{dt^2}}}e^\varphi ],`$
where the 5-brane kinetic term takes a form like that of a point particle due to our condition that the 5-brane remain parallel to the $`y^a`$ directions. Here $`\tau _1,\tau _K,\tau _5`$ denote the proper time measured along the various branes. The current term is
$`S_{current}`$ $`=`$ $`{\displaystyle 𝑑td^4xd^4y[c_1\rho _1A_\mu ^1v_1^\mu +c_K\rho _KA_\mu ^Kv_K^\mu +c_5\rho _5A_\mu ^5v_5^\mu ]},`$ (10)
where we have taken the matter to be pressureless dust as in , with $`v_1^\mu `$, $`v_K^\mu `$, and $`v_5^\mu `$ the velocities of the 1-brane, momentum, and 5-brane charge distributions (‘dust’) respectively. We take these velocities to be functions of $`t`$ only. Since the $`y^a`$ (i.e., 5,6,7,8) components of $`A_\mu ^5`$ always vanish, the corresponding components of $`v_5^\mu `$ are irrelevant. This is consistent with the fact that only the velocity of the 5-brane in the directions transverse to its world-volume are well-defined. We will continue to represent this velocity as a 9-vector, following our notation for $`v_1^\mu `$ and $`v_K^\mu `$, but with the understanding that we set the $`y^a`$ components of $`v_5^\mu `$ to zero.
In order to regularize the solution, we take $`\rho _I`$ (for $`I=1,K,5`$) to be a smooth function. In the case of $`\rho _5`$, the density will be translationally invariant in the torus directions ($`y^a`$). The limit of localized brane sources leads to the known static solutions. We follow the approach of in first deriving the effective action for smooth sources in the slow motion approximation and then taking the limit of localized brane sources.
These matter sources can be justified either by arguing (as in ) that any smooth source should be able to approximate a black hole, or by noting that $`S_{matter}+S_{current}`$ follows from the relevant parts of the Dirac-Born-Infeld and Wess-Zumino terms in the brane effective actions (see, e.g. ). Viewed in this second way, it is a part of our ansatz that the internal gauge fields are set to zero.
Due to the BPS nature of the branes, the solution (2) is static. The moduli for this system are just the brane’s spatial locations. To calculate the metric on this moduli space we consider the small velocity approximation of in which the forces between the branes remain small. The motion of the branes in this approximation is along geodesics on this moduli space so that its metric can simply be read off from the effective action.
As in , the time reversal symmetry<sup>1</sup><sup>1</sup>1The action (8) with the Chern-Simons term does not have time reversal symmetry, but the original action (1) does have this symmetry. Thus, the equations of motion for the physical fields must be time reversal invariant. can be used to argue that to first order in velocities the perturbations take the form:
$`g_{\mu \nu }dX^\mu dX^\nu `$ $`=`$ $`H_1^1H_K^1dt^2+\delta _{ij}H_5^1[dx^i+N^idt][dx^j+N^jdt]`$
$`+\delta _{ab}[dy^a+N^adt][dy^b+N^bdt]`$
$`\delta A_\alpha ^1`$ $`=`$ $`A_\alpha ^1`$
$`\delta A_\alpha ^K`$ $`=`$ $`A_\alpha ^K`$
$`\delta A_\alpha ^5`$ $`=`$ $`A_\alpha ^5,`$
where $`\alpha `$ runs over the spatial directions. Note that $`H_1,H_K,H_5`$ are now time dependent since the sources in (2) are time dependent. The perturbation in the metric appears as a non-vanishing shift $`N_\alpha `$.
The next step is to compute the $`O(v^2)`$ effective action. As in , one can show that only those $`O(v^2)`$ terms which follow from the above $`O(v)`$ expansion of the fields will in fact contribute to the equations of motion. Other $`O(v^2)`$ terms in the action do not contribute due to the fact that we are near a stationary point of the full action. For this reason, we include below only terms that arise from first order variations in the fields.
Note that an 8+1 split of the spacetime into time and space is inherent in the slow motion approximation. As a result, the fields below will be written with the index $`\alpha `$ that runs only over spatial directions. It is convenient at this point to make a change of conformal frame and to introduce a rescaled metric
$$d\stackrel{~}{s}_8^2=\stackrel{~}{g}_{\alpha \beta }dx^\alpha dx^\beta =H_5^1g_{\alpha \beta }dx^\alpha dx^\beta =dx^2+H_5^1dy^2.$$
(12)
The spatial indices $`\alpha ,\beta `$ will be raised and lowered with the rescaled metric (12).
In introducing this convention, it is important to point out that indices will arise in only two ways. One class of indices come from differential forms such as $`A_\alpha ^1`$, and $`_\alpha H_1`$. In these cases, a covariant placement of the indices is natural and the objects with lower indices $`\alpha ,\beta `$, etc. are simply the pull-back of the spacetime objects (with lower indices $`\mu ,\nu `$, etc.) to the spatial slice. All other indices appear on the velocities $`v_I^\mu `$. For such objects, a contravariant placement of the indices is natural and $`v_I^\alpha `$ represents simply the restriction to the set of spatial components. In contrast, when applying our 8+1 decomposition in the conformal frame (12) to an expression involving $`A^\mu `$ or $`v_\mu `$, one must think carefully about the factors of $`H_5`$. Despite this initial complication, the rescaled metric (12) simplifies the results sufficiently as to make its introduction worthwhile.
A long calculation leads to the $`O(v^2)`$ effective action,
$`S`$ $`=`$ $`{\displaystyle \frac{1}{16\pi G_5L^4}}{\displaystyle }dtd^4xd^4y[{\displaystyle \underset{I\{1,K,5\}}{}}\rho _IH_K\dot{H}_1\dot{H}_5H_1\dot{H}_K\dot{H}_5H_5\dot{H}_1\dot{H}_K`$
$``$ $`{\displaystyle \frac{1}{H_5\mu _𝐍}}{\displaystyle \frac{dP^𝐍dP^𝐍}{2}}+{\displaystyle \underset{I\{1,K,5\}}{}}({\displaystyle \frac{c_I\rho _I\mu _IH_5}{2}}v_I^\alpha v_{I\alpha }`$
$``$ $`{\displaystyle \frac{1}{2H_5\mu _I}}({\displaystyle \frac{H_IdP^IdP^I}{2}}dP^IdP^𝐍)+P_\alpha ^I[_t\left(\stackrel{~}{g}^{\alpha \beta }_\beta H_I\right)+c_I\rho _Iv_I^\alpha ])`$
$`+`$ $`{\displaystyle \underset{I\{1,K,5\}}{}}{\displaystyle \frac{1}{4\mu _IH_5}}\left({\displaystyle \underset{J\{1,K,5\};JI}{}}H_J{\displaystyle \frac{ϵ_{\alpha \beta \gamma \delta }(dP_{\alpha \beta }^IdP_{\gamma \delta }^J)}{2}}ϵ_{\alpha \beta \gamma \delta }(dP_{\alpha \beta }^NdP_{\gamma \delta }^I)\right)`$
$`+`$ $`{\displaystyle \frac{1}{\mu _NH_5}}{\displaystyle \frac{ϵ_{\alpha \beta \gamma \delta }(dP_{\alpha \beta }^NdP_{\gamma \delta }^N)}{4}}],`$
where we have defined
$`\mu _I`$ $`=`$ $`{\displaystyle \frac{H_1H_K}{H_I}},\mathrm{for}I\{1,K,5\}`$
$`\mu _𝐍`$ $`=`$ $`H_1H_K,`$
$`ϵ^{\alpha \beta \gamma \delta }`$ $`=`$ $`ϵ^{0\alpha \beta \gamma \delta 5678},`$ (14)
and we have introduced the one-form fields
$`P_\alpha ^I`$ $`=`$ $`A_\alpha ^I+\mu _IN_\alpha ,\mathrm{for}I\{1,K,5\},`$
$`P_\alpha ^𝐍`$ $`=`$ $`\mu _𝐍N_\alpha ,`$ (15)
with $`dP_{\alpha \beta }^I=_\alpha P_\beta ^I_\beta P_\alpha ^I`$, and $`dP^IdP^I=\stackrel{~}{g}^{\alpha \gamma }\stackrel{~}{g}^{\beta \delta }dP_{\alpha \beta }^IdP_{\gamma \delta }^I`$, etc.
In order to obtain the moduli space metric from this reduced action we need to be able to write down the velocity dependence explicitly. As described in , the values of $`P^I,P^𝐍`$ are to be determined from the constraints, which are in fact given by varying (2) with respect to $`P^I,P^𝐍`$. The existence in the delocalized case of a simple solution to these equations in terms of $`\rho _I`$ and $`v_I^\alpha `$ is rather special.
It is not at all obvious that the same should be true for the localized case. From (2), we see that the equations of motion for the localized case are
$`_\alpha \left({\displaystyle \frac{1}{H_5\mu _I}}[H_IdP^{I\alpha \beta }dP^{𝐍\alpha \beta }]{\displaystyle \frac{ϵ^{\alpha \beta \gamma \delta }}{2H_5\mu _I}}[{\displaystyle \underset{J\{1,K,5\};JI}{}}H_JdP_{\gamma \delta }^JdP_{\gamma \delta }^N]\right)`$
$`=_t\left(\stackrel{~}{g}^{\alpha \beta }_\alpha H_I\right)c_I\rho _Iv_I^\beta ,`$
(16)
$`_\alpha \left({\displaystyle \frac{1}{H_5}}[{\displaystyle \frac{2dP^{𝐍\alpha \beta }}{\mu _𝐍}}{\displaystyle \underset{I\{1,K,5\}}{}}{\displaystyle \frac{dP^{I\alpha \beta }}{\mu _I}}]{\displaystyle \frac{ϵ^{\alpha \beta \gamma \delta }}{2H_5}}[{\displaystyle \frac{2dP_{\gamma \delta }^N}{\mu _𝐍}}{\displaystyle \underset{I\{1,K,5\}}{}}{\displaystyle \frac{dP_{\gamma \delta }^I}{\mu _I}}]\right)`$
$`=0.`$ (17)
The trick to solving such equations is of course to first write the right hand side as the divergence of some antisymmetric tensor. For (2) for the case $`I=5`$ this is straightforward and proceeds along the lines of . One first notes that the right hand side is nonzero only when the index $`\beta `$ takes values in the transverse directions ($`\beta =i`$). One then uses the constraint equation (2) for $`H_5`$. By combining this constraint with current conservation, one arrives at a conservation equation for $`H_5`$ itself:
$$\dot{H}_5+v_5^i_iH_5=0.$$
(18)
By using this result, and also using the constraint for $`H_5`$ to express $`\rho _5`$ in terms of $`H_5`$, one can express the right hand side of (2) for $`I=5`$ as $`_jL^{5ji}`$ where
$$L^{5ji}=\stackrel{~}{g}^{jk}_kH_5v_5^i\stackrel{~}{g}^{ik}_kH_5v_5^j.$$
(19)
For $`I=1,K`$ the equations are more complicated. If one tries to again follow , the root of the problem is that $`H_1`$ and $`H_K`$ are coupled to $`H_5`$ through the constraints (2). Thus, even if the 1-brane isn’t moving, the field $`H_1`$ at a given point will change if we move a 5-brane. The result is that $`H_1`$ and $`H_K`$ do not satisfy simple conservation equations of the form (18).
Nonetheless, one can make progress by introducing a few more potentials. Let us first generalize (19) to $`I=1,K`$ and to an antisymmetric tensor on the full space (thus defining the $`y^a`$ components) through:
$$L^{I\alpha \beta }=\stackrel{~}{g}^{\alpha \gamma }_\gamma H_Iv_I^\beta \stackrel{~}{g}^{\beta \gamma }_\gamma H_Iv_I^\alpha $$
(20)
Note that for $`I=5`$, equation (19) is reproduced with $`L^{5ai}=L^{5ab}=0`$. In addition to (20), we also need the extra potentials
$$𝒥^I^2(_t+v_I^\alpha _\alpha )(H_IH_5),$$
(21)
where $`^2=_i_i_i.`$ Note that $`𝒥^5=0`$ due to the conservation law (18). It is useful to associate with $`𝒥^I`$ a set of antisymmetric tensor fields $`𝒥_{\alpha \beta }^I`$ defined by:
$`𝒥_{ia}^I`$ $`_i_a𝒥`$ $`𝒥_{ai}^I`$
$`𝒥_{ab}^I`$ $`0`$ $`𝒥_{ij}^I.`$ (22)
A bit of calculation then shows that the right hand side of (2) may be cast in the form,
$$_t\left(\stackrel{~}{g}^{\alpha \beta }_\alpha H_I\right)c_I\rho _Iv_I^\beta =_\alpha (L^{I\alpha \beta }+\frac{1}{H_5}𝒥^{I\alpha \beta }).$$
(23)
This allows us to obtain the following solutions:
$`2(1{\displaystyle \frac{ϵ}{2}}){\displaystyle \frac{dP^N}{\mu _𝐍}}`$ $`=`$ $`(1{\displaystyle \frac{ϵ}{2}}){\displaystyle \underset{I1,K,5}{}}{\displaystyle \frac{dP^I}{\mu _I}}`$ (24)
$`(1+{\displaystyle \frac{ϵ}{2}})dP^I`$ $`=`$ $`(1+{\displaystyle \frac{ϵ}{2}})[{\displaystyle \frac{\mu _IH_5}{H_I}}L^I{\displaystyle \frac{1}{H_I}}{\displaystyle \underset{J1,K,5}{}}{\displaystyle \frac{H_5\mu _𝐍}{H_J}}L^J],`$ (25)
where we have used the notation
$$(1\pm \frac{ϵ}{2})^{\alpha \beta \gamma \delta }=(\stackrel{~}{g}^{\gamma [\alpha }\stackrel{~}{g}^{\beta ]\delta }\pm \frac{1}{2}ϵ^{\alpha \beta \gamma \delta }),$$
(26)
and
$$\left[\left(1+\frac{ϵ}{2}\right)A\right]=\left(1+\frac{ϵ}{2}\right)_{\alpha \beta }^{\gamma \delta }A_{\gamma \delta }$$
(27)
for an antisymmetric tensor $`A`$. Using the fact that $`𝒥_{ij}=0=𝒥_{ab}`$, one can check that when the 5-brane charge vanishes, the $`(ij)`$ and $`(ab)`$ components reduce to the results of . The appearance of $`𝒥_{ai}`$ in the $`(ia)`$ components is a novel feature of our calculation.
Inserting the above results into the effective action yields:
$`S`$ $`=`$ $`{\displaystyle \frac{1}{16\pi G_5L^4}}{\displaystyle }dtd^4xd^4y[{\displaystyle \underset{I\{1,K,5\}}{}}\rho _IH_K\dot{H}_1\dot{H}_5H_1\dot{H}_K\dot{H}_5H_5\dot{H}_1\dot{H}_K`$ (28)
$`+`$ $`{\displaystyle \underset{I\{1,K,5\}}{}}\left({\displaystyle \frac{c_I\rho _I\mu _IH_5}{2}}v_I^\alpha v_{I\alpha }\right){\displaystyle \frac{H_1}{2}}(1+{\displaystyle \frac{ϵ}{2}})^{\alpha \beta \gamma \delta }(L^K+{\displaystyle \frac{1}{H_5}}𝒥^K)_{\alpha \beta }(L^5+{\displaystyle \frac{1}{H_5}}𝒥^5)_{\gamma \delta }`$
$``$ $`{\displaystyle \frac{H_K}{2}}(1+{\displaystyle \frac{ϵ}{2}})^{\alpha \beta \gamma \delta }(L^5+{\displaystyle \frac{1}{H_5}}𝒥^5)_{\alpha \beta }(L^1+{\displaystyle \frac{1}{H_5}}𝒥^1)_{\gamma \delta }`$
$``$ $`{\displaystyle \frac{H_5}{2}}(1+{\displaystyle \frac{ϵ}{2}})^{\alpha \beta \gamma \delta }(L^1+{\displaystyle \frac{1}{H_5}}𝒥^1)_{\alpha \beta }(L^K+{\displaystyle \frac{1}{H_5}}𝒥^K)_{\gamma \delta }],`$
where $`G_5=G_9/L^4.`$
Up until this point, for each type of charge $`I=1,K,5,`$ we have allowed for only one dust distribution $`\rho _I`$ with a single constant velocity $`v_I^\mu `$. However, the form (28) provides a ready generalization to the case of many independent dust distributions $`\rho _I^A`$ (representing a different stack of branes for each value of $`A`$) with independent velocities $`v_I^{A\mu }`$ for $`I=1,K`$. Note that we may write $`H_I=1+_A\stackrel{~}{H}_I^A`$ where $`\stackrel{~}{H}_I^A`$ satisfies equation (2) for the source $`\rho _I^A`$ and vanishes at infinity. The linearity of the constraint equations and of the equations of motion for $`dP^I,dP^𝐍`$ then imply that the effective action for the multi-brane case is again of the form (28) with a separate kinetic term ($`\frac{c_I\rho _I^A\mu _IH_5}{2}v_I^{A\alpha }v_{I\alpha }^A`$) included for each brane and with $`L_{\alpha \beta }^I`$ and $`𝒥^I`$ given by:
$`L^{I\alpha \beta }={\displaystyle \underset{A}{}}\left(\stackrel{~}{g}^{\alpha \gamma }_\gamma \stackrel{~}{H}_I^Av^{A\beta }\stackrel{~}{g}^{\beta \gamma }_\gamma \stackrel{~}{H}_I^Av^{A\alpha }\right)`$
$`𝒥^I{\displaystyle \underset{A}{}}^2(_t+v_I^{Ai}_i+v_I^{Ai}_a)(\stackrel{~}{H}_I^AH_5).`$ (29)
The structure here is similar to that of the delocalized case, with the main new feature being the terms of the form $`𝒥_{\alpha \beta }`$.
## 3 The Two-Body Problem
Our task now is to take a limit in which the smooth dust sources become distributions representing some set of localized branes and to then evaluate the effective action (28). The result should yield an action quadratic in velocities associated with geodesic motion through some moduli space.
As one might expect, the fully general case for localized branes is quite complicated. We therefore pick out a special two-body case for detailed analysis. Two-body problems are particularly simple due to the symmetry about the axis connecting the two bodies. This symmetry causes many terms to vanish, and the resulting effective action takes a tractable form. In particular, no term involving $`ϵ^{\alpha \beta \gamma \delta }`$ in (28) will contribute in this case. To see this, note that since $`𝒥_{ij}^I=0`$ we have $`ϵ𝒥^I=0`$. As a result, (28) shows that $`ϵ`$ always appears in the combination $`ϵ_{\alpha \beta \gamma \delta }v^\gamma u^\delta `$ where $`v^\gamma `$, $`u^\delta `$ are the velocities of the two objects. By Galilean invariance, it is sufficient to note that such terms vanish in the center of mass frame where $`v^\gamma `$ and $`u^\delta `$ are proportional.
### 3.1 The setting
Our original goal was to study the scattering of localized 1-branes and 5-branes. As noted above, an object cannot simultaneously carry localized 1-brane charge and 5-brane charge . For this reason, we take one of our two objects to be a stack of localized 1-branes carrying some longitudinal momentum, with $`v^\alpha `$ denoting the velocity of this object. We take the other object to be a stack of 5-branes, which is also allowed to carry (delocalized) 1-brane and momentum (K) charges. We denote the velocity of this object by $`u^\alpha `$. Note that the velocity components $`u^a`$ of such an object around the torus are ill-defined. It is consistent to set them to zero, and we do so for convenience. We refer to the two objects as the localized object ($`𝕃`$) and the delocalized object ($`𝔻`$), where as usual ‘delocalized’ means delocalized along the torus directions.
It is useful to decompose the various $`H_I`$ into parts corresponding to the two objects:
$$H_I=1+H_I^𝔻+H_I^𝕃,$$
(30)
where the $`H_I^{𝕃,𝔻}`$ have corresponding sources $`\rho ^𝔻,\rho ^𝕃`$, and vanish at infinity. Notice that $`H_5^𝕃=0`$. Since the delocalized part is translationally invariant in the $`y^a`$ directions, it satisfies the constraint
$$c_I\rho _I^𝔻=^2H_I^𝔻,$$
(31)
which is independent of $`H_5`$. It therefore obeys the the conservation law (18),
$$\dot{H}_I^𝔻+u^i_iH_I^𝔻=0,$$
(32)
and so does not contribute to the potentials $`𝒥^I`$.
It then follows from (2) that we have the relations
$`L^{I\alpha \beta }`$ $`=`$ $`\stackrel{~}{g}^{\alpha \gamma }_\gamma H_I^𝔻u^\beta \stackrel{~}{g}^{\beta \gamma }_\gamma H_I^𝔻u^\alpha `$
$`+\stackrel{~}{g}^{\alpha \gamma }_\gamma H_I^𝕃v^\beta \stackrel{~}{g}^{\beta \gamma }_\gamma H_I^𝕃v^\alpha `$
$`𝒥^I`$ $`=`$ $`^2(_t+v^\alpha _\alpha )(H_I^𝕃H_5).`$ (33)
After performing several integrations by parts we find the effective action to be
$`S_{eff}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi G_5L^4}}{\displaystyle }dtd^4xd^4y[{\displaystyle }_{I\{1,K,5\}}\rho _I+{\displaystyle \frac{1}{2}}[v^iv^i+v^av^a](c_1\rho _1^𝕃+c_K\rho _K^𝕃)`$ (34)
$`+`$ $`{\displaystyle \frac{1}{2}}u^iu^i\left(c_5\rho _5+c_1\rho _1^𝔻+c_K\rho _K^𝔻\right)+(u^iv^i)(u^iv^i){\displaystyle \frac{c_5\rho _5}{2}}(H_1^𝕃+H_K^𝕃+H_1^𝕃H_K^𝕃)`$
$`+`$ $`{\displaystyle \frac{1}{2}}\left[(u^iv^i)(u^iv^i)H_5+v^av^a\right]\left(c_1\rho _1^𝕃H_K^𝔻+c_K\rho _K^𝕃H_1^𝔻\right)`$
$`+`$ $`{\displaystyle \frac{1}{2}}[(_t+v^i_i)H_5][(_t+v^i_i)(H_1^𝕃H_K^𝕃)]].`$
In the expression above, a sum over $`i,j,a`$ is implicit.
As will become evident in what follows, the key feature of this action is the last term. This term turns out to be rather subtle. Note, however, that it would vanish if $`H_1^𝕃`$ and $`H_K^𝕃`$ were homogeneous on the torus, as in that case $`H_1^𝕃`$ and $`H_K^𝕃`$ each satisfy $`(_t+v^i_i)H_{1,K}^𝕃=0.`$ The integral over the torus in fact allows us to replace both $`H_1^𝕃`$ and $`H_K^𝕃`$ in this term by $`\widehat{H_1^𝕃}`$ and $`\widehat{H_K^𝕃}`$, where the hat indicates that we have removed the homogeneous mode from the Fourier expansion of each function on the torus. It will turn out to be important to note this explicitly. The reason is that, in order to evaluate this final piece in terms of the sources, we will need to write it without explicit time derivatives. In fact, the constraints and the ‘conservation’ of $`H_5`$ can be used to write this last term in the form:
$$\frac{1}{2}[(_t+v^i_i)H_5][(_t+v^i_i)(\widehat{H_1}\widehat{H_K})]=\left[(u^iv^i)\widehat{H_1^𝕃}_iH_5\right]𝒪^1_a^2\left[(u^jv^j)\widehat{H_K^𝕃}_jH_5\right],$$
(35)
where $`𝒪=\stackrel{~}{g}^{\alpha \beta }_\alpha _\beta `$. Now, convergence of the integral of the right hand side turns out to be somewhat subtle when a homogeneous part is included, and depends upon the detailed order in which certain limits are taken. However, by treating the homogeneous part separately and realizing that it will not contribute, we will avoid confusion.
In the above form one can readily take the limit in which the sources become distributions describing the desired branes,
$`c_{1,K}\rho _{1,K}^𝔻`$ $``$ $`16\pi ^2𝔮_{1,K}^𝔻\delta ^4(\stackrel{}{x}\stackrel{}{x}_5)`$
$`c_{1,K}\rho _{1,K}^𝕃`$ $``$ $`16\pi ^2L^4𝔮_{1,K}^𝕃\delta ^4(\stackrel{}{x}\stackrel{}{x}_0)\delta ^4(\stackrel{}{y}\stackrel{}{y}_0)`$
$`c_5\rho _5`$ $``$ $`16\pi ^2𝔮_5\delta ^4(\stackrel{}{x}\stackrel{}{x}_5).`$ (36)
Here $`\stackrel{}{x}_5`$ is the position of the stack of 5-branes and the delocalized 1-branes and momentum, and $`(\stackrel{}{x}_0,\stackrel{}{y}_0)`$ is the position of the localized stack of branes. $`𝔮_{1,K}^𝔻,𝔮_{1,K}^𝕃`$ are the charges of the delocalized and the localized stacks of branes, respectively. Note that from the form of (34) one can see that the details of this limit are unimportant once the branes are localized on a scale much smaller than the typical scale of variation of the functions $`H_I^𝕃`$, $`H_I^𝔻`$, $`H_5`$. Thus, for sufficiently large 4-torus and $`r_5=\sqrt{𝔮_5}`$, replacing the singular perfectly localized brane with a small cloud of well-localized 1-brane charge and momentum charge yields identical results in a regime in which supergravity is a valid description of the system.
Choosing an instantaneous coordinate system centered on the 5-brane, a decomposition into modes along the torus shows that the functions $`H_5,H_1^𝔻,H_K^𝔻,H_1^𝕃,H_K^𝕃`$ are given by
$`H_5`$ $`=`$ $`1+{\displaystyle \frac{𝔮_5}{r^2}}`$
$`H_{1,K}^𝔻`$ $`=`$ $`{\displaystyle \frac{𝔮_{1,K}^𝔻}{r^2}}`$
$`H_{1,K}^𝕃(\stackrel{}{r},\stackrel{}{y};\stackrel{}{r}_0,\stackrel{}{y}_0)`$ $`=`$ $`{\displaystyle \underset{l,\stackrel{}{q},m,n}{}}_{1,K(l\stackrel{}{q})}(r;r_0)D_{mn}^l(\psi _0,\theta _0,\varphi _0)D_{mn}^l(\psi ,\theta ,\varphi )`$ (37)
$`\times e^{i\stackrel{}{q}(\stackrel{}{y}\stackrel{}{y}_0)}`$
where $`(\psi ,\theta ,\varphi )`$ are the Euler angles on $`S^3`$, $`\stackrel{}{q}=\frac{2\pi \stackrel{}{n}}{L}`$, with $`n`$, runs over the momentum lattice of the torus, $`D_{mn}^l(\psi ,\theta ,\varphi )`$ (including both integral and half-odd integral $`l`$) are the rotation matrices which form a complete set of functions on $`S^3`$ (see Appendix), and the radial functions $`_{1,K(l\stackrel{}{q})}(r;r_0)`$ are given by,
$`\mathrm{For}\stackrel{}{q}=0,`$
$`_{1,K(l0)}(r;r_0)`$ $`=`$ $`{\displaystyle \frac{𝔮_{1,K}^𝕃}{(2l+1)}}{\displaystyle \frac{r_0^{2l}}{r^{2l+2}}}r>r_0`$ (38)
$`=`$ $`{\displaystyle \frac{𝔮_{1,K}^𝕃}{(2l+1)}}{\displaystyle \frac{r^{2l}}{r_0^{2l+2}}}r<r_0`$
$`\mathrm{and}|\stackrel{}{q}|0`$
$`_{1,K(l\stackrel{}{q})}(r;r_0)`$ $`=`$ $`2q𝔮_{1,K}^𝕃{\displaystyle \frac{1}{rr_0}}I_\mu (qr_0)K_\mu (qr)r>r_0`$ (39)
$`=`$ $`2q𝔮_{1,K}^𝕃{\displaystyle \frac{1}{rr_0}}K_\mu (qr_0)I_\mu (qr)r<r_0,`$
where $`\mu ^2=1+4l(l+1)+q^2𝔮_5=1+4l(l+1)+\frac{4\pi ^2}{L^2}n^2𝔮_5`$. Note that the homogeneous ($`\stackrel{}{q}=0`$) modes (38) satisfy the naive conservation equation:
$$\left(\frac{}{t}+v\right)\underset{l}{}_{1,K(l,0)}=0.$$
(40)
In contrast, the inhomogeneous modes (39) do not.
This gives $`H_{1,K}^𝕃(0)=\frac{𝔮_{1,K}^𝕃}{r_0^2}`$, $`H_{1,K}^𝔻(\stackrel{}{x}_0)=\frac{𝔮_{1,K}^𝔻}{r_0^2}`$, $`H_5(\stackrel{}{x}_0)=1+\frac{𝔮_5}{r_0^2}`$ so that (34) simplifies to
$`S_{eff}`$ $`={\displaystyle \frac{\pi }{G_5}}`$ $`{\displaystyle }dt\{M+{\displaystyle \frac{1}{2}}[v^iv^i+v^av^a](𝔮_1^𝕃+𝔮_K^𝕃)+{\displaystyle \frac{1}{2}}u^iu^i(𝔮_1^𝔻+𝔮_K^𝔻+𝔮_5)`$
$`+{\displaystyle \frac{1}{2}}(u^iv^i)(u^iv^i){\displaystyle \frac{𝔮_5(𝔮_1^𝕃+𝔮_K^𝕃)+𝔮_1^𝕃𝔮_K^𝔻+𝔮_K^𝕃𝔮_1^𝔻}{r_0^2}}+{\displaystyle \frac{𝔮_1^𝕃𝔮_K^𝔻+𝔮_1^𝔻𝔮_K^𝕃}{2r_0^2}}v^av^a`$
$`+{\displaystyle \frac{1}{2}}(u^iv^i)(u^iv^i){\displaystyle \frac{𝔮_5[𝔮_1^𝕃𝔮_K^𝕃+𝔮_1^𝕃𝔮_K^𝔻+𝔮_K^𝕃𝔮_1^𝔻]}{r_0^4}}\}`$
$`+{\displaystyle \frac{1}{16\pi L^4G_5}}{\displaystyle d^4xd^4y\left[(u^iv^i)\widehat{H_1^𝕃}_iH_5\right]𝒪^1_a^2\left[(u^jv^j)\widehat{H_K^𝕃}_jH_5\right]}`$
where $`M=𝔮_5+𝔮_1^𝔻+𝔮_K^𝔻+𝔮_1^𝕃+𝔮_K^𝕃`$ is the total charge/mass.
That is to say that the action is exactly the same as in the case where both branes are delocalized, except for the inclusion of terms involving $`v^a`$ and the addition of the last term,
$$=\frac{1}{16\pi L^4G_5}d^4xd^4y\left[(u^iv^i)\widehat{H_1^𝕃}_iH_5\right]𝒪^1_a^2\left[(u^jv^j)\widehat{H_K^𝕃}_jH_5\right]$$
(42)
which remains to be evaluated.
### 3.2 The Effective Action in the Near Horizon Limit
It is difficult to obtain an analytic expression from the radial integral in $``$, which involves products of three Bessel functions. However, an explicit result can be obtained in the limit $`r_0<<\sqrt{𝔮}_5`$ in which the localized branes are close to the 5-brane horizon. In this case, the Bessel functions of (39) are approximated by powers of $`r`$.
An important fact is that since the $`\stackrel{}{q}=0`$ modes (38) satisfy the conservation equation (40), they do not contribute to (42) due to the factor of $`_a^2`$ in that term. This fact is true whether or not we have $`r_0\sqrt{𝔮_5}`$. A key feature which becomes apparent here is that $``$ vanishes when either of the localized charges are set to zero. For this special case, the appearance of new moduli in the theory along the $`y`$ directions does not influence the moduli metric in the transverse directions. One expects that this is related to the fact that setting one of the charges to zero doubles the number of supersymmetries.
Explicitly, we have
$``$ $`=`$ $`{\displaystyle \frac{\pi }{(16\pi ^2L^4)^2G_5}}\times 𝔮_1𝔮_K(u^iv^i)(u^jv^j)`$ (43)
$`{\displaystyle }d^4\stackrel{}{r}d^4\stackrel{}{y}[_iH_5(\stackrel{}{r})\widehat{G}(\stackrel{}{r},\stackrel{}{y};\stackrel{}{r_0},\stackrel{}{y_0})\times `$
$`{\displaystyle }d^4\stackrel{}{r^{\prime \prime }}d^4\stackrel{}{y^{\prime \prime }}\left[G(\stackrel{}{r^{\prime \prime }},\stackrel{}{y^{\prime \prime }};\stackrel{}{r},\stackrel{}{y})_jH_5(\stackrel{}{r^{\prime \prime }})_a^2\widehat{G}(\stackrel{}{r^{\prime \prime }},\stackrel{}{y^{\prime \prime }};\stackrel{}{r_0},\stackrel{}{y_0})\right]],`$
where the Green’s function $`G(\stackrel{}{r},\stackrel{}{y};\stackrel{}{r}_0,\stackrel{}{y}_0)`$ satisfies,
$$𝒪G(\stackrel{}{r},\stackrel{}{y};\stackrel{}{r}_0,\stackrel{}{y}_0)=16\pi ^2L^4\delta (\stackrel{}{r}\stackrel{}{r}_0)\delta (\stackrel{}{y}\stackrel{}{y}_0),$$
(44)
and $`\widehat{G}`$ is the Green’s function without it’s homogeneous ($`q=0`$) part.
Expanding in terms of the modes,
$``$ $`=`$ $`{\displaystyle \frac{4\pi }{(16\pi ^2L^4)^2G_5}}𝔮_1^𝕃𝔮_K^𝕃𝔮_5^2(u^iv^i)(u^jv^j)`$ (45)
$`{\displaystyle \underset{\genfrac{}{}{0pt}{}{q_1>0,l_1,m_1,n_1}{\genfrac{}{}{0pt}{}{q_2>0,l_2,m_2,n_2}{q_30,l_3,m_3,n_3}}}{}}q_1^2D_{m_1n_1}^{l_1}(\mathrm{\Theta }_0)D_{m_2n_2}^{l_2}(\mathrm{\Theta }_0)e^{i(\stackrel{}{q}_1\stackrel{}{q}_2)\stackrel{}{y}_0}`$
$`{\displaystyle 𝑑r\left[\widehat{𝔾}_{(q_1l_1)}(r;r_0)𝑑r^{\prime \prime }𝔾_{(q_3l_3)}(r^{\prime \prime };r)\widehat{𝔾}_{(q_2l_2)}(r_0;r^{\prime \prime })\right]}`$
$`{\displaystyle 𝑑\theta 𝑑\psi 𝑑\varphi \left[\mathrm{sin}\theta D_{m_1n_1}^{l_1}(\psi ,\theta ,\varphi )D_{m_3n_3}^{l_3}(\psi ,\theta ,\varphi )b_i(\psi ,\theta ,\varphi )\right]}`$
$`{\displaystyle 𝑑\theta ^{\prime \prime }𝑑\psi ^{\prime \prime }𝑑\varphi ^{\prime \prime }\left[\mathrm{sin}\theta ^{\prime \prime }D_{m_2n_2}^{l_2}(\psi ^{\prime \prime },\theta ^{\prime \prime },\varphi ^{\prime \prime })D_{m_3n_3}^{l_3}(\psi ^{\prime \prime },\theta ^{\prime \prime },\varphi ^{\prime \prime })b_j(\psi ^{\prime \prime },\theta ^{\prime \prime },\varphi ^{\prime \prime })\right]}`$
$`{\displaystyle d^4y\left[e^{i(\stackrel{}{q}_1\stackrel{}{q}_3)\stackrel{}{y}}\right]d^4y^{\prime \prime }\left[e^{i(\stackrel{}{q}_1\stackrel{}{q}_3)\stackrel{}{y^{\prime \prime }}}\right]}`$
where $`\mathrm{\Theta }_0`$ represents the collective angular coordinates of the stack of localized 1-branes, and we have used
$$_iH_5=\frac{2𝔮_5}{r^3}\frac{x_i}{r}=\frac{2𝔮_5}{r^3}b_i(\theta ,\psi ,\varphi ),$$
(46)
which defines $`b_i`$.
The radial integral can be evaluated in the near horizon limit. We need only consider the inhomogeneous ($`\stackrel{}{q}0`$) modes (39) which are approximated by the functions,
$`𝔾_{(ql)}(r;r_0)`$ $``$ $`{\displaystyle \frac{1}{\mu }}{\displaystyle \frac{r_0^{(\mu 1)}}{r^{\mu +1}}}r>r_0`$ (47)
$``$ $`{\displaystyle \frac{1}{\mu }}{\displaystyle \frac{r^{\mu 1}}{r_0^{(\mu +1)}}}r<r_0.`$
This gives us,
$$=\frac{4\pi }{(16\pi ^2)^2G_5}\frac{𝔮_1^𝕃𝔮_K^𝕃𝔮_5^2}{r_0^4}(u^iv^i)(u^jv^j)\underset{q>0,l_1,l_2,l_3}{}q^2\zeta (\mu _1,\mu _2,\mu _3)\times 𝔄_{ij}^{l_1l_2l_3}(\mathrm{\Theta }_0)$$
(48)
where, $`\zeta (\mu _1,\mu _2,\mu _3)`$ is given by
$`\zeta (\mu _1,\mu _2,\mu _3)`$ $`=`$ $`{\displaystyle \frac{4}{\mu _1\mu _2\mu _3}}{\displaystyle \frac{1}{(2+\mu _1+\mu _2)(2+\mu _1+\mu _2)}}`$
$`\times {\displaystyle \frac{(\mu _1+\mu _2)(\mu _1+\mu _3)(\mu _2+\mu _3)(\mu _1+\mu _2+\mu _3)+\mu _3(\mu _1+\mu _2)+4\mu _1\mu _2}{(1+\mu _1+\mu _3)(1+\mu _1+\mu _3)(1+\mu _2+\mu _3)(1+\mu _2+\mu _3)}},`$
and
$`𝔄_{ij}^{l_1l_2l_3}(\mathrm{\Theta }_0)`$ $`=`$ $`{\displaystyle \underset{m_1,m_2,m_3,n_1,n_2,n_3}{}}D_{m_1n_1}^{l_1}(\mathrm{\Theta }_0)D_{m_2n_2}^{l_2}(\mathrm{\Theta }_0)`$
$`{\displaystyle 𝑑\theta 𝑑\psi 𝑑\varphi }\left[\mathrm{sin}\theta D_{m_1n_1}^{l_1}(\psi ,\theta ,\varphi )D_{m_3n_3}^{l_3}(\psi ,\theta ,\varphi )b_i(\psi ,\theta ,\varphi )\right]`$
$`{\displaystyle 𝑑\theta ^{\prime \prime }𝑑\psi ^{\prime \prime }𝑑\varphi ^{\prime \prime }}\left[\mathrm{sin}\theta ^{\prime \prime }D_{m_2n_2}^{l_2}(\psi ^{\prime \prime },\theta ^{\prime \prime },\varphi ^{\prime \prime })D_{m_3n_3}^{l_3}(\psi ^{\prime \prime },\theta ^{\prime \prime },\varphi ^{\prime \prime })b_j(\psi ^{\prime \prime },\theta ^{\prime \prime },\varphi ^{\prime \prime })\right].`$
Note that $`\zeta `$ is symmetric in $`\mu _1,\mu _2.`$ Each component of $`b_i(\psi ,\theta ,\varphi )`$ is a sum over two rotation matrices, so that the angular integration is easily performed using an identity (66) from the Appendix.
In order to evaluate $`𝔄_{ij}(\mathrm{\Theta }_0)`$, we note that there are two determining vectors in the transverse directions: the relative transverse velocity of the stack of localized 1-branes with respect to the 5-brane $`(v^iu^i)`$, and the transverse separation vector $`\stackrel{}{x}_0`$ between the 5 and 1-brane. We may then choose coordinates so that both sets of branes lie in the 1-2 plane.
In order to ease the calculation, we consider the instantaneous frame in which the 5-brane is at the origin, and the stack of localized 1-branes is on the 1-axis. Symmetry about the 1-axis dictates that the off-diagonal part of $`𝔄_{ij}`$ is zero, and that $`𝔄_{22}=𝔄_{33}=𝔄_{44}`$. On the 1-axis, $`\theta _0,\psi _0,\varphi _0=0`$, and thus $`D_{mn}^l(0,0,0)=\delta _{mn}`$ reduces the number of summations. This simplifies the angular part and we find (see Appendix),
$`(u^iv^i)(u^jv^j)𝔄_{ij}^{l_1l_2l_3}(\mathrm{\Theta }_0)`$ $`=`$ $`{\displaystyle \frac{(8\pi ^2)^2}{2}}\delta ({\displaystyle \frac{1}{2}},l_1,l_3)\delta ({\displaystyle \frac{1}{2}},l_2,l_3)`$
$`[(F+H)^{l_1l_2l_3}(v^ru^r)^2(FH)^{l_1l_2l_3}r_0^2(v^\varphi u^\varphi )^2]`$
where $`F^{l_1l_2l_3}`$ and $`H^{l_1l_2l_3}`$ are given by (6) in the Appendix and $`\delta (j,k,l)`$ is the triangle condition,
$`\delta (j,k,l)`$ $`=`$ $`1\mathrm{for}j+kl|jk|,\mathrm{and}\mathrm{j}+\mathrm{k}\mathrm{l}\mathrm{an}\mathrm{integer}`$
$`\delta (j,k,l)`$ $`=`$ $`0,\mathrm{otherwise}.`$ (52)
Thus,
$$=\frac{\pi }{2G_5}\frac{𝔮_1^𝕃𝔮_K^𝕃𝔮_5}{r_0^4}[(𝔣_0+𝔥_0)(v^ru^r)^2(𝔣_0𝔥_0)r_0^2(v^\varphi u^\varphi )^2],$$
(53)
where
$`𝔣_0`$ $`=`$ $`{\displaystyle \underset{q>0,l_1,l_2,l_3}{}}𝔮_5q^2\zeta (\mu _1,\mu _2,\mu _3)F^{l_1l_2l_3}\delta ({\displaystyle \frac{1}{2}},l_1,l_3)\delta ({\displaystyle \frac{1}{2}},l_2,l_3)`$
$`𝔥_0`$ $`=`$ $`{\displaystyle \underset{q>0,l_1,l_2,l_3}{}}𝔮_5q^2\zeta (\mu _1,\mu _2,\mu _3)H^{l_1l_2l_3}\delta ({\displaystyle \frac{1}{2}},l_1,l_3)\delta ({\displaystyle \frac{1}{2}},l_2,l_3).`$ (54)
As noted above, the $`q=0`$ (‘homogeneous’) modes do not contribute to (3.2). Note that for large $`r_5/L`$ the $`n^2=q^2L^2/4\pi ^2>0`$ contributions are highly suppressed by the correspondingly large values of $`\mu `$ even for the lowest term $`n=1`$. Thus, both $`𝔣_0`$ and $`𝔥_0`$ vanish in the limit of large $`r_5/L`$. On the other hand, for small $`r_5/L`$, our $`\mu `$ (and therefore the quantity to be summed in (3.2)) depends only weakly on the integer $`n^2=L^2q^2/4\pi ^2`$ and many terms contribute with equal weight. Thus, $`𝔣_0`$ and $`𝔥_0`$ are correspondingly large in this limit. In section 4 we will discuss in more detail the physically appropriate way to take the $`L\mathrm{}`$ limit. However, for now we simply note that, since $`n`$ appears in $`\mu `$ only through $`\frac{𝔮_5}{L^2}4\pi ^2n^2`$, the growth of $`𝔣_0,𝔥_0`$ for large $`L`$ is of the form $`\frac{𝔮_5}{L^2}_{n^2\frac{L^2}{4\pi ^2𝔮_5}}n^2d^4n\left(\frac{L^2}{𝔮_5}\right)^2.`$
Thus, the effective Lagrangian we obtain to leading order in the near horizon limit is,
$`L_{eff}`$ $``$ $`{\displaystyle \frac{1}{2r_0^2}}\left(𝔮_1^𝕃𝔮_K^𝔻+𝔮_1^𝔻𝔮_K^𝕃\right)v^av^a`$
$`+`$ $`{\displaystyle \frac{1}{2r_0^4}}𝔮_5(𝔮_1^𝕃𝔮_K^𝔻+𝔮_1^𝔻𝔮_K^𝕃+𝔮_1^𝕃𝔮_K^𝕃(1+(𝔣_0+𝔥_0))(v^ru^r)^2`$
$`+`$ $`{\displaystyle \frac{1}{2r_0^4}}𝔮_5(𝔮_1^𝕃𝔮_K^𝔻+𝔮_1^𝔻𝔮_K^𝕃+𝔮_1^𝕃𝔮_K^𝕃(1(𝔣_0𝔥_0))r_0^2(v^\varphi u^\varphi )^2.`$
Note that the zeroth order velocity contribution here is a constant potential equal to the total mass. Hence the dynamics of the system is determined entirely by the geodesics on the moduli space metric.
### 3.3 The Moduli Space
It is useful to cast the effective action in the center of mass coordinates in which the relative velocity is $`\omega ^\alpha =v^\alpha u^\alpha `$. The moduli space metric to leading order in the near horizon limit is therefore,
$`ds^2`$ $``$ $`{\displaystyle \frac{𝔮_5}{r_0^4}}(𝔮_1^𝕃𝔮_K^𝔻+𝔮_1^𝔻𝔮_K^𝕃+𝔮_1^𝕃𝔮_K^𝕃(1+[𝔣_0+𝔥_0]))dr_0^2`$ (56)
$`+`$ $`{\displaystyle \frac{𝔮_5}{r_0^4}}(𝔮_1^𝕃𝔮_K^𝔻+𝔮_1^𝔻𝔮_K^𝕃+𝔮_1^𝕃𝔮_K^𝕃(1[𝔣_0𝔥_0]))r_0^2d\mathrm{\Omega }_3^2`$
$`+`$ $`{\displaystyle \frac{(𝔮_1^𝕃𝔮_K^𝔻+𝔮_1^𝔻𝔮_K^𝕃)}{r_0^2}}dy^ady^a,`$
where $`d\mathrm{\Omega }_3^2`$ is the metric on the unit 3-sphere. Relevant quantities are the total mass $`M=𝔮_5+𝔮_1^𝔻+𝔮_K^𝔻+𝔮_1^𝕃+𝔮_K^𝕃`$, the reduced mass $`𝔪=\frac{(𝔮_1^𝕃+𝔮_K^𝕃)(𝔮_5+𝔮_1^𝔻+𝔮_K^𝔻)}{M}`$, and the center of mass velocity $`V^\alpha =\frac{(𝔮_1^𝕃+𝔮_K^𝕃)v^\alpha +(𝔮_5+𝔮_1^𝔻+𝔮_K^𝔻)u^\alpha }{M}`$. However, the center of mass terms do not appear explicitly in (56) since the center of mass part of the metric is constant in these coordinates and we have restricted our analysis to the leading order $`r_0^4`$ contribution. The metric (56) has a warped product structure, with the transverse radial direction warping the metric in the internal directions. At first sight, it may appear odd that the $`dy^ady^a`$ term does not depend on $`𝔮_5`$. Recall however, that a fundamental string (without longitudinal momentum) should respond to the string metric, at least in the test string approximation. Thus, the above result might be expected from the fact that, in the string frame, the metric for a Neveu-Schwarz fivebrane is simply $`dy^ady^a`$ in the torus directions.
An important question which arises at this stage is whether the metric has singularities. Notice that if $`𝔣_0𝔥_0>0`$, then, by tuning the values of the charges one could make the coefficient of $`d\mathrm{\Omega }_3^2`$ in (56) negative<sup>2</sup><sup>2</sup>2We thank Andy Strominger for pointing this out.. Since (56) describes the leading near-horizon behavior, such an effect could not be compensated by the neglected terms. However, as we demonstrate in the Appendix, this is not the case: namely, $`𝔣_0𝔥_0<0`$.
We have only calculated the metric in the near horizon limit $`r_00`$. There, a change of coordinates $`r_0=\rho ^1`$ illustrates the fact that, as in the delocalized case, small $`r_0`$ is really a large asymptotic region of the relative motion moduli space. The difference between the present case and the delocalized case is that the coefficients of the $`d\rho ^2`$ and $`\rho ^2d\mathrm{\Omega }_3^2`$ terms do not agree. As a result, our transverse moduli space is not quite asymptotically flat in this region. Nonetheless, curvature scalars do go to zero for small $`r_0`$.
More specifically, let us consider the special case when the $`T^4`$ velocity is zero. Again suppose that the motion takes place in some plane so that only one angle $`\varphi `$ on the 3-sphere is relevant. We first rewrite the effective action in terms of the parameters $`\chi `$ and $`\xi `$ as,
$$L_{eff}=\frac{1}{2r_0^4}\chi \dot{r}^2+\frac{1}{2r_0^2}\xi \dot{\varphi }^2.$$
(57)
If the conserved energy and momentum for this system are $``$ and $``$ respectively, the effective radial motion in the near horizon region is governed by the equation,
$$\dot{r}_{0}^{}{}_{}{}^{2}+\frac{2}{\chi }[\frac{^2}{2\xi }r_0^6r_0^4]=0.$$
(58)
The classically accessible regions are those for which the effective potential $`U_{eff}=\frac{2}{\chi }[\frac{^2}{2\xi }r_0^6r_0^4]0`$ and the turning point for the radial motion occurs when
$$r_t=\sqrt{\frac{2\xi }{^2}},$$
(59)
and at $`r_t=0`$. Thus, $`U_{eff}0`$ only for $`r<r_t`$, so that the branes are confined to this region. $`U_{eff}`$, moreover, has a minima at
$$r_m=\sqrt{\frac{2}{3}}r_t.$$
(60)
¿From (59), it would appear that there is a turning point for the motion for any value of the angular momentum, thus at variance with the delocalized case where the branes can sometimes escape to infinity. However, it must be noted that (59) is valid only in the near horizon region: for sufficiently small angular momentum, $`r_t`$ lies outside the near horizon region.
On physical grounds we expect minimal interaction of the objects at large distances, so that at large $`r`$ the metric should be asymptotically flat as in . Thus, black hole scattering should have the familiar qualitative behavior of with a critical impact parameter, depending on the various charges, $``$, and $``$, which separates coalescing orbits from orbits for which the branes escape to relative infinity.
As we noted before, $``$ vanishes when the localized momentum $`𝔮_K^𝕃`$ is set to zero. In this case we need not limit our analysis to the near horizon region and the effective action (3.1) yields the moduli metric,
$`ds^2`$ $`=`$ $`(𝔪+{\displaystyle \frac{[𝔮_5𝔮_1^𝕃+𝔮_1^𝕃𝔮_K^𝔻]}{r_0^2}}+{\displaystyle \frac{𝔮_5𝔮_1^𝕃𝔮_K^𝔻}{r_0^4}})dr_0^2`$ (61)
$`+`$ $`(𝔪+{\displaystyle \frac{[𝔮_5𝔮_1^𝕃+𝔮_1^𝕃𝔮_K^𝔻]}{r_0^2}}+{\displaystyle \frac{𝔮_5𝔮_1^𝕃𝔮_K^𝔻}{r_0^4}})r_0^2d\varphi _0^2`$
$`+`$ $`(𝔪+{\displaystyle \frac{𝔮_1^𝕃𝔮_K^𝔻}{r_0^2}})dy^ady^a.`$
The transverse part of this metric coincides with the results of the black hole calculation for this set of charges.
In particular, when $`𝔮_K^𝔻=0`$, this moduli metric reduces to the particularly simple form,
$$ds^2=[𝔪+\frac{𝔮_5𝔮_1^𝕃}{r_0^2}][dr_0^2+r_0^2d\mathrm{\Omega }_3^2]+𝔪dy^ady^a,$$
(62)
thus reproducing the probe calculation of .
With non-zero localized momenta, however, even to leading order in the near horizon limit, (56) differs from the black hole moduli space calculations of in detail, even though some gross features are preserved. In particular, the transverse moduli space metric for relative motion is no longer conformally related to the standard Euclidean space metric given by the isotropic coordinates. We also note that the coefficients $`𝔣_0`$ and $`𝔥_0`$ are now functions of the ratio $`r_5/L`$. This remains true even when the extra charges $`𝔮_1^𝔻,𝔮_K^𝔻`$ on the 5-brane are set to zero, and hence can be seen as a generic feature of brane localization.
## 4 Discussion
Our results describe the moduli space for a stack of localized 1-branes interacting with a stack of 5-branes. Both branes are allowed to carry momentum in the direction along the 1-brane, and the 5-brane is also allowed to carry 1-brane charge. All charges on the 5-brane are necessarily delocalized along the 5-brane.
If the localized branes are replaced with a system in which either the one-brane charge or the momentum charge is delocalized, the structure of the moduli space simplifies greatly and reduces to previously known forms (e.g. ). When the momentum vanishes, the fact the the moduli space is independent of whether the 1-brane charge is localized might be expected from the (4,4) nonrenormalization theorem described in . That the simple form persists in the presence of delocalized momentum charge is interesting, since momentum charge breaks the same supersymmetries whether or not it is localized.
Another interesting point is that localization affects the structure of the moduli space even in the near 5-brane limit. Recall that when a one-brane approaches a 5-brane, there is a sense in which it ‘spontaneously’ delocalizes. Because of this, one might have expected the moduli space for localized 1-branes to go over to that of delocalized one-branes in the near 5-brane limit. However, this is not the case. The reason for this is that (see ) the one-branes only appear to spontaneously delocalize from the viewpoint of an observer far from the one-brane. When one examines the solutions in the immediate vicinity of the one-brane, it is clear that the one-brane is in fact localized.
Thus, the effective action is sensitive to the region near the one-brane and thus to the localization. It was shown in how the spontaneous delocalization is described in the dual field theory, but it is less clear which field theory observable would encode the fact that the one-brane is localized as viewed by a nearby observer. As a result, it would be interesting to discover how our moduli space metric can be understood from the dual field theory description.
In this near 5-brane limit, we were able to study the structure of the moduli space for this system in some detail. Our results differ from those of previously known, less localized cases , through a modification of the three-charge term<sup>3</sup><sup>3</sup>3Since the coefficient of this term now involves a complicated function of $`𝔮_5/L^2`$, it is not clear that the terminology “three-charge term” is strictly speaking appropriate. Nevertheless, it is a convenient way to refer to this term.. Although our setup is somewhat different, it is interesting to note that a three-charge term is responsible for the puzzle described in . Thus, such terms may warrant further consideration in the future.
Although the three-charge term is modified relative to the less localized case, scattering in the localized moduli space must exhibit the same qualitative behavior as in . ¿From the analysis of , we know that for the delocalized case there is a critical impact parameter beyond which widely separated branes always coalesce. In the near horizon limit, however, we see that this critical impact parameter cannot be calculated.
As the coefficients $`𝔣_0`$ and $`𝔥_0`$ of our three-charge terms are complicated functions of the ratio $`r_5/L`$, it is enlightening to discuss their behavior in various limits. We have seen that they are large ($`L^4/𝔮_5^2`$) for $`r_5L`$ with $`𝔮_1,𝔮_K`$ fixed. For $`r_5L`$ and $`𝔮_1,𝔮_K`$ fixed, the behaviors of $`𝔣_0`$ and $`𝔥_0`$ are controlled by the behavior of $`\zeta `$ for the lowest modes with $`n^2=1`$. Since, $`\zeta `$ scales like $`\mu ^5`$ for large $`\mu `$, we see that $`𝔣_0,𝔥_0\left(\frac{L^2}{𝔮_5}\right)^{\frac{3}{2}}`$ in this limit. Such scalings would correspond to, for example, changing the charge on the fivebrane while holding all other parameters fixed.
Changing the size of the torus, however, is not naturally described by such a limit. Presumably, it is more appropriate to change the size of the torus holding fixed the ten-dimensional parameters. This is equivalent to holding fixed our 9-dimensional parameters as there is no need to rescale the size $`L_z`$ of the remaining circle. As the dimensions of one-brane and momentum charge in ten dimensions is naturally $`(length)^6`$, to hold fixed the ten-dimensional parameters we should scale each of $`𝔮_1,𝔮_K`$ as $`L^4.`$ We should also include the overall factor of $`1/G_5`$ in the effective action, and holding fixed the ten-dimensional Newton’s constant will cause $`G_5`$ to also scale as $`L^4.`$ Taken together with the divergence of $`𝔣_0,𝔥_0`$, we see that our term makes a finite non-zero contribution in this large $`L`$ limit although the usual three-charge term becomes vanishingly small.
On the other hand, for a small torus with ten-dimensional parameters held fixed, our new terms scale as $`L^1`$. While we see that these new terms do become large in this limit, the standard three-charge term in fact scales as $`L^4`$, so that our modification becomes irrelevant.
It would be of interest to understand our moduli space metric as the target space of a supersymmetric sigma model in the spirit of . Although the effective theory we consider includes extended objects, freezing the $`T^4`$ moduli reduces it to one of point particles so that the relevant sigma model will be 1-dimensional as in . Moreover, in a general moduli potential
$$d^9xH_1H_KH_5,$$
(63)
where $`d^9x`$ is the Euclidean measure associated with isotropic coordinates $`x`$, was proposed for a large class of 3-charge brane solutions preserving 4 supercharges. The localization of our charges means that our solution falls outside the class of solutions considered there, but nevertheless it does preserve four supercharges. It is therefore of interest to know how their proposed scheme may be extended to include the localized case. A short calculation shows that a naive attempt to use (63) directly in our context would predict that, in the usual isotropic coordinates, the transverse part of the moduli space metric (56) for single brane scattering to be simply a conformal factor multiplied by the standard Euclidean metric. As discussed in section 3.3, this is not the case <sup>4</sup><sup>4</sup>4Nonetheless, the spherical symmetry of our two-body transverse relative moduli space means that it is conformally flat in different coordinates. This observation allows one to construct a moduli potential, showing that our two-body moduli space is appropriately supersymmetric..
In the introduction, a possible connection was mentioned to the work of which endeavors to associate internal states of black holes with the multi-black hole moduli space. As in their work for the delocalized case, we find an asymptotic region of the moduli space when the branes are nearly coincident. Thus, as one would expect, the moduli space for localized branes also has a continuum of low energy states. In the black hole case, a superconformal structure was discovered which allowed a new choice of Hamiltonian with a discrete spectrum and finite density of states. It would be interesting to know if such a structure arises in this case as well, though we leave this as an open question for the moment. Taking into account the properties of localized branes may lead to further developments for this program.
## 5 Acknowledgments
We would like to thank Sumit Das, Bernard Kay, Jeremy Michelson, Ashoke Sen, and David Tong for helpful discussions. We also thank Andy Strominger and George Papadopoulos for pointing out errors in a previous draft and we thank Jan Gutowski for his patient discussion of calculations associated with the Chern-Simons term. S. Surya was supported by a Visiting Fellowship from the Tata Institute of Fundamental Research. D. Marolf is an Alfred P. Sloan Research Fellow and was supported in part by funds from Syracuse University and from NSF grant PHY-9722362.
## 6 Appendix
We use the Euler angles on the 3-sphere, $`(\psi ,\theta ,\varphi )`$, where $`0\theta \pi `$, $`0\varphi ,\psi 2\pi `$. The transformation to Cartesian coordinates is,
$`x_1=r\mathrm{cos}{\displaystyle \frac{\theta }{2}}\mathrm{cos}{\displaystyle \frac{\varphi +\psi }{2}}`$
$`x_2=r\mathrm{cos}{\displaystyle \frac{\theta }{2}}\mathrm{sin}{\displaystyle \frac{\varphi +\psi }{2}}`$
$`x_3=r\mathrm{sin}{\displaystyle \frac{\theta }{2}}\mathrm{cos}{\displaystyle \frac{\varphi \psi }{2}}`$
$`x_4=r\mathrm{sin}{\displaystyle \frac{\theta }{2}}\mathrm{cos}{\displaystyle \frac{\varphi \psi }{2}}.`$ (64)
To calculate $`𝔄_{ij}`$, in the 1-2 plane, we use the following:
$`b_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}[D_{\frac{1}{2}\frac{1}{2}}^{\frac{1}{2}}+D_{\frac{1}{2}\frac{1}{2}}^{\frac{1}{2}}]`$
$`b_2`$ $`=`$ $`{\displaystyle \frac{1}{2i}}[D_{\frac{1}{2}\frac{1}{2}}^{\frac{1}{2}}D_{\frac{1}{2}\frac{1}{2}}^{\frac{1}{2}}]`$
where $`b_i(\psi ,\theta ,\varphi )`$ is given by (46). The following identity, from , can be used to perform the angular integrals in (3.2),
$$𝑑\theta 𝑑\psi 𝑑\varphi \left[sin\theta D_{mn}^lD_{m_1n_1}^{l_1}D_{m_2n_2}^{l_2}\right]=8\pi ^2\times \left(\begin{array}{ccc}\hfill l& \hfill l_1& \hfill l_2\\ \hfill m& \hfill m_1& \hfill m_2\end{array}\right)\times \left(\begin{array}{ccc}\hfill l& \hfill l_1& \hfill l_2\\ \hfill n& \hfill n_1& \hfill n_2\end{array}\right).$$
(66)
The result is that $`F^{l_1l_2l_3}`$ and $`H^{l_1l_2l_3}`$ in (3.2) are given by,
$`F^{l_1l_2l_3}`$ $``$ $`Q^{l_1l_2l_3}{\displaystyle \underset{m_1}{}}{\displaystyle \frac{(l_2+m_1+1)!(l_1m_1)!}{(l_2m_11)!(l_1+m_1)!}}`$
$`H^{l_1l_2l_3}`$ $``$ $`Q^{l_1l_2l_3}{\displaystyle \underset{m_1}{}}{\displaystyle \frac{(l_2+m_1)!(l_1+m_1)!}{(l_2m_1)!(l_1m_1)!}}[{\displaystyle \frac{(l_3m_1+\frac{1}{2})!}{(l_3+m_1\frac{1}{2})!}}]^2,`$
where,
$$Q^{l_1l_2l_3}=\frac{(\frac{1}{2}+l_1+l_3)!(\frac{1}{2}+l_2+l_3)!}{(l_1+l_3+\frac{3}{2})!(l_2+l_3+\frac{3}{2})!(\frac{1}{2}+l_1l_3)!(\frac{1}{2}+l_2l_3)!(\frac{1}{2}+l_3l_1)!(\frac{1}{2}+l_3l_2)!}$$
(68)
and we have used the Wigner closed expression for the Clebsch Gordon coefficients (see chapter 3, for example).
In order to show that $`𝔣_0𝔥_00`$, we use the fact that the triangle conditions, $`\delta (\frac{1}{2},l_1,l_3)`$ and $`\delta (\frac{1}{2},l_2,l_3)`$ impose rather severe restrictions on the sums over $`l_1,l_2,l_3`$. This allows us to restrict to the four possible cases for each term,
$`\mathrm{Case}\mathrm{\hspace{0.17em}1}:`$ $`l_3=l_1{\displaystyle \frac{1}{2}};`$ $`l_2=l_1`$
$`\mathrm{Case}\mathrm{\hspace{0.17em}2}:`$ $`l_3=l_1+{\displaystyle \frac{1}{2}};`$ $`l_2=l_1`$
$`\mathrm{Case}\mathrm{\hspace{0.17em}3}:`$ $`l_3=l_1+{\displaystyle \frac{1}{2}};`$ $`l_2=l_1+1`$
$`\mathrm{Case}\mathrm{\hspace{0.17em}4}:`$ $`l_3=l_1{\displaystyle \frac{1}{2}};`$ $`l_2=l_11`$
This helps us evaluate,
$`\mathrm{Case}\mathrm{\hspace{0.17em}1}`$ $`:(HF)^{J+1,J+1,J+\frac{1}{2}}`$ $`={\displaystyle \frac{1}{12}}{\displaystyle \frac{(2J+1)}{(2J+3)(J+1)}}`$
$`\mathrm{Case}\mathrm{\hspace{0.17em}2}`$ $`:(HF)^{J,J,J+\frac{1}{2}}`$ $`={\displaystyle \frac{1}{12}}{\displaystyle \frac{(2J+3)}{(2J+1)(J+1)}}`$
$`\mathrm{Case}\mathrm{\hspace{0.17em}3}`$ $`:(HF)^{J,J+1,J+\frac{1}{2}}`$ $`={\displaystyle \frac{1}{12}}{\displaystyle \frac{1}{(J+1)}}`$
$`\mathrm{Case}\mathrm{\hspace{0.17em}4}`$ $`:(HF)^{J+1,J,J+\frac{1}{2}}`$ $`={\displaystyle \frac{1}{12}}{\displaystyle \frac{1}{(J+1)}},`$
where we have put $`l_3=J+\frac{1}{2}`$ for each case. One may check that, the summation over $`l_1,l_2,l_3`$ reduces to a sum over $`J`$ which takes on integer and half odd-integer values, $`J0`$. The corresponding form of $`\zeta (\mu _1,\mu _2,\mu _3)`$ for each case $`s\{1,2,3,4\}`$, is a clumsy expression, which we denote as $`\zeta _s`$.
Now, it suffices to show that $`𝔣_0𝔥_0<0`$ term by term in the sum over $`q`$. Define $`c=1+q^2𝔮_5`$, and
$$𝒰_s(c)=(c1)\underset{J}{}\zeta _s(HF)^s,$$
(71)
where
$$𝔥_0𝔣_0=\underset{q>0}{}\underset{s}{}𝒰_s(c).$$
(72)
In order to test for positivity, we consider a truncated summation up to $`J=100`$ for each case (6), which we refer to as $`\stackrel{~}{𝒰}_s(c)`$. We then plot the various $`\stackrel{~}{𝒰}_s(c)`$’s in Fig(1), as functions of $`c`$. Note that $`c>1`$ for the inhomogeneous modes we are considering, and that $`c=1`$ doesn’t contribute. Finally, plotting $`_s\stackrel{~}{𝒰}_s(c)`$, in Fig(2), we find a distinctively positive function. We have also found that the result given by truncating the series at $`J=5`$ to be essentially the same as that shown below, so that we believe the numerical results to be accurate and the sum to converge rapidly. Figs (3, 4, 5) demonstrate this convergence.
Thus, we conclude that term by term in $`q`$, $`𝔣_0𝔥_0<0`$.
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# Cosmological Evolution Models for QSO/AGN Luminosity Functions: Effects of Spectrum-Luminosity Correlation and Massive Black Hole Remnants
## 1 INTRODUCTION
It has been well established that the observed comoving space density of optically bright QSOs reaches a peak at a critical redshift $`z2`$ (e.g. Peterson 1997, Hartwick & Shade 1990, Weedman 1986). A similar trend is also seen in the X-ray evolution (Miyaji et al. 2000 and references therein). In both cases, the evolution of the luminosity function (LF) is roughly accounted for by the number-conserving luminosity evolution (e.g. Mathez et al. 1976) in which the luminosities of QSOs first gradually decrease at $`z>2`$ from their births near $`z>4`$ and rapidly decline at $`z<2`$. It has recently been questioned whether the X-ray evolution can be adequately described by the pure luminosity evolution (Hasinger 1998, Miyaji et al. 2000). Although the critical redshift at which the QSO activities show a suddden transition is firmly established, it is unclear what determines such a redshift. There have been numerous attempts and models leading to debates as to what physical processes in QSO/AGN engines and/or their surroundings determine the basic characteristics of the cosmological evolution of QSO/AGN (e.g. Caditz et al. 1991, Small & Blandford 1992, Fukugita & Turner 1995, Yi 1996, Haehnelt et al. 1998, Nulsen & Fabian 2000).
We intend to explore most of the existing classes of QSO/AGN evolution in their broad categories. The phenomenological scenarios we explore in this work can be roughly classified as follows. First, we consider a class of models in which a single long-lived ($`10^9`$ yr) QSO population evolves throughout the cosmological time after birth at high redshifts $`z>4`$ (Mathez 1976, Yi 1996, Peterson 1997 and references therein). Second, we study the models in which many short-lived (a few $`10^8`$ yr) QSO populations form and evolve. In these models, the overall observed evolutionary trend is a result of the collective evolution of many different generations of QSO/AGN (Haehnelt et al. 1998). Third, a model of recurrent QSO/AGN activities in which some long-term variabilities of QSO/AGN emission dominates their recurrent activities (Siemiginowska & Elvis 1997). As a specific model of the third type, we assume that the variabilities are mainly caused by the accretion disk instability of the dwarf novae type (Frank et al. 1992) as is occasionally discussed in connection with the origin of high luminosity QSO activities.
We consider two fundamentally different assumptions which are motivated by the recent works in the general area of acrretion disk physics. These assumptions have direct implications on the luminosity evolution of the QSO/AGN and as a consequence they are indirectly testable.
The first one of the two is that the change of accretion flow around a central black hole, which powers the QSO/AGN activities, is mainly driven by the change in mass accretion rate. Such a change then results in the correlated spectral and luminosity evolution of QSO. This assumption is based on the observed behavior of black hole X-ray binaries (BHXBs). Some BHXBs show prominent spectral changes which are strongly correlated with the luminosity changes (e.g. Rutledge et al. 1999, Choi et al. 1999a, 1999b, 2000). Although the details of the models accounting for this behavior vary significantly, it is widely believed that the changes in the mass accretion rate is the underlying cause of this phenomenon. For instance, the ADAF models (Narayan et al. 1998, Yi 1996 and references therein) have enjoyed a relative success in explaining the hard X-ray emitting states of BHXBs while the thin disk models (Frank et al. 1992) have been widely applied to the soft X-ray emitting states. In these models, the continuous changes in the spectral state is interpreted as changes in the accretion rate and the accretion disk’s physical conditions (e.g. Esin, McClintock, & Narayan 1997, and reference therein, Rutledge et al. 1999). We assume that the QSOs’ luminosity evolution is accompanied by a strong spectral evolution caused by the accretion flow transition (see Choi et al. 1999a for details). The spectral state is determined by the BH mass and the physical accretion rate. The dominating effect is obviously that the resulting luminosity evolution shows distinct evolutionary behavior in different energy bands. The nearly mass scale-invariant nature of the accretion flow properties further support this assumption (Narayan & Yi 1995).
The second one is that the QSOs luminosities are interpreted as fixed fractions of the bolometric luminosity and each band’s luminosity is simply given as a fixed fraction of the bolometric luminosity regardless of changes in the accretion rate. The latter approach is closer to those of the conventional studies carried out so far (Peterson 1997). Although this assumption is not clearly supported by any physical models of accretion flows, the lack of detailed QSO spectral information in different bands makes it hard to rule out this simple but convenient approach.
In our previous work (see Choi, Yang, & Yi 1999b, 2000 for details), we looked into the pure luminosity evolution model (Mathez 1976, Peterson 1997) with the explicit inclusion of the spectrum-luminosity correlation. In this model, all QSOs are long-lived and they become gradually dimmer with their comoving space number density conserved throughout the evolution, i.e. from their roughly synchronous births to the present epoch. We arrived at a conclusion that the first assumption, i.e. the accretion flow transition, can be accommodated in the QSO evolution with relatively good fits to the observed luminosity evolution. According to this model, however, the smaller mass black holes with the masses $`<10^8M_{}`$ in galactic nuclei (Magorrian et al. 1998) cannot be direct remnants of the past QSO activities as the QSOs’ black holes grow much more massive than these black holes (cf. Wandel 1999). In this model, the QSO remnants have to exist in the nuclei of rare massive galaxies. Then, the often discussed massive dark objects in ordinary galaxies (Magorrian et al. 1998) should have formed and grown without experiencing the QSO phase, which does not appear to be a popular proposition.
In the present work, we explore a number of QSO/AGN evolution models and make both qualitative and quantitative comparisons among them. The three types of the models we consider are (i) the single long-lived population, pure luminosity evolution model (Yi 1996, Choi et al. 1999b), (ii) the density evolution model in which the QSOs’ evolution is the superposition of multiple generations of short-lived ($``$ a few $`10^610^8`$ yr) QSOs (Haehnelt et al. 1998), and (iii) the model in which the QSO luminosities undergo recurrent variabilities (Siemiginowska & Elvis 1997). The second type of the models have been studied recently (e.g. Haehnelt & Rees 1993, Haehnelt, Natarajan, & Rees 1998) and they are specifically based on the idea that the hierarchical build-up of normal galaxies and the evolution of the AGNs/QSOs are closely connected to each other. It has been claimed that such models are supported by the recent observational evidences such as the existence of massive black holes (BHs) with $`10^{610}M_{}`$ (e.g. Franceschini et al. 1998), and the strong correlation in their masses between the supermassive BH and the spheroidal components of nearby galaxies (e.g. Magorrian et al. 1998, cf. Wandel 1999). In addition, we also explore a specific physical model in which the QSOs’ long-term variabilities caused by accretion disk instabilities contribute considerably to the observed QSO LFs (Siemiginowska & Elvis 1997). We use these classes of models and make comparisons among them using the derived LFs and MFs along with the observational data.
In sections 2 and 3, we summarize the evolution models and describe how we determine model parameters in each model. We derive the resulting analytical LFs of QSO/AGN and MFs of BH remnants. In section 4, we draw our conclusions and discuss their implications on how to interpret the observational data.
## 2 The Multiple Population Model for Short-Lived QSOs
In this type of scenario, the formation and evolution of normal galaxies occur within a hierarchical merging of dark matter halos (Peebles 1993), which is closely connected to the formation and evolution of QSOs experiencing short active phases (Haehnelt & Rees 1993, Kauffmann & Haehnelt 2000, Monaco et al. 2000). It has been generally assumed that a QSO is born (e.g. Nulsen & Fabian 2000 and references therein) or re-activated (e.g. Small & Blandford 1992) when two galaxies merge and that mergers provide the fuel for the newly formed central massive BH. It remains unclear whether during the short, active emission phase, the QSOs’ emission spectra show any rapid spectral evolution while their luminosities rise and fall on time scales much shorter than the cosmological evolution time scale. In other words, one could ask how the assumed spectral evolution, which we mentioned above, affects the QSO LFs in this scenario. In the multiple population models, the cosmologically evolving QSO population is composed of many short-lived generations and undergoes successive rise and fall of QSO activities. If the luminosity-spectrum correlation is applied to the QSO generations, the evolution of an individual short-lived QSO should exhibit some appreciable changes in the spectral emission state with the decreasing mass accretion rate or luminosity after merge-driven trigger of the QSO activity. One of the most significant constraint on this scenario comes from the observed LFs (e.g. ROSAT samples by Miyaji et al. 2000) and the MF derived from the radio LF of E/SO galaxies cores (e.g. Salucci et al. 1999). The AGN LFs in the soft X-ray band have been updated by more extensive analyses with more recent and expanded ROSAT Bright Survey and ROSAT Deep Survey (Miyaji et al. 2000). Together with the earlier results reported by Hasinger (1998), these LFs show an apparent excess at the faintest soft X-ray luminosities $`<10^{42}`$ erg s<sup>-1</sup> at the redshift epoch of $`z=0.00.2`$. In this work, we do not consider this excess in deriving the best fits based on the fact that at such a low luminosity level the distinction between QSO/AGN and less active bright galaxies and/or Seyferts is quite obscure (e.g. Yi & Boughn 1998, 1999 and references therein). For these low luminosity AGN, the observed low-energy X-ray background (XRB) could provide an integral constraint on their LFs (e.g. Franceschini et al. 1999).
We describe the several assumptions and parameters involved in deriving the LFs and constructing an evolution model for QSOs. We first need to specify an individual QSO activity in terms of its spectra and luminosities (see, Choi, Yang, & Yi 1999b for details). We assume that each QSO begins their activity with a newly formed black hole of mass, M, and gas accretion flow with a rate, $`\dot{M}`$, which shows a generally decreases as the QSO evolves. The rate decrease is taken for simplicity to be of the exponential form with the characteristic e-folding time scale $`t_{evol}`$ (e.g. Haiman & Menou 2000)
$$\dot{M}=\dot{M}_oexp\left(t/t_{evol}\right)$$
(1)
where $`t=0`$ is taken to be the QSO trigger time (e.g. Yi 1996). The $`t_{evol}`$ is taken as a fraction of the cosmic time $`t_{age}10^{10}`$ yr for a flat universe (i.e. $`q_o=0.5`$) with no cosmological constant and $`H_o=50`$ km s<sup>-1</sup> Mpc<sup>-1</sup> and essentially corresponds to the lifetime of a QSO. The accretion time scale $`t`$ scales with redshift
$$t(1+z)^{1.5}$$
(2)
and the accretion rate,
$$\dot{m}=\dot{M}/\dot{M}_{Edd}\dot{M}/M$$
(3)
is expressed in units of the mass-dependent Eddington accretion rate.
### 2.1 The Spectral Evolution Model for Multiple QSO Populations (SEM)
As $`\dot{m}`$ decreases after a merging or some other QSO trigger event, QSOs experience two types of spectral transition, from “Very High” state (VHS, slim disk, and $`\dot{m}>1`$) to “High” state (HS, thin disc, and $`0.01\dot{m}1`$) and subsequently, from “High” state to “Low” state including “Off” state (LS/OS, ADAF, and $`\dot{m}<0.01`$). These transitions are caused by the physical changes in the accretion flows and the accompanied luminosity change is strongly correlated with the spectral state change (Yi 1996, Choi et al. 1999a, 1999b, 2000, Narayan et al. 1998). This essentially describes the main ingredients in the assumption of spectral evolution. Based on this simple assumption, we show the expected spectral energy distributions in various spectral states of a QSO with a BH mass of $`10^8M_{}`$ which include the emission from ADAFs (see Figure 1(a)). The hard X-ray emitting ADAF emission has been pointed out as a possible source of the diffuse X-ray background (Yi & Boughn 1998, Fabian & Rees 1995, Di Matteo et al. 1999). Figure 1(b) shows the expected luminosity evolution of a long-lived QSO with the evolution time scale
$$t_{evol}=0.50t_{age}6.4\times 10^9yr$$
(4)
where the initial BH mass is $`10^8M_{}`$ and the initial $`\dot{m}=1`$. The different energy bands show different luminosity evolution trends (Choi et al. 1999b), which suggest that the spectral evolution could be distinguishable from the simple non-spectral evolution case in which the luminosity in each band is simply proportional to the bolometric luminosity (see a dash-dotted line in Figure 1(b)). The transition of accretion flows from the thin disk type (HS) to the ADAF type (LS/OS) occurs at the critical accretion rate $`\dot{m}_c=0.3\alpha ^2`$ where $`\alpha =0.3`$ (Yi 1996, Narayan & Yi 1995) and such a transition makes QSO luminosities decrease/increase rapidly at $`z<1`$. In particular, the LFs in the hard X-ray band are affected differently by the spectral change (from HS to LS/OS), which is quite distinct from the way other energy bands behave. The bolometric luminosity evolution is primarily determined by the redshift dependence of the evolutionary time scale defined above. That is, the rapid decline of the QSO’s bolometric luminosity at low redshifts is caused by the particular redshift dependence of the evolution time scale we have adopted.
We must also specify the several major parameters necessary for the construction of the LF in the SEM model. The main physical parameters are: (i) the merger rate of host galaxies or the birth rate of QSOs, (ii) the life time of an individual QSO, and (iii) the mass distribution of seed BHs at birth. The birth rate in this work is equivalent to the formation rate of new seed BHs at different redshifts. We assume that the birth rate is a simple function of the redshift for which we adopt a simple functional form. For simplicity, we adopt the power-law of the form
$$N(z)=N_o[(1+z)/(1+z{}_{o}{}^{})]^s$$
(5)
where $`N_o`$ is the number of QSOs belonging to the first generation at the epoch, $`z=z_o`$, and $`s`$ determines the basic evolutionary trend of the birth rate as a function of the redshift.
We note that all of the characteristic parameters we find in each evolution model can cause substantial differences to the QSO number density and LF. The best fit values of the parameters are determined by comparing the resulting LF fits in the soft X-ray band with the observed LFs of AGN in the same band. Simultaneously, the theoretically predicted mass distributions are compared with the estimated MF from radio data of E/S0 galaxies. In this sense, the mass distribution of QSO remnants provides an additional constraint. Such a constraint is of course invalid if a significant fraction of black holes at galactic nuclei are not formed from QSO activities.
Figure 2 clearly shows the dependence of the QSO LF on the three main parameters. In other words, in the SEM models considered, the LFs are essentially determined by the number density in each QSO generation, the number of generations, g, constituting QSOs population, and the evolutionary time scale. QSO LFs shown in this figure refer to those in the soft X-ray energy band, which is identical to the observed soft X-ray band. We assume that the initial accretion rate $`\dot{m}`$ of an individual QSO dose not exceed the Eddington limit, $`\dot{m}=1`$. The comoving number density is normalized to match that of Miyaji et al. (2000) in panel (b). Panel (a1) gives the best fit model which is in good agreement with the observed LFs (see panel (b)). This particular model assumes with the simple power-law evolution law ($`s=1`$) with a broad initial mass distribution of seed BHs, and 100 QSO generations over the range of $`0.18z4.3`$. The $`t_{evol}`$ for all QSOs is given as $`0.03t_{age}3.9\times 10^8`$ yr. This is in the lifetime range ($`10^610^8`$ yrs) of the luminous phase of QSOs suggested by Haiman & Hui (2000). The initial masses of seed BHs in each generation are randomly chosen by a single power law with a slope of $`2.5`$ and their masses are spread from $`10^{710}M_{}`$ in the first high $`z`$ generation to $`10^{69}M_{}`$ in the last low $`z`$ generation. Although these assumptions are rather arbitrary, the best fit model is clearly contrasted with less successful models in fitting the observed LF in panel (b).
Panel (c1) shows the predicted MF of BHs at six redshifts obtained from LFs (in panel (a1)) which are in turn derived in the SEM model and panel (c2) illustrates that in the single QSO population model (hereafter SES model and see below). The solid line with the cross symbols in panel (c1) and (c2) represents the estimated MF from radio data of E/SO galaxies in Salucci et al. (1999). According to this model, just a few % of the present-day galaxies with the comoving number density $`10^2`$ h$`{}_{50}{}^{}_{}^{3}`$Mpc<sup>-3</sup> (Magorrian et al. 1998) would have passed through the QSO phase (cf. Wandel 1999). On the other hand, due to the longevity of the QSOs and the prolonged accretion phase, it is difficult for a single populations model to explain the low mass BHs with masses of $`10^{67}M_{}`$ in nearby spiral galaxies (Choi et al. 1999b, 2000 and references therein). In the single population model, the massive dark objects in the nearby galactic centers cannot have passed through the long QSO phase (see panel (c2)). It is still possible that some rare massive black holes in giant elliptical galaxies residing in galaxy clusters (Fabian & Canizares 1988, Mahadevan 1997) could have grown through a long-lasting QSO phase. As the number of QSO generations in the SEM model decreases, the resulting LF become similar to those of the single population model as anticipated. The apparent transition from SEM to the single population model occurs when the number of QSO generations falls to $`10`$ and the QSO evolution time scale $`t_{evol}3.9\times 10^9`$ yr (in panel (a4)).
One of the distinguishable conclusions concerning this model is that the number density of bright QSOs is not a constant but shows a strong evolution over the cosmological time scale (Miyaji et al. 2000, Wisotzki 2000a). The evolutionary time scale, $`t_{evol}`$, and its redshift dependence play an important role in causing this feature. In essence, a shorter $`t_{evol}`$ allows $`\dot{M}`$ to decline more rapidly and the BHs to gain less mass. For instance, for $`t_{evol}=0.03t_{age}3.9\times 10^8`$ yr, remnant BH masses are bigger than the initial masses by a factor $`10`$. As a result, QSOs experience the sudden changes of the spectral states and a large number of QSOs at any given redshift become too faint to be observed (i.e. $`<10^{42}`$ ergs s<sup>-1</sup>). This sudden decrease in the number of luminous QSOs gets more marked at lower redshifts due to the particular redshift dependence of $`t_{evol}`$. Therefore, unless the power $`s`$ of the power-law density evolution reverses its sign from positive to negative (from $`+1`$ to $`1`$), the evolution trend of number density in the derived LF cannot match the observed one. In short, the observed LF evolution requires that the number of QSOs formed at lower redshifts be higher than that at higher redshifts (in panel (a1) and (a2)). The parameters used in panel (a2) except for the power $`s`$ are as given in panel (a1). Panel (a3) shows that the derived LFs do not evolve significantly if QSOs formed in all generations have the same initial mass distribution in the range of $`10^6`$ to $`10^9M_{}`$. By comparing panel (a1) with panel (a3), we see that when the QSOs with less massive seed BHs are born and refueled more abundantly at lower redshifts, the resulting LFs in the SEM model are generally in better agreement with the observed LFs.
Panels (a1), (a4), (a5), and (c1) show the effects of the number of generations and the evolutionary time scale which directly affect the LF evolution and the MF of the BH remnants. The parameters except for the $`t_{evol}`$ and the number of generations used in panels (a4) and (a5) are identical to those given in panel (a1). If the number of QSO generations is small, the QSO evolution could show some discontinuous evolutionary behavior. In this case, the birth and death of a generation of QSOs could show up as a discrete evolutionary episode. This type of the discrete evolutionary behavior is smeared out when the number of QSO generations is large enough or the lifetime of QSOs is long enough. For instance, when QSOs in each generation evolve with $`t_{evol}3.9\times 10^8`$ yr, $`100`$ generations are enough to ensure that the QSO evolution appears smooth and continuous.
### 2.2 Constraint from Mass Function of Supermassive Black Hole Remnants
The MF provides an additional constraint on the QSO evolution. In panel (c1), the solid, dash-dotted, and dashed lines (without symbols) correspond to the predicted MFs of BH remnants in the models with the LFs in panel (a1), (a4) and (a5), respectively. Because QSOs with a short lifetime do not substantially grow in mass during accretion, in order to match the comoving number density implied by the BH remnant distribution in nearby galaxies, the shorter QSO lifetime demands the higher QSO number density. Panel (c1), however, shows that such a compensation can be ruled out in certain cases based on the MF. Assuming that the QSOs shine with an X-ray luminosity efficiency of 10% of the bolometric luminosity
$$L=\eta \dot{M}c^2$$
(6)
with the bolometric radiative efficiency $`\eta 0.1`$ during QSO phase, the predicted MF can not account for the high comoving mass density implied by the nearby galaxies (Salucci et al. 1999).
We find that the comoving number densities of the BH remnants inferred from the best-fitted soft X-ray LF (in panel (a1)) are smaller than those estimated for the putative BHs in the bulges of nearby galaxies (e.g. Magorrian et al. 1998, Salucci et al. 1999) by a factor of about $`100`$. We also find that the same number densities are smaller by a factor of about $`10`$ than those implied by optically bright QSOs under the assumption that they are powered by accretion onto supermassive BHs with an assumed accretion efficiency of $`10\%`$ (e.g. Phinney 1997, Haehnelt et al. 1998 and reference therein). These differences are significant and there do exist some discrepancies between our estimates and the other earlier estimates. Most notably, we use the soft X-ray LF to match the MF of BHs while other studies adopted the B-band LF to match the MF of massive dark objects. In order to allow the predicted MFs to satisfy observational constraint, a number of conditions have to be met. First, the lifetime of sample QSOs must be short (a few $`10^6`$ to $`10^7`$ yr) and the QSOs have to be already massive enough (about $`10^8`$ to $`10^{11}M_{}`$) before accretion-driven mass growth occurs. Such massive black holes could form without emitting detectable photons if their pre-QSO evolution occurs in the form of the super-Eddington accretion with very low radiative efficiencies (Haehnelt et al. 1998). Second, the X-ray fraction of the bolometric luminosity, $`f_{XR}`$, and/or the bolometric radiative efficiency, $`\eta `$ must be smaller than that previously assumed (Haehnelt et al. 1998). For instance, in the absence of any spectral evolution in the context of the multiple QSO population model (hereafter NEM model), the best-fit value of $`f_{XR}\eta `$ is smaller than the usually assumed value by a factor of about 100.
### 2.3 No Spectral Evolution Model with Multiple Populations (NEM)
In Figure 3 (a1), we plot the soft X-ray LF giving the best-fit MF of BH remnants (dashed line in Figure 3 (c)) for the model of no-spectral evolution of multiple QSO population (hereafter NEM model). In Figure 3 (a2) and (c), the best-fit soft X-ray LF and the MF (dash-dotted line) corresponding to this LF are shown. The resulting LF in Figure 3 (a1) is in very poor agreement with the observed LF in Figure 3 (b). Figure 3 demonstrates the difficulty of reconciling the LFs and MFs and fitting them simultaneously.
### 2.4 Discrepancies between MF and LF
There are several possibilities which could potentially account for the discrepancies between the LF and the MF. These possibilities essentially rely on the formation and accretion history for supermassive BHs.
First, the mass density derived from the Massive Dark Objects (MDOs) mass estimates (e.g. Magorrian et al. 1998) could be grossly overestimated due to some systematic uncertainties in the estimates of $`L_{sph}/L_{tot}`$, $`M_{MDO}/M_{sph}`$, etc. (e.g. Salucci et al. 1999). Second, it might be unreasonable for all galaxies to contain MDOs which are actually supermassive BHs. In other words, just a few % of the present-day galaxies with the comoving number density $`10^2`$ h$`{}_{50}{}^{}_{}^{3}`$Mpc<sup>-3</sup> (Magorrian et al. 1998) could have passed through the QSO phase powered by accretion on to supermassive BHs. Third, the low optical emission efficiency could be responsible for low number densities observed in the LFs (Haehnelt et al. 1998). Haehnelt et al. (1998) speculated on two possibilities based on the possibility of the low optical emission efficiency. One of them is the existence of a population of obscured QSO by dust based on the unified Seyfert scheme (e.g. Madau et al. 1994, Comastri et al. 1995, Hasinger 2000, Fabian 1999) and the other is the existence of a population of galaxies undergoing ADAF-type accretion (Narayan & Yi 1995, Fabian & Rees 1995, Yi 1996). We emphasize that although we do not use the B-band LF but the soft X-ray LF to match the observed soft X-ray LF, our treatment of accretion and QSO emission does include the evolution through the ADAF phase (via spectral evolution). Despite this improvement which addresses the low efficiency ADAF flows, the resulting comoving number density of BH remnants is still far short of the required mass density of BHs in nearby galaxies.
The gas and dust obscuration could affect the observed QSO luminosities not only in the optical band but also in the soft X-ray band (Hasinger 2000, Fabian 1999). The effects of obscuration in QSO number counting could be significant and its indirect evidence could be found in the fact that the space density of obscured AGN is about 3 times higher than that of unobscured AGN in the X-ray background model (Comastri et al. 1995, Hasinger 1998). The effects of obscuration on emission could differ significantly from band to band. The space density obtained for the soft X-ray of AGN by Miyaji et al. (2000) exceeds the one observed by Boyle et al. (1991, 2000), which suggests that the optical and X-ray LFs evolve quite distinctly (cf. Hatziminaoglou et al. 1998, Wisotzky 2000b). Therefore, we can plausibly conclude that the effects of the low efficiency accretion as discussed by Haehnelt et al. (1998) are at most marginally significant whereas the possibility of obscured accretion and emission appear to be more important in reconciling the QSO LFs and MFs.
On the other hand, due to the long accretion phase and huge mass gains of QSOs, it is difficult for a single population model to account for the smaller mass BHs with masses of $`10^{67}M_{}`$ in nearby spiral galaxies in terms of the QSO remnants which have evolved through the QSO phase (see Figure 2 (c2)). It is also difficult to rule out the possibility that only the more massive black hole remnants have gone through the QSO phase and these giant QSO remnants are found in present-day giant elliptical galaxies. LF in the SEM model with the generations of $`10`$ and $`t_{evol}`$ of $`3.9\times 10^9`$ yr (in Figure 2 (a4) ) is already nearly indistinguishable from one in the single population model with spectral evolution.
### 2.5 Comparison among Various Evolution Models
Figure 4 makes a comparison among the derived LFs in various energy bands including hard X-ray energies ($`210`$ keV), soft X-ray energies ($`0.52`$ keV), and $`4400\AA `$ for all the QSO evolution models we test. The evolution parameters in each model are determined by requiring them to make the derived the soft X-ray LF in good agreement with the observed LFs of AGN in the soft X-ray band (Miyaji et al. 2000). Such a fitting also fixes the overall normalization of the comoving number density. The best-fit parameters for various models are shown in the Table 1.
The MFs of BH remnants in Figure 5 result from these best-fitted LFs shown in Figure 4. Figure 4 and 5 clearly demonstrate the difficulties in satisfying the LF and MF constraints simultaneously. In Figure 4, panel (a1) and (b1) show the derived LFs in various energy bands for SES and SEM models, respectively. It appears that the SEM model gives the LFs in the soft X-ray energy band which are in better agreement with the observed number densities and luminosities evolution of QSOs than the SES model. Such an assessment critically relies on the fact that at $`z=1.62.3`$ the number densities of the observed bright QSOs (Miyaji et al. 2000) apparently decrease only very gradually.
It is generally difficult to interprete specific best-fit values of the QSO evolution parameters in the multiple population model although in the hierarchical galaxy formation scenario the QSO evolution occurs as a part of the general galaxy formation. In particular, the duration of each QSO generation and the number of QSO generations over the cosmological age of the universe are not well constrained. We find that just as in the SES model, the SEM model also gives different LF evolution trends in different energy bands, which is clearly a distinguishable and potentially observable feature caused by the strong spectral evolution correlated with the luminosity evolution. The optical ($`4400\AA `$) LF tends to evolve more rapidly than the X-ray LF at lower redshifts and the B band LF of QSOs observed by Boyle et al. (1991) in Figure 4. At lower redshifts, the adopted spectral evolution prescription (Choi et al. 1999b and references therein) causes the QSOs to be optically much less brighter than the observed QSOs. The distinctive evolutionary feature at the hard X-ray energies shows up in the SES model. This feature is essentially summarized as the the reversal in the direction of luminosity function evolution. Such an apparently puzzling feature is solely due to the accretion flow transition from a thin disk to an ADAF. In the SEM model, however, this effect is cancelled out due to a successive evolution of multiple generations, a result which doesn’t sensitively depend on the exact number of QSO generations as long as the SEM model has a large enough generations to be qualified as a multiple population model. It is interesting to point out that we can infer from the existence or non-existence of this particular feature whether the QSO population is composed of multiple generations or a single generation.
### 2.6 Multi-Band QSO Luminosity Evolution
We now look into how the multiple population model is distinguished from the single population model if the QSO luminosity evolution is not accompanied by the spectral evolution. In the case of no-spectral evolution, the luminosity, $`L`$, in each band is simply proportional to the bolometric luminosity, $`L_{bol}`$, i.e.
$$L=fL_{bol}=f(\eta \dot{M}c^2),$$
(7)
where $`\eta 0.1`$ is the radiative efficiency. We adopt the fraction $`f`$ of each band as 0.43, 0.47, 0.06, and 0.04 in the bands $`4400\AA ,1216\AA `$, soft X-ray, and hard X-ray, respectively. These fractions have been chosen solely based on the need to fit the observed LF evolution in each band and hence at this point, we don’t have any compelling physical motivation for these numbers. As anticipated, the evolution in all bands essentially have similar trends. The fraction of the bolometric luminosity radiated in the optical band of $`4400\AA `$, $`f_{4400\AA }`$, is about four times larger than the conventionally assumed value, $`f_B=0.1`$ (e.g. Haehnelt et al. 1998) in our case. We have not considered detailed corrections for dust absorption.
In the absence of the strong spectral evolution induced by the accretion flow transition to low efficiency ADAFs, the luminosity evolution is generally much slower than in the case of the spectral evolution. An accelerated evolution occurs in which luminosities drop dramatically below the critical mass accretion rate, $`\dot{m}_c`$. ADAFs appear with their characteristic hard spectra as shown in Figure 1 (b). In the single population models with no-spectral evolution, the obtained best-fit QSO evolutionary time scale is shorter than that of the spectral evolution case. Figure 4 (a2) and (b2) show the best-fit LFs in various energy bands in the single population no-spectral evolution model (hereafter NES model) and no-spectral evolution with multiple populations (NEM model), respectively. All parameters used in the case of the non-spectral evolution are given in Table 1. The resulting MF of BH remnants in each model is shown in Figure 5. Comparing the SEM model with the NEM model, we find that the optical LF in NEM model gives a better fit to the observed one (B band LF of QSO observed by Boyle et al. (1991)) than that of the SEM model even if there are no substantial differences in the evolving energy-dependent profiles of LFs in this NE- series models. Despite the physical motivation for the spectral evolution, we conclude that the NEM model gives the evolving shape of LFs which fits the observed optical ($`4400\AA `$) as well as soft X-ray LFs even if, in this NE- series models, both of LFs show the same evolution rates, differently from the observed one. Wisotzki (2000b) attempts to reconcile the discrepant evolution rates of optical and X-ray LFs by modifying an optical LF (i.e. a new determination of the optical K correction for QSOs). It still remains problematic that the BH number density corresponding to the observed soft X-ray LF (Miyaji et al. 2000) is too high to match the number density observed by Boyle et al. (1991). It appears that our results indicate that the optical and X-ray LFs evolve separately with different number densities. For instance, Haiman & Menou (2000) suggest the possibility that the ratio of optical to X-ray emission of quasars could evolve with redshift due to the cosmological structure formation and accompying evolution of the QSO environment.
The two models, one with the spectral evolution and the other without, have substantial differences in the hard X-ray luminosity evolution. Consequently, as previously predicted (Choi et al. 1999b), the observational LF data in the hard X-ray, which will be obtained by Chandra X-ray observatory could play a role in determining the importance of the spectral evolution. In conclusion, the NEM model is more plausible in accounting for the observed LF. This result partially supports the possibility that the QSO population is composed of the multiple short-lived (a few 10<sup>8</sup> yrs) generations of which the lifetime is short enough to have no significant spectral evolution over each generation’s evolution. This is in agreement with the observational result that the spectra in the X-ray bands show no conspicuous variation in the spectral index with redshift (e.g. Blair et al. 2000).
## 3 QSO LFs in the Disk Instability Model
The observational and theoretical similarities between Galactic X-ray binaries and AGN in emission mechanisms and spectral behavior lead us not only to assume the QSO spectrum-luminosity correlation but also to expect the QSO accretion flows’ time-dependent phenomena similar to those of the Galactic binary systems. Most notably, the occasional outbursts from black hole binaries in the Galaxy could be relevant for AGN/QSO emission from scaled-up accretion flows around supermassive black holes. The enormous differences between Galactic black hole systems and AGN/QSO should be reflected in time scales of similar phenomena such as the recurrent outbursts. If the variability due to the thermal-viscous ionization instabilities in accretion disks in cataclysmic variables or X-ray transients does operate in accretion disk in AGN on longer time scales (Siemiginowska, Czerny, & Kostyunin 1996 and reference therein, Mineshige & Shields 1990), the observed QSO/AGN LFs could be affected by the recurrent outbursts. Siemiginowska & Elvis (1997) assumed that all QSOs are subject to this variability and calculated the LF of a population of identical sources with a single mass and an accretion rate while the possible ADAF-like emission at the low luminosity end of the light curve was ignored.
Here, we consider the disk instability model for QSOs in which QSOs’ disk emission is periodically modulated with recurrent outbursts separated by quiescence. We apply our simple assumptions on the luminosity evolution of QSOs to this disk instability model and derive the LFs and MFs for four sets of evolution models of QSOs; (i) Disk instability induced spectral evolution model for a single QSO population (DSES model), (ii) Disk instability induced non-spectral evolution model of a single QSO population (DNES model), (iii) Disk instability induced spectral evolution model of multiple QSO populations (DSEM model), and (iv) Disk instability induced non-spectral evolution model of multiple QSO populations (DNEM model). The resulting LFs and MFs for these disk instability models are included in Figure 4 and 5. In order to obtain the soft X-ray LF in good agreement with the observed LFs of AGN in the soft X-ray band (Miyaji et al. 2000), the parameters are optimized and the comoving number density is normalized. The best-fit parameters determined for various models are listed in Table 1.
We assume that a single QSO undergoes repeated bright and faint phases. These bright and faint phases correspond to outbursts and quiescent states respectively. When the spectral evolution model is adopted, the behavior of low-luminosity part is essentially determined by the ADAF emission. A number of new parameters are introduced in the disk instability models. The major model parameters include the amplitude of accretion rate variation, duty cycle between the active phase and the quiescent phase, and the number of bursts. We assume that the amplitude of accretion rate variation is constant throughout the evolution while the maximum and minimum accretion rates decrease exponentially with the e-folding time scale $`t_{evol}`$ as defined in Section 2. In the model calculations, the maximum accretion rate does not exceed the Eddington accretion rate, $`\dot{m}=1`$, and the variability is taken to be periodic in cosmic time. Due to the redshift dependence of time scale, the variability in the redshift unit depends on the redshift epoch. The duty cycle, the ratio of the elapsed time of QSO in the active phase and the quiescent phase, $`t_{act}/t_q`$, essentially characterizes the outburst time scale in the disk instability model. There are two important simplifying assumptions we adopt here in order to make the calculations possible. First, an identical set of period and amplitude of variabilities is applied to QSOs regardless of the overall QSO evolution time scale. Second, because the number of parameters required in a relatively simple DSES model is already excessive, the inter-dependence among these parameters is ignored for practical calculational purposes.
### 3.1 Single Population Model with Disk Instability
Figure 6 shows the parameter dependence of the soft X-ray LFs for DSES model. The thermal instability in the accretion disk occurs continuously over the cosmological time scale. See the QSO light curve for a BH mass of $`10^8M_{}`$ in panel (a). The amplitude of accretion rate variation, $`\mathrm{\Delta }`$ log $`\dot{m}`$ is set at 1 and the accretion rate changes without much gain in mass. Because the BH gains most of its mass during an early active phase, the masses of BHs in this model do not grow substantially. Therefore, it is required that the sample QSOs experiencing thermal-viscous ionization instabilities in accretion disks have more massive seed BHs in order to fit the observed LF and MF. The maximum of accretion rate starts from 1 in Eddington unit at the time birth and subsequently it falls exponentially below the critical rate, log $`\dot{m}_c=2`$ causing the transition of accretion flow from thin disk (HS) to ADAF (LS/OS). A QSO evolving with the larger $`\mathrm{\Delta }`$ log $`\dot{m}`$ experiences the ADAF phase earlier and longer than one with the smaller amplitude. The ratio of $`t_{act}/t_q=0.25`$ and the number of outbursts of $`10^3`$ mean that QSO spends about $`20\%`$ in the active phase and about $`80\%`$ in the quiescent phase and the time scale for a burst in the active phase, $`t_{act}`$, is of the order of $`10^6`$ yr, respectively.
From the light curve in panel (a) in Figure 6, we see the dominant redshift dependence of time scale which shows that in the redshift plot the recurrence time scale shortens toward low redshifts. Panel (b1) gives the LFs which fit the observed ones in this DSES model. In Figure 4 (a3) and (a4), the LFs in various energy bands for DSES and DNES models are displayed. In these disk instability models, the number density of bright QSOs does not remain constant and it decreases just as in the SEM and NEM models. This is simply because the dim QSOs in quiescent phase effectively disappear in number counting. Based on the comparisons between SES and DSES models and between NES and DNES models, we see that there are little differences between them in the overall shape of LFs, which in part implies that the adopted disk instability model (see the best-fit parameters in Table 1) may not be sophisticated enough to catch some potential subtle differences. Within the accuracy of the current models, it is plausible that the QSO accretion disks may indeed undergo the thermal instability and that the observed LFs may well be accounted for by the disk instability model. This possibility is not directly testable in individual QSOs because the variability time scales expected for AGN/QSO in our model are simply too long (e.g. $`10^510^6`$ yr) and is not directly observable.
From panel (b1) vs. (b4) and panel (b2) vs. (b3) in Figure 6, we see the effects of the amplitude of accretion rate variation on the LFs in DSES model. The amplitude at the level of $`\mathrm{\Delta }`$ log $`\dot{m}=4`$ allows the QSO accretion disk to undergo recurrent spectral changes between the HS and the LS/OS. Starting from its birth, a QSO with a time-varying accretion disk spends most ($`80\%`$) of its evolution in the quiescent phase. The quiescent phase is comprised of LS and OS as the mass accretion rate fluctuates. A large number of QSOs remain in the very dim luminosity state while dramatically increasing the undetected low luminosity galactic nuclei. This increase in the number density which is normalized to match the comoving number density of Miyaji et al. (2000) depends on the amplitude of the mass accretion rate fluctuations. For instance, the derived comoving number density of QSOs (both observed and dormant) in the case of $`\mathrm{\Delta }`$ log $`\dot{m}=4`$ is roughly 10 times higher than that in the case of $`\mathrm{\Delta }`$ log $`\dot{m}=1`$. From panel (b1) vs. (b2) and panel (b3) vs. (b4), we show the dependence of the LF on the duty cycle in the DSES model, and from from panel (b1) vs. (b5), the dependence on the burst frequency. The number of bursts $`10^4`$ means that $`t_{act}`$ for a burst is of the order of $`10^5`$ yr. As the duty cycle shortens, the time elapsed in the quiescent phase of QSO increases and therefore, the more QSOs tend to stay in the low luminosity part of LF at any given epoch. The resulting LFs in the panels, however, show that the number of bursts and the duration of the duty cycle do not affect the evolving shape of the LF significantly in contrast to the effects of the amplitude of accretion rate variation.
### 3.2 The Multiple Population Model with Disk Instability
In Figure 4, panel (b3) and (b4) show the best-fit LFs in various energy bands for DSEM and DNEM models. MFs of remnant BHs corresponding to these LFs appear in Figure 5. Among the multiple population models, there is little difference in the evolving shape of LFs. It is remarkable that even the evolving shape of LF of the DSEM model with $`\mathrm{\Delta }`$ log $`\dot{m}=4`$ is not conspicuously different from those of the other multiple population models. This is interesting because the derived comoving number density for this very large amplitude case is substantially bigger (by a factor $`10`$) than those of the other models. This example demonstrates that the observed LF alone cannot rule out a large number of dormant and hence undetected QSO/AGN at any given epoch. However, this large comoving number density value is still short of the required level in the MF. We conclude that the disk instability model cannot simultaneously account for the LF and MF.
## 4 DISCUSSIONS
We have explored diverse types of possible phenomenological scenarios for the cosmological evolution of QSOs. Based on the basic paradigm in which QSO/AGN are powered by accreting supermassive black holes (Frank et al. 1992), we have considered various spectral states correlated with the luminosity level (Choi et al. 1999b) and constructed the analytical LFs in the various energy bands. Using the MF of remnant BHs in nearby galactic nuclei as an additional constraint, we have examined whether there exists a model simultaneously satisfying both LF and MF. The best and hence most promising scenario along with the best fit parameters points to the following formation and accretion history for supermassive BHs and QSOs.
1. In the multiple population models, due to the rapid decrease of $`\dot{M}`$ caused by the short evolutionary time scale and to the redshift dependence of the evolutionary time scale, the number density of bright QSOs has to evolve strongly over the cosmological time. According to the derived soft X-ray LFs in a good agreement with observed LFs (Miyaji et al. 2000), the QSOs with less massive seed BHs must have been born more abundantly and/or refueled more efficiently at lower redshifts.
2. While the space number density obtained from the soft X-ray LF of AGN by (Miyaji et al. 2000) is about only 3 times higher than that observed by Boyle et al. (1991), the resulting comoving space densities in the BH remnants inferred from the derived best-fit soft X-ray LFs in all QSO evolution models, which we have tested, are smaller than those estimated for the putative BHs in nearby galaxies by a factor of about 10 to 100. No models we have explored can derive the LF and MF satisfying both of the two observational constraints well enough. This could imply that it is highly unlikely that all galaxies contain MDOs which are actually supermassive BH remnants of QSOs. Alternatively, there could be a population of highly obscured AGN which harbor growing BHs (Fabian 1999, Hasinger 2000).
3. There is little difference in evolving shapes of soft X-ray LFs among all the models we considered. This is inevitable given the fact that our models have at least been required to fit the observed LFs of QSO/AGN in the soft X-ray bands (Miyaji et al. 2000). We, however, find substantial differences in optical ($`4400\AA `$) LFs and hard X-ray LFs, which is a natural consequence of the assumed spectrum-luminosity correlation (Choi et al. 1999a, 1999b, 2000). In the case of the spectral evolution, the optical LF evolves more rapidly than the X-ray LF especially at lower redshifts or when QSOs become faint. There is a great difference among the evolving features of the hard X-ray LFs in various models. This is one of the strongest outcomes in the spectral evolution models. Expected observational data by hard X-ray missions such as the Chandra X-ray observatory could play a role in providing a much needed, additional observational constraint. Such a constraint will shed light on which QSO models are more plausible and provide a statistical test on whether QSOs have experienced the spectral transition over the cosmological time scale. Such a transition is obviously untestable in individual galaxies due to the prohibitively long evolution and transition time scales.
4. The multiple population model with no spectral evolution gives the evolving shape of LFs which fit the observed B band LF of QSO derived by Boyle et al. (1991) in the optical ($`4400\AA `$) band as well as the soft X-ray LF. It is by far the best model among the various evolution models we have tested. This result supports the possibility that the QSO population is composed of many short-lived (a few 10<sup>8</sup> yr) generations while in each short generation QSOs do not experience any significant spectral evolution.
This work was supported by a KRF grant No. 1999-001-D00365. JY was also supported by the MOST through the National R & D program (99-N6-01-01-A-06) for women’s universities. IY wishes to thank Anna I. Yi and Fojyik Yi for numerous discussions and helpful suggestions and Ethan Vishniac for hospitality while this work was in progress.
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