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# Optimization of entanglement witnesses
## I Introduction
Quantum entanglement , which is an essence of many fascinating quantum mechanical effects , is a very fragile phenomenon. It is usually very hard to create, maintain, and manipulate entangled states under laboratory conditions. In fact, any system is usually subjected to the effects of external noise and interactions with the environment. These effects turn pure state entanglement into mixed state, or noisy entanglement. The separability problem, that is, the characterization of mixed entangled states, is highly nontrivial and has not been accomplished so far. Even the apparently innocent question: Is a given state entangled and does it contain quantum correlations, or is it separable, and does not contain any quantum correlations? will, in general, be very hard (if not impossible!) to answer.
Mathematically, mixed state entanglement can be described as follows. A density operator $`\rho 0`$ acting on a finite Hilbert space $`H=H_AH_B`$ describing the state of two quantum systems A and B is called separable (or not entangled) if it can be written as a convex combination of product vectors; that is, in the form
$$\rho =\underset{k}{}p_k|e_k,f_ke_k,f_k|,$$
(1)
where $`p_k0`$, and $`|e_k,f_k|e_k_A|f_k_B`$ are product vectors. Conversely, $`\rho `$ is nonseparable (or entangled) if it cannot be written in this form. Physically, a state described by a separable (nonseparable) density operator $`\rho `$ can always (never) be prepared locally. Most of the applications in quantum information are based on the nonlocal properties of quantum mechanics, and therefore on nonseparable states. Thus, a criterion to determine whether a given density operator is nonseparable, i.e. useful for quantum information purposes, or not is of crucial importance. On the other hand, PPTES are objects of special interest since they represent so–called bound entangled states, and therefore provide an evidence of irreversibility in quantum information processing .
For low dimensional systems there exist operationally simple necessary and sufficient conditions for separability. In fact, in $`H=\mathrm{I}\mathrm{C}^2\mathrm{I}\mathrm{C}^2`$ and $`H=\mathrm{I}\mathrm{C}^2\mathrm{I}\mathrm{C}^3`$ the Peres–Horodecki criterion establishes that $`\rho `$ is separable iff its partial transpose is positive. Partial transpose means a transpose with respect to one of the subsystems . For higher dimensional systems all operators with non–positive partial transposition are entangled. However, there exist positive partial transpose entangled states (PPTES) . Thus, the separability problem reduces to finding whether density operators with positive partial transpose are separable or not .
In the recent years there has been a growing effort in searching for necessary and sufficient separability criteria and checks which would be operationally simple . Several necessary or sufficient conditions for separability are known. A particularly interesting necessary condition is given by the so–called range criterion . According to this criterion, if the state $`\rho `$ acting on a finite dimensional Hilbert space is separable then there must exist a set of product vectors $`\{|e_k,f_k\}`$ that spans the range $`R(\rho )`$ such that the set of partial complex conjugated product states $`\{|e_k,f_k^{}\}`$ spans the range of the partial transpose of $`\rho `$ with respect to the second system, i.e., $`\rho ^{T_B}`$. Among the PPTES that violate this criterion there are particular states with the property that if one subtracts a projector onto a product vector from them, the resulting operator is no longer a PPTES . In this sense, these states lie in the edge between PPTES and entangled states with non–positive partial transposition, and therefore we will call them “edge” PPTES. The analysis of the range of density operators initiated in Ref. has turned out to be very fruitful. In particular, it has led to an algorithm for the optimal decomposition of mixed states into a separable and an inseparable part , and to a systematic method of constructing examples of PPTES using unextendible product bases . For low rank operators it has allowed to show that one can reduce the separability problem to the one of determining the roots of certain complex polynomial equations .
From a different point of view, a very general approach to analyze the separability problem is based on the so–called entanglement witnesses (EW) and positive maps (PM) . Entanglement witnesses are operators that detect the presence of entanglement. Starting from these operators one can define PM’s that also detect entanglement. An example of a PM is precisely partial transposition . The importance of EW stems from the fact that a given operator is separable iff there exists an EW that detects it . Thus, if one was able to construct all possible EW (or PM) one would have solve the problem of separability. Unfortunately, it is not known how to construct EW that detect PPTES in general. The only result in this direction so far has been given in Ref. , although some preliminary results exist in the mathematical literature . Starting from a PPTES fulfilling certain properties (related to the existence of unextendible basis of product vectors ), it has been shown how to construct an EW (and the corresponding PM) that detects it. Perhaps, one of the most interesting goals regarding the separability problem is to develop a constructive and operational approach using EW and PM that allows us to detect mixed entanglement.
In this paper we realize this goal partially: we introduce a powerful technique to construct EW and PM that, among other things, allows us to study the separability of certain density operators . In particular, we show how to construct optimal EW; that is, operators that detect the presence of entanglement in an optimal way. We specifically concentrate on non–decomposable EW, which are those that detect the presence of PPTES. Furthermore, we present a way of constructing optimal EW for edge PPTES. Our method generalizes the one introduced by Terhal to the case in which there are no unextendible basis of product vectors. When combined with our previous results regarding subtracting product vectors from PPTES, the construction of non–decomposable optimal EW starting from “edge” PPTES gives rise to a novel sufficient criterion for non–separability of general density operators with positive partial transposition. We illustrate our method by constructing optimal EW that detect some known examples of PPTES in $`H=\mathrm{I}\mathrm{C}^2\mathrm{I}\mathrm{C}^4`$. The corresponding PM constitute the first examples of PM with minimal “qubit” domain, or – equivalently – minimal hermitian conjugate codomain.
This paper is organized as follows. In Section II we review the definition of EW and fix some notation. In Section III we study general EW. We define optimal witnesses and find a criterion to decide whether an EW is optimal or not. In Section IV we restrict the results of Section III to non–decomposable EW. In particular, we show how to optimize them by subtracting decomposable operators. In Section V we give an explicit method to optimize both, general and non–decomposable EW. We also show how to construct non–decomposable EW, and that this leads to a sufficient criterion of non–separability. The construction and optimization is based on the use of “edge” PPTES. In Section VI we extend our results to positive maps. In Section VII we illustrate our methods and results starting from the examples of PPTES given in Ref. . The paper also contains two appendices. In Appendix A we describe in detail a method to check whether an EW is optimal or not. In Appendix B we discuss separately some important properties of the edge PPTES, and show that they provide a canonical decomposition of mixed states with positive partial transpose.
## II Definitions and notation
We say that an operator $`W=W^{}`$ acting on $`H=H_AH_B`$ is an EW if :
$`e,f|W|e,f0`$ for all product vectors $`|e,f`$;
has at least one negative eigenvalue (i.e. is not positive);
$`\mathrm{tr}(W)=1`$.
The first property (I) implies that $`\rho _W\mathrm{tr}(W\rho )0`$ for all $`\rho `$ separable. Thus, if we have $`\rho _W<0`$ for some $`\rho 0`$, then $`\rho `$ is nonseparable. In that case we say that $`W`$ detects $`\rho `$. The second one (II) implies that every EW detects something, since in particular it detects the projector on the subspace corresponding to the negative eigenvalues of $`W`$. The third property (III) is just normalization condition that we need in order to compare the action of different EW .
In this paper we will denote by $`K(\rho )`$ and $`R(\rho )`$ the kernel and range of $`\rho `$, respectively. The partial transposition of an operator $`X`$ will be denoted by $`X^T`$ . On the other hand, we will encounter several kinds of operators (EW, positive operators, decomposable operators, etc) and vectors. In order to help to identify the kind of operators and vectors we use, and not to overwhelm the reader by specifying at each point their properties, we will use the following notation:
* $`W`$ will denote an EW.
* $`P,Q`$ will denote positive operators. Unless specified they will have unit trace \[$`\mathrm{tr}(P)=\mathrm{tr}(Q)=1`$\].
* $`D`$ will denote a decomposable operator. That is, $`D=aP+bQ^T`$, where $`a,b0`$. Unless stated, all decomposable operators that we use will have unit trace (i.e., $`b=1a`$).
* $`\rho `$ will denote a positive operator (not necessarily of trace 1).
* $`|e,f`$ will denote product vectors with $`|eH_A`$ and $`|fH_B`$. Unless especified, they will be normalized.
## III General entanglement witnesses
In this Section we first give some definitions directly related to EW. Then we introduce the concept of optimal EW. We derive a criterion to determine when an EW is optimal. This criterion will serve us to find an optimization procedure for these operators.
### A Definitions
Given an EW, $`W`$, we define:
* $`D_W=\{\rho 0,\text{ such that }\rho _W<0\}`$; that is, the set of operators detected by $`W`$.
* Finer: Given two EW, $`W_1`$ and $`W_2`$, we say that $`W_2`$ is finer than $`W_1`$, if $`D_{W_1}D_{W_2}`$; that is, if all the operators detected by $`W_1`$ are also detected by $`W_2`$.
* Optimal entanglement witness (OEW): We say that $`W`$ is an OEW if there exist no other EW which is finer.
* $`P_W=\{|e,fH,\text{ such that }e,f|W|e,f=0\}`$; that is, the set of product vectors on which $`W`$ vanishes. As we will show, these vectors are closely related to the optimality property.
Note the important role that the vectors in $`P_W`$ play regarding entanglement (for a method to determine $`P_W`$ in practice, see Appendix A). If we have an EW, $`W`$, which detects a given operator $`\rho `$, then the operator $`\rho ^{}=\rho +\rho _w`$ where
$$\rho _w=\underset{k}{}p_k|e_k,f_ke_k,f_k|$$
(2)
with $`p_k0`$, and $`|e_k,f_kP_W`$ is also detected by $`W`$. In fact, this means that any operator of the form (2) is in the border between separable states and non–separable states, in the sense that if we add an arbitrarily small amount of $`\rho `$ to it we obtain a non–separable state. Thus, the structure of the sets $`P_W`$ characterizes the border between separable and non–separable states. In fact, from the results of this Section it will become clear that we can restrict ourselves to the structure of the sets of $`P_W`$ corresponding to OEW’s.
### B Optimal entanglement witnesses
According to Ref. $`\rho `$ is nonseparable iff there exists an EW which detects it. Obviously, we can restrict ourselves to the study of OEW. For that, we need criteria to determine when an EW is optimal. In this subsection we will derive a necessary and sufficient condition for this to happen (Theorem 1 below). In order to do that, we first have to introduce some results that tell us under which conditions an EW is finer than another one.
Lemma 1: Let $`W_2`$ be finer than $`W_1`$ and
$$\lambda \underset{\rho _1D_{W_1}}{inf}\left|\frac{\rho _1_{W_2}}{\rho _1_{W_1}}\right|.$$
(3)
Then we have:
If $`\rho _{W_1}=0`$ then $`\rho _{W_2}0`$.
If $`\rho _{W_1}<0`$, then $`\rho _{W_2}\rho _{W_1}`$.
If $`\rho _{W_1}>0`$ then $`\lambda \rho _{W_1}\rho _{W_2}`$.
$`\lambda 1`$. In particular, $`\lambda =1`$ iff $`W_1=W_2`$.
Proof: Since $`W_2`$ is finer than $`W_1`$ we will use the fact that for all $`\rho 0`$ such that $`\rho _{W_1}<0`$ then $`\rho _{W_2}<0`$.
(i) Let us assume that $`\rho _{W_2}>0`$. Then we take any $`\rho _1𝒟_{W_1}`$ so that for all $`x0`$, $`0\stackrel{~}{\rho }(x)\rho _1+x\rho 𝒟_{W_1}`$. But for sufficiently large $`x`$ we have that $`\stackrel{~}{\rho }(x)_{W_2}`$ is positive, which cannot be since then $`\rho (x)𝒟_{W_2}`$.
(ii) We define $`\stackrel{~}{\rho }=\rho +|\rho _{W_1}|1𝐥0`$. We have that $`\stackrel{~}{\rho }_{W_1}=0`$. Using (i) we have that $`0\rho _{W_2}+|\rho _{W_1}|`$.
(iii) We take $`\rho _1D_{W_1}`$ and define $`\stackrel{~}{\rho }=\rho _{W_1}\rho _1+|\rho _1_{W_1}|\rho 0`$, so that $`\stackrel{~}{\rho }_{W_1}=0`$. Using (i) we have $`|\rho _1_{W_1}|\rho _{W_2}|\rho _1_{W_2}|\rho _{W_1}`$. Dividing both sides by $`|\rho _1_{W_1}|>0`$ and $`\rho _{W_1}>0`$ we obtain
$$\frac{\rho _{W_2}}{\rho _{W_1}}\left|\frac{\rho _1_{W_2}}{\rho _1_{W_1}}\right|.$$
(4)
Taking the infimum with respect to $`\rho _1D_{W_1}`$ in the rhs of this equation we obtain the desired result.
(iv) From (ii) immediately follows that $`\lambda 1`$. On the other hand, we just have to prove that if $`\lambda =1`$ then $`W_1=W_2`$ (the only if part is trivial). If $`\lambda =1`$, using (i) and (iii) we have that $`\rho _v_{W_1}\rho _v_{W_2}`$ for all $`\rho _v=|e,fe,f|`$ projector on a product vector. Since $`\mathrm{tr}(W_1)=\mathrm{tr}(W_2)`$ we must have $`\mathrm{tr}[(W_1W_2)\rho _v]=0`$ for all $`\rho _v`$, since we can always find a product basis in which we can take the trace. But now, for any given $`\rho 0`$ we can define $`\stackrel{~}{\rho }(x)=\rho +x1𝐥`$ such that for large enough $`x`$, $`\stackrel{~}{\rho }(x)`$ is separable . In that case we have $`\stackrel{~}{\rho }(x)_{W_1}=\stackrel{~}{\rho }(x)_{W_2}`$ which implies that $`\rho _{W_1}=\rho _{W_2}`$, i.e. $`W_1=W_2`$. $`\mathrm{}`$
Corollary 1: $`D_{W_1}=D_{W_2}`$ iff $`W_1=W_2`$.
Proof: We just have to prove the only if part. For that, we define $`\lambda `$ as in (3). On the other hand, defining
$$\stackrel{~}{\lambda }\underset{\rho _2D_{W_2}}{inf}\left|\frac{\rho _2_{W_1}}{\rho _2_{W_2}}\right|$$
(5)
we have that $`\stackrel{~}{\lambda }1`$ since $`W_1`$ is finer than $`W_2`$ (Lemma 1(iv)). Equivalently,
$$1\underset{\rho _1D_{W_1}}{sup}\left|\frac{\rho _1_{W_2}}{\rho _1_{W_1}}\right|\lambda 1,$$
(6)
where for the last inequality we have used that $`W_2`$ is finer than $`W_1`$. Now, since $`\lambda =1`$ we have that $`W_1=W_2`$ according to Lemma 1(iv). $`\mathrm{}`$
Next, we introduce one of the basic results of this paper. It basically tell us that EW is finer than another one if they differ by a positive operator. That is, if we have an EW and we want to find another one which is finer, we have to subtract a positive operator.
Lemma 2: $`W_2`$ is finer than $`W_1`$ iff there exists a $`P`$ and $`1>ϵ0`$ such that $`W_1=(1ϵ)W_2+ϵP`$.
Proof: (If) For all $`\rho D_{W_1}`$ we have that $`0>\rho _{W_1}=(1ϵ)\rho _{W_2}+ϵ\rho _P`$ which implies $`\rho _{W_2}<0`$ and therefore $`\rho D_{W_2}`$. (Only if) We define $`\lambda `$ as in (3). Using Lemma 1(iv) we have $`\lambda 1`$. First, if $`\lambda =1`$ then according to Lemma 1(iv) we have $`W_1=W_2`$ (i.e., $`ϵ=0`$). For $`\lambda >1`$, we define $`P=(\lambda 1)^1(\lambda W_1W_2)`$ and $`ϵ=11/\lambda >0`$. We have that $`W_1=(1ϵ)W_2+ϵP`$, so that it only remains to be shown that $`P0`$. But this follows from Lemma 1(i–iii) and the definition of $`\lambda `$, $`\lambda =inf_{\rho _1D_{W_1}}\left|\frac{\rho _1_{W_2}}{\rho _1_{W_1}}\right|`$ . $`\mathrm{}`$
The previous lemma provides us with a way of determining when an EW is finer than another one. With this result, we are now at the position of fully characterizing OEW.
Theorem 1: $`W`$ is optimal iff for all $`P`$ and $`ϵ>0`$, $`W^{}=(1+ϵ)WϵP`$ is not an EW \[does not fulfill (I)\].
Proof: (If) According to Lemma 2, there is no EW which is finer than $`W`$, and therefore $`W`$ is optimal. (Only if) If $`W^{}`$ is an EW, then according to Lemma 2 $`W`$ is not optimal. $`\mathrm{}`$
The previous theorem tells us that $`W`$ is optimal iff when we subtract any positive operator from it, the resulting operator is not positive on product vectors. This result is not very practical because of two reasons: (1) for a given $`P`$ it is typically very hard to check whether there exists some $`ϵ>0`$ such that $`WϵP`$ is positive on all product vectors; (2) it may be difficult to find a particular $`P`$ that can be subtracted from $`W`$ among all possible positive operators. In Appendix A we show how to circumvent these two drawbacks in practice: we give a simple criterion to determine when a given $`P`$ can be subtracted from $`W`$. This allows us to determine which are the positive operators which can be subtracted from a given EW.
In the rest of this subsection we will present some simple results related to these two questions. First, it is clear that not every positive operator $`P`$ can be subtracted from an EW, $`W`$. In particular, the following lemma tells us that it must vanish on $`P_W`$.
Lemma 3: If $`PP_W0`$ then $`P`$ cannot be subtracted from $`W`$.
Proof: There exists some $`|e_0,f_0P_W`$ such that $`e_0,f_0|P|e_0,f_0>0`$. Substituting this product vector in the condition I for any $`WϵP`$ we see that the inequality is not fulfilled for any $`ϵ>0`$, i.e. $`P`$ cannot be subtracted. $`\mathrm{}`$
Corollary 2: If $`P_W`$ spans $`H`$ then $`W`$ is optimal.
Note that, as announced at the beginning of this Section, the set $`P_W`$ plays an important role in determining the properties of the separable states which lie on the border with the entangled states. We see here, that this set also plays an important role in determining whether an EW is optimal or not.
On the other hand, in order to check whether a given operator $`P`$ can be subtracted or not from $`W`$, one has to check whether there exists some $`ϵ>0`$ such that $`e,f|WϵP|e,f>0`$ for all $`|e,f`$. The following lemma gives an alternative way to do this. In fact, it gives a necessary and sufficient criterion for an EW to be optimal. For a given $`|eH_A`$, we will denote by $`W_ee|W|e`$.
Lemma 4: $`W`$ is optimal iff for all $`|\mathrm{\Psi }`$ orthogonal to $`P_W`$
$$ϵ\underset{|eH_A}{inf}\left[\mathrm{\Psi }|eW_e^1e|\mathrm{\Psi }\right]^1=0.$$
(7)
Proof: (If) Let us assume that $`W`$ is not optimal; that is, there exists $`W^{}W`$, finer than $`W`$. Then, according to Lemma 2 we have that there exists $`ϵ_0>0`$ and $`P0`$ such that $`W^{}=(Wϵ_0P)/(1ϵ_0)`$. Imposing that $`W^{}`$ is positive on product vectors (i.e. $`W_e^{}0`$ for all $`|eH_A`$) we obtain $`0e|Wϵ_0P|eW_eϵ_0\lambda _\mathrm{\Psi }e|\mathrm{\Psi }\mathrm{\Psi }|e`$, where $`|\mathrm{\Psi }`$ is any eigenstate of $`P`$ with nonzero eigenvalue $`\lambda _\mathrm{\Psi }`$. According to Ref. , this last operator is positive iff both: (i) $`e|\mathrm{\Psi }`$ is in the range of $`e|W|e`$, which imposes that $`|\mathrm{\Psi }`$ is orthogonal to $`P_W`$; (ii) $`\lambda _\mathrm{\Psi }ϵ_0\left[\mathrm{\Psi }|eW_e^1e|\mathrm{\Psi }\right]^1`$, which imposes that $`ϵ\lambda _\mathrm{\Psi }ϵ_0>0`$ for that given $`|\mathrm{\Psi }`$. (Only if) Let us assume that there exists some $`|\mathrm{\Psi }`$ orthogonal to $`P_W`$ such that $`ϵ>0`$. Then, using the same arguments one can show that $`W^{}(Wϵ|\mathrm{\Psi }\mathrm{\Psi }|)/(1ϵ)W`$ is an EW. According to Lemma 2, $`W^{}`$ is finer than $`W`$, so that $`W`$ is not optimal. $`\mathrm{}`$.
### C Decomposable entanglement witnesses
There exists a class of EW which is very simple to characterize, namely the decomposable entanglement witnesses (d–EW) . Those are EW that can be written in the form
$$W=aP+(1a)Q^T,$$
(8)
where $`a[0,1]`$. As it is well known (see next section), these EW cannot detect PPTES. In any case, for the sake of completeness, we will give some simple properties of optimal d–EW.
Theorem 2: Given a d–EW, $`W`$, if it is optimal then it can be written as $`W=Q^T`$, where $`Q0`$ contains no product vector in its range.
Proof: Since $`W`$ is decomposable, it can be written as $`W=aP+(1a)Q^T`$. $`W^{}WaP`$ is also a witness, which according to Lemma 2 is finer than $`W`$, and therefore $`W`$ is not optimal. On the other hand, if $`|e,fR(Q)`$ then for some $`\lambda >0`$ we have that $`W(Q\lambda |e,fe,f|)^T`$ is finer than $`W`$, and therefore this last is not optimal. $`\mathrm{}`$
This previous result can be slightly generalized as follows:
Theorem 2’: Given a d–EW, $`W`$, if it is optimal then it can be written as $`W=Q^T`$, where $`Q0`$ and there is no operator $`PR(Q)`$ such that $`P^T0`$.
Proof: Is the same as in previous theorem. $`\mathrm{}`$
Corollary 3: Given a d–EW, $`W`$, if it is optimal then $`W^T`$ is not an EW \[does not fulfill (II)\].
Proof: Using Theorem 2 we have that $`W=Q^T`$ with $`Q0`$. Then $`W^T=Q0`$, which does not satisfy property (ii). $`\mathrm{}`$
## IV Non–decomposable entanglement witnesses
In the previous section we have been concerned with EW in general. As mentioned above, when studying separability we just have to consider those EW that can detect PPTES. In order to characterize them, one defines non–decomposable witnesses (nd–EW) as those EW which cannot be written in the form (8) . This Section is devoted to this kind of witnesses. The importance of nd–EW in order to detect PPTES is reflected in the following
Theorem 3: An EW is non–decomposable iff it detects PPTES.
Proof: (If) Let us assume that the EW is decomposable. Then it cannot detect PPT, since if $`\rho ,\rho ^T0`$ we have $`\mathrm{tr}[(aP+(1a)Q^T)\rho ]=a\mathrm{tr}(P\rho )+(1a)\mathrm{tr}(Q\rho ^T)0`$. (Only if) The set of decomposable witnesses is convex and closed, and $`W`$, as a set containing one point, is a closed convex set itself. Thus, from Hahn–Banach theorem it follows that there exists an operator $`\rho `$ such that: (i) $`\mathrm{tr}[\rho (aP+(1a)Q^T)]0`$ for all $`P,Q0`$, $`a[0,1]`$; (ii) $`\mathrm{tr}(\rho W)<0`$. From (i), taking $`a=1`$ we infer that $`\rho 0`$; on the other hand, taking $`a=0`$ we obtain that $`\mathrm{tr}[\rho ^TQ]0`$ for all $`Q0`$, and therefore $`\rho ^T0`$. Thus, $`W`$ detects $`\rho `$ which is a PPTES. $`\mathrm{}`$
Corollary 4: Given an operator $`D`$, it is decomposable iff $`\mathrm{tr}(D\rho )0`$ for all $`\rho ,\rho ^T0`$.
### A Definitions
In this Subsection we introduce some definitions which are parallel to those given in the previous Section. Given a nd–EW, $`W`$, we define:
* $`d_W=\{\rho 0,\text{ such that }\rho ^T0\text{ and }\rho _W<0\}`$; that is, the set of PPT operators detected by $`W`$.
* Non–decomposable-finer (nd–finer): Given two nd–EW, $`W_1`$ and $`W_2`$, we say that $`W_2`$ is nd–finer than $`W_1`$, if $`d_{W_1}d_{W_2}`$; that is, if all the operators detected by $`W_1`$ are also detected by $`W_2`$.
* Non–decomposable optimal entanglement witness (nd–OEW): We say that $`W`$ is an nd–OEW if there exist no other nd–EW which is nd–finer.
* $`p_W=\{|e,fH,\text{ such that }e,f|W|e,f=0\}`$; that is, the product vectors on which $`W`$ vanishes.
Note again the important role that the vectors in $`p_W`$ play regarding PPTES. If we have a nd–EW, $`W`$, which detects a given PPTES $`\rho `$, then the operator $`\rho ^{}=\rho +\rho _w`$ where $`\rho _w`$ has the form (2) with $`p_k0`$, and $`|e_k,f_kp_W`$ also describes a PPTES. Thus, any operator of the form (2) lies in the border between separable states and PPTES.
### B Optimal non–decomposable entanglement witness
The goal of this section is to find a necessary and sufficient condition for a nd–EW to be optimal. We start by proving a similar result to the one given in Lemma 1, but for nd–EW:
Lemma 1b: Let $`W_2`$ be nd–finer than $`W_1`$,
$$\lambda \underset{\rho _1d_{W_1}}{inf}\left|\frac{\mathrm{tr}(W_2\rho _1)}{\mathrm{tr}(W_1\rho _1)}\right|,$$
(9)
and now both, $`\rho ,\rho ^T0`$. Then we have have (i–iv) as in Lemma 1.
Proof: The proof is basically the same as in Lemma 1 and will be omitted here.
Corollary 1b: Given two nd–EW, $`W_{1,2}`$, then $`d_{W_1}=d_{W_2}`$ iff $`W_1=W_2`$.
Proof: The proof is basically the same as Corollary 1 and will be omitted here.
Lemma 2b: Given two nd–EW, $`W_{1,2}`$, $`W_2`$ is nd–finer than $`W_1`$ iff there exists a decomposable operator $`D`$ and $`1>ϵ0`$ such that $`W_1=(1ϵ)W_2+ϵD`$.
Proof: (If) Given any $`\rho ,\rho ^T0`$, we have that if $`\rho d_{W_1}`$ then $`0>\rho _{W_1}=(1ϵ)\rho _{W_2}+ϵ\rho _D\rho _{W_2}`$, where in the last inequality we have used that $`\rho _D0`$ since $`D`$ is decomposable (see Corollary 4). Therefore $`\rho d_{W_2}`$. (Only if) We define $`\lambda `$ as in (9), so that $`\lambda 1`$ according to Lemma 1b(iv). If $`\lambda =1`$ we have $`W_1=W_2`$. If $`\lambda >1`$ we define $`D=(\lambda 1)^1(\lambda W_1W_2)`$ and $`ϵ=11/\lambda `$. We have that $`W_1=(1ϵ)W_2+ϵD`$, so that it only remains to be shown that $`D`$ is decomposable. But from Lemma 1b(i–iii) and the definition of $`\lambda `$ it follows that $`\rho _D0`$ for all $`\rho ,\rho ^T0`$. Using Corollary 4 we then have that $`D`$ is decomposable. $`\mathrm{}`$.
Now we are able to fully characterize nd–OEW.
Theorem 1b: Given an nd–EW, $`W`$, it is nd–optimal iff for all decomposable operators $`D`$ and $`ϵ>0`$, $`W^{}=(1+ϵ)WϵD`$ is not an EW \[does not fulfill (I)\].
Proof: Is the same as for Theorem 1. $`\mathrm{}`$
Theorems 1 and 1b allow us to relate OEW and nd–OEW. In this way we can directly translate the results for general OEW to nd–OEW. We have
Theorem 4: Given a nd–EW, $`W`$, $`W`$ is a nd–OEW iff both $`W`$ and $`W^T`$ are OEW.
Proof: (If) Let us assume that $`W`$ is not a nd–OEW. Then, according to Theorem 1b there exists $`ϵ>0`$ and a decomposable operator $`D`$ such that $`W^{}=(1+ϵ)WϵD`$ is a nd–EW. We can write $`D=aP+(1a)Q^T`$, with $`a[0,1]`$. If $`a0`$, then $`W_1=(1+aϵ)WaϵP`$ fulfills $`e,f|W_1|e,fe,f|W^{}|e,f0`$, and therefore, according to Lemma 2, $`W`$ is not optimal. If $`a1`$ then $`W_2=[1+(1a)ϵ]W^T(1a)ϵQ`$ fulfills $`e,f|W_2|e,fe,f|(W^{})^T|e,f0`$, i.e. is an EW and therefore $`W^T`$ is not optimal. (Only if) According to Theorem 1b, if $`W`$ is nd–optimal then for all $`D=aP+(1a)Q^T`$, with $`a[0,1]`$, and all $`ϵ>0`$ we have that $`W^{}=(1ϵ)WϵD`$ does not satisfy (I). Taking $`a=1`$ we have for all $`P`$ and $`ϵ>0`$, $`W_1=(1ϵ)WϵP`$ does not fulfill (I), and therefore (Theorem 1) $`W`$ is optimal; analogously, taking $`a=0`$ we have that $`W^T`$ is optimal also. $`\mathrm{}`$
Corollary 5: $`W`$ is a nd–OEW iff $`W^T`$ is an nd–OEW.
## V Optimization
In this Section we give a procedure to optimize EW which is based on the results of the previous Sections.
### A Optimization of general entanglement witnesses
Our method is based in the following lemma. It tells us how much we can subtract from an EW. Here we will denote by $`W_e=e|W|e`$ and $`P_e=e|P|e`$ where $`|eH_A`$, by $`[\mathrm{}]_{\mathrm{min}}`$ the minimum eigenvalue, and by $`[\mathrm{}]_{\mathrm{max}}`$ the maximum eigenvalue. On the other hand, $`X^{1/2}`$ will denote the square root of the pseudoinverse of $`X`$ .
Lemma 5: If there exists some $`P`$ such that $`PP_W=0`$ and
$`\lambda _0`$ $``$ $`\underset{|eH_A}{inf}\left[P_e^{1/2}W_eP_e^{1/2}\right]_{\mathrm{min}}`$ (10)
$`=`$ $`\left(\underset{|eH_A}{sup}\left[W_e^{1/2}P_eW_e^{1/2}\right]_{\mathrm{max}}\right)^1>0.`$ (11)
then
$`W^{}(\lambda )(W\lambda P)/(1\lambda )`$ (12)
with $`\lambda >0`$ is an EW iff $`\lambda \lambda _0`$.
Proof: Let us find out for which values of $`\lambda 0`$, $`W^{}(\lambda )`$ is an EW. We have to impose condition (I), which can be written as $`e|W^{}(\lambda )|e0`$, i.e.
$$W_e\lambda P_e0.$$
(13)
Multiplying by $`P_e^{1/2}`$ on the right and left of this equation we obtain $`P_e^{1/2}W_eP_e^{1/2}\lambda `$, which immediately gives that $`\lambda \lambda _0`$ given in the first part of Eq. (10). On the other hand, multiplying by $`W_e^{1/2}`$ on the right and left of Eq. (13) we obtain $`W_e^{1/2}P_eW_e^{1/2}1/\lambda `$, which immediately gives that $`\lambda \lambda _0`$ given in the second equality of Eq. (10). $`\mathrm{}`$
Lemma 5 provides us with a direct method to optimize EW by subtracting positive operators for which the elements of $`P_W`$ are contained in their kernels. The method thus consists of: (1) determining $`P_W`$; (2) choosing an operator $`P`$ so that $`PP_W=0`$ and determining $`\lambda `$ using (10); (3) if $`\lambda 0`$ then we subtract the operator $`P`$ according to Lemma 5. Continuing in the same vein we will reach an OEW. In Appendix A we show how to accomplish steps (1) and (2) in practice.
### B Optimization of non–decomposable entanglement witnesses
For nd–EW we have the following generalization of Lemma 5:
Lemma 5b: Given a nd–EW, $`W`$, if there exists some decomposable operator $`D`$ such that $`Dp_W=0`$ and
$`\lambda _0\underset{|eH_A}{inf}\left[D_e^{1/2}W_eD_e^{1/2}\right]_{\mathrm{min}}=`$ (14)
$`\left(\underset{|eH_A}{sup}\left[W_e^{1/2}D_eW_e^{1/2}\right]_{\mathrm{max}}\right)^1>0.`$ (15)
then
$`W^{}(\lambda )(W\lambda D)/(1\lambda )`$ (16)
with $`\lambda >0`$ is a nd–EW iff $`\lambda \lambda _0`$.
Proof: Is the same as for Lemma 5.
With the help of Lemma 5b we can optimize nd–EW by subtracting decomposable operators as follows: (1) determining $`p_W`$ and $`p_{W^T}`$; (2) choosing $`P,Q`$ so that $`Pp_W=0`$ and $`Qp_{W^T}=0`$, building $`D=aP+(1a)Q^T`$ with $`a[0,1]`$, and determining $`\lambda _0`$ using (14); (3) if $`\lambda _00`$ then we subtract the operator $`D`$ according to Lemma 5b.
### C Detectors of “edge” PPTES
In the previous subsections we have have given two optimization procedures. In both of them, starting from a general EW one can obtain one which is optimal (or nd–optimal). It may well happen that the EW found in this way is non–decomposable even though the original one was decomposable. To check that one simply has to use Corollary $`3`$; that is, check whether $`W^T`$ is an EW or not. In case it is, then the OEW $`W`$ is non–decomposable. However, nothing guarantees that the final EW is non–decomposable if the original one is not. In this subsection we describe a general method to construct nd–EW using the optimization procedures introduced earlier. This method generalizes the one presented in Ref. .
We are going to use the results presented in Ref. . There, we have already used and discussed the “edge” PPTES, without naming them, however. Let us now introduce the following definition:
Definition: \[see Ref. \] A PPTES $`\delta `$ is an “edge” PPTES if for all product vector $`|e,f`$ and $`ϵ>0`$, $`\delta ϵ|e,fe,f|`$ is not a PPTES.
This definition implies that that the “edge” states lie on the boundary between the PPTES and entangled states with non–positive partial transpose. In this subsection we will show how, out of an “edge” PPTES, we can construct a nd–OEW that detects it. As we mentioned in the introduction, “edge” PPTES are of special importance. In particular, they allow to provide a canonical form to write PPTES in arbitrary Hilbert spaces. For these reasons, some of the properties of the “edge” PPTES are discussed in Appendix B.
In order to check whether a PPTES $`\delta `$ is an “edge” PPTES we can use the range criterion (see also ). That is, $`\delta `$ is an “edge” PPTES iff for all $`|e,fR(\delta )`$, $`|e,f^{}R(\delta ^{T_B})`$.
Let $`\delta `$ be an “edge” PPTES, and let us denote by $`P_1`$ the projector onto $`K(\delta )`$ and by $`Q_1`$ the projector onto $`K(\delta ^T)`$. We define
$$W_\delta =a(P_1+Q_1^T),$$
(17)
where $`a=1/\mathrm{tr}(P_1+Q_1)`$. Let us also define
$$ϵ_1\underset{|e,f}{inf}e,f|W_\delta |e,f.$$
(18)
Then we have
Lemma 6: Given an “edge” PPTES $`\delta `$, then $`W_1W_\delta ϵ_11𝐥`$ is a nd–EW, where $`ϵ_1`$ and $`W_\delta `$ are defined in (18,17), respectively.
Proof: We have that $`e,f|W_\delta |e,f=a(e,f|P_1|e,f+e,f^{}|Q_1|e,f^{})0`$. This quantity is zero iff $`e,f|P_1|e,f=e,f^{}|Q_1|e,f^{}=0`$. But this is not possible since $`\delta `$ is an “edge” PPTES. Thus, $`e,f|W_\delta |e,f>0`$ for all $`|e,f`$. Defining $`ϵ_1`$ as in (18), and taking into account that $`e,f|W_\delta |e,f`$ is a continuous function of (the coefficients of) $`|e,f`$ and that the set in which we are taken the infimum is compact, we obtain $`ϵ_1>0`$. Then we obviously have that $`W_1`$ fulfills properties (I) and (III). On the other hand, $`\delta _{W_1}a(\delta _{P_1}+\delta ^T_{Q_1})ϵ_1<0`$, since $`P_1\delta =Q_1\delta ^T=0`$. Thus, $`W_1`$ detects a PPTES, and therefore, according to Theorem 3 is non–decomposable.$`\mathrm{}`$
Note that Lemma 6 provides an important generalization of the method of Terhal , based on the use of unextendible product bases . Our method works in Hilbert spaces of arbitrary dimensions, and in particular when $`dim(H_A)=2`$ (in $`2\times N`$ dimensional systems) for which unextendible product basis do not exist. By combining Lemma 6 and the optimization procedure introduced earlier, we obtain a way of creating nd–OEW. Once we have $`W_1`$ we find $`p_{W_1}`$ and $`p_{W_1^T}`$. We denote by $`P_2`$ and $`Q_2`$ the projector operators orthogonal to these two sets, respectively,
$$ϵ_2=\underset{|e,f}{inf}\frac{e,f|W_1|e,f}{e,f|P_2+Q_2^T|e,f},$$
(19)
and $`W_2W_1ϵ_2(P_2+Q_2^T)`$. According to Lemma 2b we have that $`W_2`$ is nd–finer than $`W_1`$. Now we can define $`p_{W_2},p_{W_2^T},P_3,Q_3`$ and $`W_3`$ in the same way, and continue in this vein until for some $`k`$, $`ϵ_k=0`$. If $`W_k`$ is not yet optimal, we still have to find other projectors such that we can optimize as explained in the previous subsections.
In Section VII we illustrate this method with a family of edge PPTES from Ref. . In fact, as we will mention in that Section, we have checked that the optimization method typically works as well by starting with three random vectors, and following a similar procedure to the one indicated here. This means that in our construction method we do not need in practice to start from a “edge” PPTES.
### D Sufficient condition for PPTES
In this subsection we use the results derived in the previous one to construct a sufficient criterion for non–separability of PPTES. As shown in Ref. , given an operator $`\rho 0`$, with $`\rho ^T0`$, we can always decompose it in the form
$$\rho =\rho _s+\delta ,$$
(20)
where $`\rho _s`$ is separable and $`\delta `$ is an “edge” PPTES. More details concerning this decomposition, and in particular its canonical optimal form are presented in Appendix B. In this section we use this decomposition together with the following
Lemma 7: Given a non–separable operator $`\rho =\rho _s+\delta `$, where $`\rho _s0`$ is separable then for all EW, $`W`$, such that $`\rho _W<0`$ we have that $`\delta _W<0`$.
Proof: Obvious from the definition of EW. $`\mathrm{}`$
Lemma 7 tells us that if $`\rho `$ is non–separable, then there must exist some EW that detect both $`\delta `$ and $`\rho `$. Actually, it is clear that there must exist an OEW with that property. In particular, if $`\rho ^T0`$, it must be a nd–OEW. In the previous subsection we have shown how to build them out of “edge” PPTES. Thus, given $`\rho `$ we can always decompose it in the form (20), construct an OEW that detects $`\delta `$ and check whether it detects $`\rho `$. In that case, we will have that $`\rho `$ is non–separable. Thus, this provides a sufficient criterion for non–separability.
We stress the fact that for PPTES only a special class of states, namely the class of “edge” PPTES, is responsible for the entanglement properties. In fact, one should stress that very many of the examples of PPTES discussed so far in the literature belong to the class of “edge” PPTES: the $`24`$ family from , the $`nn`$ states obtained via unextendible product basis construction , the $`33`$ states obtained via the chess-board method (b), and projections of continuous variable PPTES onto finite dimensional subspaces (c).
## VI Positive maps
It is known that PM allow for necessary and sufficient conditions for separability (or, equivalently, entanglement) of bipartite mixed states . PM’s have been also applied in the context of distillation of entanglement and information theoretic analysis of separability . In this Section we will use the isomorphism between operators and linear maps to extend the properties derived for witnesses to PM . We will first review some of the definitions and properties of linear maps.
Let us consider a linear map $`:B(H_A)B(H_C)`$. We say that $``$ is positive if for all $`YB(H_A)`$ positive, $`(Y)0`$. One can extend a linear map as follows. Given $`:B(H_A)B(H_C)`$, we define its extension $`1_B:B(H_A)B(H_B)B(H_C)B(H_B)`$ according to $`1_C(_iY_iZ_i)=_i(Y_i)Z_i`$, where $`Y_iB(H_A)`$ and $`Z_iB(H_B)`$. A linear map is completely positive if all extensions are positive. The classification and characterization of positive (but not completely positive) maps is an open question (see, e.g. Ref. ).
An example of positive (but not completely positive) map is transposition (in a given basis $`O_A`$); that is, the map $`_T`$ such that $`_T(Y)=Y^T`$. The corresponding extension is the partial transposition . A map $``$ is called decomposable if it can be written as $`=_1+_2_T`$, where $`_{1,2}`$ are completely positive.
One can relate linear maps with linear operators in the following way. We will assume $`d_A\mathrm{dim}(H_A)\mathrm{dim}(H_C)`$, but one can otherwise exchange $`H_A`$ by $`H_C`$ in what follows. Given $`XB(H_AH_C)`$ and an orthonormal basis $`O_A=\{|k\}_{k=1}^{d_A}`$ in $`H_A`$, we define the linear map $`_X:B(H_A)B(H_C)`$ according to
$$(Y)=\mathrm{tr}_A(X^{T_A}Y),$$
(21)
for all $`YB(H_A)`$, where $`\mathrm{tr}_A`$ denotes the trace in $`H_A`$ and the partial transpose is taken in the basis $`O_A`$. Similarly, given a linear map we can always find an operator $`X`$ such that (21) is fulfilled. For instance, if we choose $`T=(|\mathrm{\Psi }\mathrm{\Psi }|)^{T_A}`$, where
$$|\mathrm{\Psi }=\underset{k=1}{\overset{d_A}{}}|k_A|k_C,$$
(22)
then the corresponding map $`_T`$ is precisely the transposition in the basis $`O_A`$.
Given a linear map $`_X`$, one can easily show the following relations: (a) $`_X`$ is completely positive iff $`X0`$; (b) $`_X`$ is positive but not completely positive iff $`X`$ is an EW \[except for the normalization condition (III)\]; (c) $`_X`$ is decomposable iff $`X`$ is decomposable. Thus, the problem of studying and classifying PM is very much related to the one of EW. Furthermore, PM can be also used to detect entanglement . Let us consider the extension $`\overline{}_X_X1:B(H_A)B(H_B)B(H_C)B(H_B)`$, where we take $`d_B\mathrm{dim}(H_B)=\mathrm{dim}(H_C)`$. Then we have that given $`\rho B(H_AH_B)`$,
$$\rho _X=\mathrm{\Psi }|\overline{}_X(\rho )|\mathrm{\Psi },$$
(23)
where
$$|\mathrm{\Psi }=\underset{k=1}{\overset{d_B}{}}|k_C|k_B.$$
(24)
Thus, if an EW, $`W`$, detects $`\rho `$, then $`\overline{}_W(\rho )`$ is a non–positive operator. Consequently, $`\rho 0`$ is entangled iff there exists a PM such that acting on $`\rho `$ gives a non–positive operator. In that case we say that the PM “detects” $`\rho `$. Actually, PM are “more efficient” in detecting entanglement than EW. The reason is that it may happen that $`\overline{}_X(\rho )`$ is non–positive but still $`\rho _X0`$.
It is convenient to define finer and optimal PM as for EW. That is, given two PM, $`_{1,2}`$, we say that $`_2`$ is finer than $`_1`$ if it detects more. We say that a PM, $``$, is optimal if there exists no one that is finer. In the same way we can define nd–finer and nd–optimal.
The results presented in the previous sections can be directly translated to PM given the following fact.
Lemma 8: If $`W_2`$ is finer (nd–finer) than $`W_1`$ then $`_{W_2}`$ is finer (nd–finer) than $`_{W_1}`$.
Proof: Using Lemma 2 we can write $`W_1=(1ϵ)W_2+ϵP`$. According to (21) we have that $`_{W_1}=(1ϵ)_{W_2}+ϵ_P`$. Since $`_P(\rho )0`$ for all $`\rho 0`$, we have that $`_{W_2}`$ is finer than $`_{W_1}`$. Using Lemma 2b we can also prove that it is nd–finer. $`\mathrm{}`$
¿From this lemma it follows that optimizing EW implies optimizing PM. In fact, the constructions that we have given in the Section VC can be viewed as ways of constructing non–decomposable PM. In fact, since the method works for $`dim(H_A)=2`$, the resulting PM $`:B(H_A)B(H_C)`$ has a minimal “qubit” domain, or – equivalently – minimal hermitian conjugate codomain. Up to our knowledge, our method is the first one that permits to construct non–decomposable PM with these characteristics.
## VII Illustration
In this section we explicitly give construct a nd–OEW out of edge PPTES. We use, as a starting point, the family of PPTES introduced in ).
### A Family of “edge” PPTES
We consider $`H_A=\mathrm{I}\mathrm{C}^2`$ and $`H_B=\mathrm{I}\mathrm{C}^4`$, and denote by $`\{|k\}_{k=0}^{d_\alpha }`$ ($`\alpha =A,B`$) an orthonormal basis in these spaces, respectively. Most of the time we will write the operators in those bases; that is, as matrices. For operators acting in $`H_AH_B`$ we will always use the following order $`\{|0,0,|0,1,\mathrm{},|1,0,|1,1,\mathrm{}\}`$. On the other hand, all partial transposes will be taken with respect to $`H_B`$.
We consider the following family of positive operators
$`\rho _b={\displaystyle \frac{1}{7b+1}}\left(\begin{array}{cccccccc}b& 0& 0& 0& 0& b& 0& 0\\ 0& b& 0& 0& 0& 0& b& 0\\ 0& 0& b& 0& 0& 0& 0& b\\ 0& 0& 0& b& 0& 0& 0& 0\\ 0& 0& 0& 0& \frac{1+b}{2}& 0& 0& \frac{\sqrt{1b^2}}{2}\\ b& 0& 0& 0& 0& b& 0& 0\\ 0& b& 0& 0& 0& 0& b& 0\\ 0& 0& b& 0& \frac{\sqrt{1b^2}}{2}& 0& 0& \frac{1+b}{2}\end{array}\right),`$ (33)
where $`b[0,1]`$. For $`b=0,1`$ those states are separable, whereas for $`0<b<1`$, $`\rho _b`$ is an “edge” PPTES. This can be shown by checking directly that they violate the range criterion of Ref. , i.e. the definition given in Section IVC.
If we take the partial transpose in the basis $`\{|k\}`$, the density operators $`\rho _b`$ have the property that $`\rho _b^T=U_B\rho _bU_B^{}`$ with $`U_B=(\sigma _x)_{03}(\sigma _x)_{12}`$. Here, the subscript $`ij`$ denotes the subspace, $`_{Bij}_B`$ spanned by $`\{|i,|j\}`$ and $`\sigma _x`$ is one of the Pauli-operators. Note that $`U_B`$ is a real unitary operator acting only on $`H_B`$. This immediately implies that
$$\stackrel{~}{\rho }_b^T=\stackrel{~}{\rho }_b,$$
(34)
where $`\stackrel{~}{\rho }_b=V_B\rho _bV_B^{}`$ and $`V_B=1\sqrt{2}[(1𝐥+i\sigma _x)_{03}(1𝐥+i\sigma _x)_{12}]`$. We will use the property (34) to simplify the problem of constructing the nd–OEW. Thus, we will concentrate from now on the operators $`\stackrel{~}{\rho }_b`$ . Obviously, $`\stackrel{~}{\rho }_b`$ is an “edge” PPTES for $`1>b>0`$.
The projector onto the kernel of $`\stackrel{~}{\rho }_b`$, $`P_1`$, is invariant under the transformation $`T_{AB}=T_AT_B`$, where
$`T_A`$ $`=`$ $`\left(\begin{array}{cc}1& 0\\ 0& e^{i2\pi /3}\end{array}\right),`$ (37)
$`T_B`$ $`=`$ $`\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& \mathrm{cos}(2\pi /3)& \mathrm{sin}(2\pi /3)& 0\\ 0& \mathrm{sin}(2\pi /3)& \mathrm{cos}(2\pi /3)& 0\\ 0& 0& 0& 1\end{array}\right).`$ (42)
Note that $`T_B`$ is a real matrix. Later on we will need its eigenstates with real coefficients; they are $`|0\pm |3`$. Note also that $`T_{AB}^3=1𝐥`$.
### B Construction of nd–EW’s
We use now the methods developed in Section V to obtain a nd–OEW starting from $`\stackrel{~}{\rho }_b`$. That is, we define $`W_b=P_1+P_1^T`$, where $`P_1`$ is the projector onto $`K(\stackrel{~}{\rho }_b)=K(\stackrel{~}{\rho }_b^T)`$. Our procedure consists of first subtracting the identity to obtain $`W_1=W_bϵ_11𝐥`$. Then, we subtract $`P_2+Q_2^T`$, $`P_3+Q_3^T`$, etc. In the n–th step we will have
$$W_n=W_{n1}ϵ_n(P_n+Q_n^T),$$
(43)
where $`P_n`$ ($`Q_n`$) is the projector orthogonal to the space spanned by $`P_{W_{n1}}`$ ($`P_{W_{n1}^T}`$). We will use the symmetries of $`\stackrel{~}{\rho }_b`$ to better understand the structure of $`W_n`$.
$`W_n=W_n^T`$. We can prove this by induction. First, it is clear that $`W_1=W_1^T`$. Let us now assume that $`W_{n1}=W_{n1}^T`$. Then we show that $`W_n=W_n^T`$. For that, we just have to show that the subspace spanned by $`P_{W_{n1}}`$ is the same that the one spanned by $`P_{W_{n1}^T}`$, so that $`Q_n=P_n`$. But this is clear since $`W_{n1}=W_{n1}^T`$. $`\mathrm{}`$
$`T_{AB}W_nT_{AB}^{}=W_n`$. We prove this by induction. First, for $`W_1=P_1+P_1^Tϵ_11𝐥`$ we have that $`T_{AB}W_1T_{AB}^{}=T_{AB}P_1T_{AB}^{}+T_{AB}P_1^TT_{AB}^{}ϵ_11𝐥=W_1`$, since $`T_{AB}P_1^TT_{AB}^{}=(T_{AB}P_1T_{AB}^{})^T`$ (given the fact that $`T_B`$ is real) and $`P_1`$ is invariant under $`T_{AB}`$. Then, let us assume that $`T_{AB}W_{n1}T_{AB}^{}=W_{n1}`$. In order to show that $`T_{AB}W_nT_{AB}^{}=W_n`$ we just have to show that $`P_n`$ is invariant under $`T_{AB}`$, or, equivalently, that the subspace spanned by $`P_{W_{n1}}`$ is invariant under $`T_{AB}`$. But this follows immediately from the fact that $`T_{AB}W_{n1}T_{AB}^{}=W_{n1}`$. $`\mathrm{}`$
Starting the property (a) it follows that the vectors $`|e,fP_{W_n}`$ will have $`|f`$ real (unless we have degeneracies). This can be seen by noticing that those vectors minimize $`e,f|W_n|e,f`$; defining $`W_ee|W_n|e`$, we have that $`W_e^T=W_e=W_e^{}`$ is symmetric, and therefore the eigenstate corresponding to its minimum eigenvalue can be chosen to be real. On the other hand, starting from the property (b) it follows that if $`|e,fP_{W_n}`$ then $`T_{AB}^{}|e,f,T_{AB}^2|e,fP_{W_n}`$. According to that, we will typically have two kinds of product vectors in $`P_{W_n}`$:
$`|e,f`$ is an eigenstate of $`T_{AB}^{}`$ with $`|f`$ real: There are only $`4`$ possible product vectors which fulfill these conditions: $`\{|0,|1\}\{|0+|3,|0|3\}`$.
$`|e,f`$ is not an eigenstate of $`T_{AB}^{}`$: Then, we will also have: $`T_{AB}^{}|e,f`$ and $`(T_{AB}^{})^2|e,fP_W`$.
We have carried out this procedure for $`\stackrel{~}{\rho }_b`$ and found nd–OEW for each $`b`$. We find that for the optimal EW we have two vectors of the kind (1) and six of the kind (2). In total we find eight product vectors in $`P_W`$, which span the whole Hilbert space and therefore the corresponding EW are optimal (see Corollary 2). This means that any operator of the form (2) with $`|e_k,f_kP_W`$ the product vectors we have found, and $`p_k>0`$ will be a full range separable density operator that lies on the boundary between separable and PPTES. Up to our knowledge, this constitutes the first example of those operators . We have also created the PM corresponding to the nd–OEW, which are the first examples of non–decomposable PM with minimal “qubit” domain, or – equivalently – minimal hermitian conjugate codomain.
In Fig. 1 we show for which $`b^{}`$ $`\stackrel{~}{\rho }_b^{}`$ is still detected by the nd–OEW created out of $`\stackrel{~}{\rho }_b`$. We find that for a given $`b`$, the optimal witness that we create detects all $`\stackrel{~}{\rho }_{\stackrel{~}{b}}`$ for $`\stackrel{~}{b}b^{}`$. Thus, in the figure we plot $`b^{}`$ as a function of $`b`$. As explained above, the corresponding positive map detects more than the witness itself. In the figure one can also see how much is detected by the positive map.
Obviously, the witnesses that we create do not only detect the density operators $`\stackrel{~}{\rho }_b`$. For instance one can check how much one can add the identity to certain $`\stackrel{~}{\rho }_b`$ but still keeping the state entangled. That is for which $`\lambda `$, $`\stackrel{~}{\rho }_b+\lambda 1𝐥`$ is still detected by the witness. This is shown in the following figure.
Finally, let us note that we have observed using numerical calculations that if one starts with a random projector, $`P`$ of rank $`3`$, and optimizes the decomposable operator $`WP+P^{T_B}`$ in the same way as the one described here, then one will end up with a nd–OEW $`\stackrel{~}{W}`$, where $`p_{\stackrel{~}{W}}`$ is complete. This means that in order to create nd–OEW one does not need to know in practice an edge PPTES. In another words, optimization itself is a way to reach nondecomposableness.
### C Analytical procedure
In this subsection we will present an analytical way to create nd–EW’s. Furthermore we will given an example of such a witness, which detects $`\rho _b`$ for all $`b(0,1)`$. From Fig. $`1`$ we see that the witness which detects most is the one we created out of $`\stackrel{~}{\rho }_b`$, where $`b`$ is very close to $`1`$. We will work with the original $`\rho _b`$ (33).
We consider two hermitian operators $`A`$ and $`B`$, with $`A`$ positive on product vectors, i.e., $`e,f|A|e,f0`$, whereas $`B`$ does not have to. As before we denote by $`P_A`$ ($`P_B`$) the (not necessarily complete) set of product vectors on which $`A`$ ($`B`$) vanishes. We require that for all $`|e,fP_A`$, $`e,f|B|e,f0`$. Then we define $`W(x)\frac{1}{x}(A+xB)`$ for any real $`x`$. So we have the following
Lemma 9: If $`\text{lim}_{x0}\rho _{W(x)}<0`$ then $`\rho `$ is entangled.
Proof: We prove that $`\text{lim}_{x0}e,f|W(x)|e,f0`$. This implies that if $`\rho `$ is separable, then $`\text{lim}_{x0}\rho _{W(x}0`$. Let us therefore distinguish two cases: (i) if $`|e,fP_A`$ then we have that $`\text{lim}_{x0}e,f|W(x)|e,f=e,f|B|e,f`$, which is, per assumption, positive. (ii) $`|e,fP_A`$ then we have $`\text{lim}_{x0}e,f|W(x)|e,f=\text{lim}_{x0}\frac{a}{x}+b`$, where $`a=e,f|A|e,f>0`$ and $`b=e,f|B|e,f`$. Thus this limit tends to infinity, which proves the statement.$`\mathrm{}`$
Note that $`W(x)`$ is not an EW since it is not necessarily positive on product vectors. However, one can make it positive by adding the identity operator to convert it into an EW.
Corollary 6: Given any $`x_0>0`$, then $`W(x_0)\frac{1}{x_0}(A+x_0B)+\lambda _{x_0}1𝐥`$ , with $`\lambda _{x_0}=\text{ min}_{|e,f}e,f|\frac{1}{x_0}(A+x_0B)|e,f`$ is an EW.
Let us now illustrate how we can use Lemma 9 to detect all the states $`\rho _b`$. We define
$`A=\left(\begin{array}{cccccccc}0& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 2& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 2& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\end{array}\right),`$ (52)
$`B=\left(\begin{array}{cccccccc}1& 0& 0& 1& 0& 2& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 2\\ 1& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 1\\ 2& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 2& 0& 1& 0& 0& 1\end{array}\right).`$ (61)
One can easily show that $`A=|\psi \psi |+(|\varphi \varphi |)^{T_B}`$, where $`|\psi =|01|12`$ and $`|\varphi =|02|11`$. Thus this operator is positive on product vectors, since it is decomposable. Let us now use unnormalized states in order to present the set of product vectors on which $`A`$ vanishes, i.e. $`P_A`$. $`P_A=P_{A_1}P_{A_2}`$, where $`P_{A_1}=\{(|0+\alpha |1)(x|0+y|3)\alpha ,x,y\}`$ and $`P_{A_2}=\{(|0+e^{i\mathrm{\Phi }}|1)[x|0+y|3+z(|1+e^{i\mathrm{\Phi }}|2]\mathrm{\Phi },x,y,z\}`$. The operator $`B`$ has to be positive on those product vectors, i.e., $`|e,fP_A`$, $`e,f|B|e,f0`$. In order to show that this is indeed like that let us distinguish the two cases: $`|e,fP_{A_1}`$ and $`|e,fP_{A_2}`$. In the first case we have that
$`e|B|e=\left(\begin{array}{cccc}1+|\alpha |^2& 2\alpha & 0& 1|\alpha |^2\\ 2\alpha ^{}& 1+|\alpha |^2& 0& 0\\ 0& 0& 1+|\alpha |^2& 2\alpha \\ 1|\alpha |^2& 0& 2\alpha ^{}& 1+|\alpha |^2\end{array}\right)`$ (66)
and so $`e,f|B|e,f=|x+y|^2+|\alpha |^2|xy|^20`$. If $`|e,fP_{A_2}`$ then
$`e|B|e=\left(\begin{array}{cccc}2& 2e^{i\mathrm{\Phi }}& 0& 0\\ 2e^{i\mathrm{\Phi }}& 2& 0& 0\\ 0& 0& 2& 2e^{i\mathrm{\Phi }}\\ 0& 0& 2e^{i\mathrm{\Phi }}& 2\end{array}\right)`$ (71)
which is a positive operator and so $`e,f|B|e,f0`$. So those two operators $`A`$ and $`B`$ fulfill all the required properties. Furthermore one can show that $`\rho _b_A=0`$ and $`\rho _b_B<0`$ for all $`0<b<1`$. Thus we have that $`\text{lim}_{x0}W(x)\rho _b<0`$ for all $`0<b<1`$, where we defined $`W(x)=\frac{1}{x}(A+xB)`$.
As mentioned above we can use now $`W(x)`$ in order to create other PPTES just by adding product vectors on which $`W(x)`$ vanishes. To find the product vectors we can add, all we need to do is to determine the intersection between $`P_A`$ and $`P_B`$. Since $`P_B=\{(|0+e^{i\varphi }|1)[a(|0+e^{i\varphi }|1+b(|2+e^{i\varphi }|3)]\varphi ,a,b\}`$ we have that $`SP_AP_B=P_{A_2}P_B=\{(|0+e^{i\varphi }|1)(|0+e^{i\varphi }|1+e^{i2\varphi }|2+e^{i3\varphi }|3)\varphi \}`$. Note that $`S`$ spans a $`5`$ dimensional subspace and that the orthogonal subspace is spanned by the vectors $`\{|02+|13,|01+|12,|00+|11\}`$.
## VIII Conclusions
Entanglement witnesses allow us to study the separability properties of density operators. We have defined OEW, which are those that detect entanglement in an optimal way. We have given necessary and sufficient conditions for an EW to be optimal, and we have shown a way to construct them. We have also concentrated on nd–EW, which are those that detect PPTES. We have extended the definitions of optimality and the optimization procedure to those EW. It turns out that one can optimize nd–EW by subtracting decomposable operators. We have also given an explicit method to construct nd–EW starting from “edge” PPTES. We have also mentioned that this method works by starting out from random operators. We have extended our techniques to PM, and therefore given a method to systematically construct non–decomposable positive maps. We have illustrated our methods with a family of “edge” PPTES acting on $`\mathrm{I}\mathrm{C}^2\mathrm{I}\mathrm{C}^4`$. The corresponding PM constitute the first examples of PM with minimal “qubit” domain, or – equivalently – minimal hermitian conjugate codomain. We have also constructed the first examples of separable states of full range that lie on the boundary between separable and PPTES. These states can be used for experimental realization of PPTES .
In this paper we have also introduced the “edge” PPTES, which violate the range criterion of separability. As shown in Appendix B, the “edge” PPTES allow us to construct a canonical form of PPTES in Hilbert spaces of arbitrary dimensions. They also allow us to give a novel sufficient condition for non–separability which applies to operators with positive partial transpose. It is based on the fact that among all PM (or EW) only the subset $`\{\mathrm{\Lambda }_{edge}\}`$ of those PM that detect edge PPTES are needed to study the separability of PPTES. This opens many interesting questions. Is it possible that in the set $`\{\mathrm{\Lambda }_{edge}\}`$ there is some map that is globally finer than the transposition? In another words, is there a map detecting the entanglement of all the states with non–positive partial transpose? What is the minimal subset of $`\{\mathrm{\Lambda }_{edge}\}`$ providing such condition? Is it finite?
Finally, let us consider the implications of the our results for the very interesting problem of locality of PPTES. There is a conjecture that those states can be local in the sense that they admit a local hidden variable (LHV) model for any set of possible local measurements. The problem is not trivial given the fact that it may be important to take into account the role of sequential measurements and the possible existence of many copies. Quite recently it has been shown that PPTES satisfy Bell-type of inequalities introduced by Mermin . It is not difficult to convince oneself that the set of states admitting LHV model for any fixed type of measurements is a convex set. Furthermore, extending the reasoning from it is easy to see that the set of separable states admits LHV models for any possible set of measurements. Hence, taking into account the results of this paper it follows that in order to prove, or to disprove locality of PPTES it is enough to study only “edge” PPTES.
Note that the “edge” states have typically very small rank (the minimal rank is four in $`33`$ systems, see Ref. ). There have been no examples of LHV models for states of low rank, so far. Thus, perhaps completely new techniques will be needed to study this problem. In this case the most symmetric PPTES provided recently (c) seem to be the best suitable for the first test.
## IX Acknowledgments
This was has been supported in part by the Deutsche Forschungsgemeinschaft (SFB 407 and Schwerpunkt ”Quanteninformationsverarbeitung”), the DAAD, the Austrian Science Foundation (SFB “control and measurement of coherent quantum systems”), the ESF PESC Programm on Quantum Information, TMR network ERB–FMRX–CT96–0087, the IST Programme EQUIP, and the Institute for Quantum Information GmbH.
## A Optimality of EW
In this Appendix we study necessary and sufficient conditions for an EW to be optimal. According to Theorem 1 of Section III we have that an EW, $`W`$, is optimal iff no positive operator can be subtracted from $`W`$ while keeping property (I). This condition can be reexpressed in terms of the infimum of some scalar products in Lemma 4. This infimum is, in general, difficult to calculate (at least analytically). In this Section we will give a different method to determine whether an EW is optimal or not. This method will turn out to be very simple for the case in which $`dim(H_A)=2`$. The idea is to find the conditions such that a given operator $`P0`$ can/cannot be subtracted from an EW. This will give us automatically a criterion to determine when $`W`$ is optimal.
In all this appendix we will use that given an EW, $`W`$, and an operator $`P0`$ we say that $`P`$ cannot be subtracted from $`W`$ if for all $`\lambda >0`$, $`W\lambda P`$ does not fulfill (I). In other words, there exist $`|e(\lambda )H_A`$ and $`|f(\lambda )H_B`$ such that
$$e(\lambda ),f(\lambda )|(W\lambda P)|e(\lambda ),f(\lambda )<0.$$
(A1)
Note that $`e(\lambda ),f(\lambda )|P|e(\lambda ),f(\lambda )`$ must be strictly positive, so that (A1) can be expressed as
$$\underset{\lambda 0}{lim}\frac{e(\lambda ),f(\lambda )|W|e(\lambda ),f(\lambda )}{e(\lambda ),f(\lambda )|P|e(\lambda ),f(\lambda )}=0.$$
(A2)
In the first subsection we will introduce some definitions and notation. In the second one we give a method to determine the set of product vectors $`P_W`$, on which $`W`$ vanishes. In the third subsection we find a necessary and sufficient condition under which an operator cannot be subtracted from an EW. We will see that there must exist a vector $`|e_0,f_0P_W`$, some other vectors $`|e_1`$ and $`|f_1`$, and certain phases $`\varphi _{e,f}`$ and $`\theta `$ such that some quantity is zero. In the next subsection we will see that the problem can be reduced to finding only the vectors $`|e_{0,1}`$ and $`|f_{0,1}`$. Finally, we will show that if $`dim(H_A)=2`$ we just have to find $`|e_0`$ and $`|f_0`$, which is very simple.
### 1 Definitions and notation
In order to prove the results of this appendix in a compact and readable form we have made an extensive numbers of definitions.
We will always denote by $`|e_0,f_0`$ a product vector in $`P_W`$, and by $`|e_1H_A`$ and $`|f_1H_B`$ two vectors orthogonal to $`|e_0`$ and $`|f_0`$, respectively. We will use the following notation:
$$W_{i,j}^{k,l}=e_i,f_j|W|e_k,f_l,(i,j,k,l=0,1).$$
(A3)
and we will write
$`W_{1,0}^{0,1}`$ $`=`$ $`|W_{1,0}^{0,1}|e^{i\varphi _0}`$ (A5)
$`W_{0,0}^{1,1}`$ $`=`$ $`|W_{0,0}^{1,1}|e^{i\varphi _1}.`$ (A6)
We will also define the following operators:
$`w_{i,j}^e`$ $``$ $`e_i|W|e_j,`$ (A8)
$`w_{i,j}^f`$ $``$ $`f_i|W|f_j.`$ (A9)
The following vectors will be used in the context of Eq. (A2):
$`|e(ϵ)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{1+|\mathrm{cos}(\theta )ϵ|^2}}}(|e_0+ϵ\mathrm{cos}(\theta )e^{i\varphi _e}|e_1),`$ (A11)
$`|f(ϵ)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{1+|\mathrm{sin}(\theta )ϵ|^2}}}(|f_0+ϵ\mathrm{sin}(\theta )e^{i\varphi _f}|f_1),`$ (A12)
where $`ϵ`$ is a real number, and $`\varphi _{e,f}[0,\pi )`$ and $`\theta [0,\pi /2]`$ are certain constants. Given a product vector $`|e(ϵ),f(ϵ)`$ and an operator, $`W`$, we will expand $`e(ϵ),f(ϵ)|W|e(ϵ),f(ϵ)`$ by collecting terms with the same powers in $`ϵ`$; that is, except for a normalization constant,
$$e(ϵ),f(ϵ)|W|e(ϵ),f(ϵ)\underset{i=1}{\overset{4}{}}ϵ^iA_i(W),$$
(A13)
where
$`A_0(W)`$ $`=`$ $`W_{0,0}^{0,0},`$ (A15)
$`A_1(W)`$ $`=`$ $`2\mathrm{R}\mathrm{e}\left[\mathrm{cos}(\theta )e^{i\varphi _e}W_{0,0}^{1,0}+\mathrm{sin}(\theta )e^{i\varphi _f}W_{0,0}^{0,1}\right],`$ (A16)
$`A_2(W)`$ $`=`$ $`\mathrm{cos}^2(\theta )W_{1,0}^{1,0}+\mathrm{sin}^2(\theta )W_{0,1}^{0,1}`$ (A19)
$`+2\mathrm{sin}(\theta )\mathrm{cos}(\theta )`$
$`\times \mathrm{Re}\left[e^{i(\varphi _e\varphi _f)}W_{1,0}^{0,1}+e^{i(\varphi _e+\varphi _f)}W_{0,0}^{1,1}\right],`$
$`A_3(W)`$ $`=`$ $`2\mathrm{sin}(\theta )\mathrm{cos}(\theta )`$ (A21)
$`\times \mathrm{Re}\left[\mathrm{cos}(\theta )e^{i\varphi _f}W_{1,0}^{1,1}+\mathrm{sin}(\theta )e^{i\varphi _e}W_{0,1}^{1,1}\right],`$
$`A_4(W)`$ $`=`$ $`\mathrm{sin}^2(\theta )\mathrm{cos}^2(\theta )W_{1,1}^{1,1}.`$ (A22)
On the other hand, we will define
$$|\mathrm{\Psi }_{0,1}\mathrm{sin}(\theta )e^{i\varphi _f}|e_0,f_1+\mathrm{cos}(\theta )e^{i\varphi _e}|e_1,f_0.$$
(A23)
Finally, the following quantity will play an important role in determining whether there exist vectors and parameters for which (A2):
$$X(W)W_{1,0}^{1,0}W_{0,1}^{0,1}(|W_{1,0}^{0,1}|+|W_{0,0}^{1,1}|)^2.$$
(A24)
### 2 Determining $`P_W`$
As stated in Lemma 3, not every positive operator $`P`$ can be subtracted from an EW, $`W`$; it must vanish on $`P_W`$. Thus, in order to choose $`P`$ one has to know the set $`P_W`$. In this subsection we give a method to determine it.
We start by characterizing the vectors in $`P_W`$:
Lemma A1: Given an operator $`W`$ satisfying (I), then $`|e_0,f_0P_W`$ iff
$`e_0|W|e_0|f_0`$ $`=`$ $`0,`$ (A26)
$`f_0|W|f_0|e_0`$ $`=`$ $`0,`$ (A27)
Proof: (If) We just apply $`f_0|`$ to Eq. (A26). (Only if) Since $`W`$ fulfills (I) then $`W_{e_0}e_0|W|e_0`$ must be positive. Thus, $`f_0|W_{e_0}|f_0=0`$ implies Eq. (A26). In the same way we obtain Eq. (A27). $`\mathrm{}`$
In practice, for a given $`W`$ the set $`P_W`$ can be found as follows. Due to the fact that $`W`$ is an EW we have that for any $`|eH_A`$, $`W_ee|W|e`$ must be a positive operator (i.e. $`f|W_e|f0`$ for all $`|fH_B`$). Thus, the determinant $`\mathrm{det}(W_e)0`$. According to Lemma A1, this determinant is zero iff there exists some $`|f_0H_B`$ such that $`f_0|W_{e_0}|f_0=0`$, i.e., if $`|e_0,f_0P_W`$. That is, the determinant as a function of $`|e`$ has a minimum (which is zero) at $`|e_0`$. We can use this fact to find $`|e_0`$. Then, we can easily obtain $`|f_0`$ via Eq. (A26). We can expand an unnormalized state $`|e`$ in an orthonormal basis $`\{|k\}`$ as
$$|e=\underset{k=1}{\overset{dim(H_A)}{}}c_k|k,$$
(A28)
and impose that the corresponding determinant is zero. This gives us a polynomial equation for the coefficients $`c_k`$, i.e.
$$P(c_k,c_k^{})=0.$$
(A29)
We also impose that, given the fact that the determinant is a minimum,
$$\frac{}{c_k}P(c_k,c_k^{})=\frac{}{c_k^{}}P(c_k,c_k^{})=0,$$
(A30)
which also give a set of polynomial equations. These equations can be solved using the method mentioned in Ref. .
### 3 Necessary and sufficient conditions for subtracting an operator
In this subsection we give a necessary and sufficient condition for an operator $`P`$ to be subtractable from an EW. We start out by giving some properties of the coefficients $`A(W)`$ defined above (A13).
Lemma A2: Given $`W`$ satisfying (I) and $`|e_0,f_0P_W`$, then for all $`|e_1H_A`$ and $`|f_1H_B`$ we have
$`A_0(W)=A_1(W)=0`$.
$`A_2(W)0`$.
If $`A_2(W)=0`$ then $`A_3(W)=0`$.
Proof: (i) It is a direct consequence from Lemma A1. In order to prove (ii–iii) we use the fact that $`W`$ satisfies (I). We define $`|e(ϵ)`$ and $`|f(ϵ)`$ as in (A 1). We impose that $`e(ϵ),f(ϵ)|W|e(ϵ),f(ϵ)0`$. Using the expansion (A13) and taking into account (i), we have $`A(ϵ)A_2(W)+ϵA_3(W)+ϵ^2A_4(W)0`$ for all $`ϵ`$. This automatically implies (ii), since otherwise for sufficiently small $`ϵ`$ we would have $`A(ϵ)<0`$. It also implies (iii), since if $`A_3(W)<0`$ ($`A_3(W)>0`$) then for sufficiently small $`ϵ>0`$ ($`ϵ<0`$) we would have $`A(ϵ)<0`$. $`\mathrm{}`$
Now, we are at the position of giving a necessary and sufficient condition under which an operator cannot be subtracted from an EW:
Lemma A3: Given $`P`$ fulfilling $`PP_W=0`$, it cannot be subtracted from $`W`$ iff there exists $`|e_0,f_0P_W`$, $`|e_1|e_0`$, $`|f_1|f_0`$, $`\varphi _{e,f}`$, and $`\theta `$ such that $`A_2(W)=0`$ but $`A_2(P)0`$.
Proof: (If) We define $`|e(\lambda )`$ and $`|f(\lambda )`$ as in (A 1). Using Lemma A2(i) we have $`A_0(W)=A_0(P)=A_1(W)=A_1(P)=0`$. Using Lemma A2(iii) we have that $`A_3(W)=0`$. Thus, we can write the limit (A2) as
$$\underset{\lambda 0}{lim}\frac{\lambda ^2A_4(W)}{A_2(P)+\lambda A_3(P)+\lambda ^2A_4(P)}$$
(A31)
which obviously tends to zero given that $`A_2(P)0`$. (Only if) There exist two normalized vectors $`|\stackrel{~}{e}(\lambda )`$ and $`|\stackrel{~}{f}(\lambda )`$ (continuous functions of $`\lambda `$) fulfilling (A2). Taking the limit $`\lambda 0`$ in this expression we have that $`\stackrel{~}{e}(0),\stackrel{~}{f}(0)|W|\stackrel{~}{e}(0),\stackrel{~}{f}(0)=0`$, and therefore $`|e_0,f_0|\stackrel{~}{e}(0),\stackrel{~}{f}(0)P_W`$. This means that we can always choose $`|\stackrel{~}{e}(\lambda )=|e[ϵ(\lambda )]`$ and $`|\stackrel{~}{f}(\lambda )=|f[ϵ(\lambda )]`$ given in (A 1), where $`|e_1|e_0`$ and $`|f_1|f_0`$ are two normalized vectors, $`lim_{\lambda 0}ϵ(\lambda )=0`$, and $`e(ϵ),f(ϵ)|P|e(ϵ),f(ϵ)0`$. We use (A 1) to expand the numerator and denominator of (A2) as in (A13). According to Lemma A2(i) we have that $`A_0(W)=A_0(P)=A_1(W)=A_1(P)=0`$. Thus, we must have
$$\underset{ϵ0}{lim}\frac{A_2(W)+ϵA_3(W)+ϵ^2A_4(W)}{A_2(P)+ϵA_3(P)+ϵ^2A_4(P)}=0.$$
(A32)
This implies $`A_2(W)=0`$ and $`A_2(P)0`$. Note that if both $`A_2(W)=A_2(P)=0`$ then, according to Lemma A2(iii) we have that $`A_3(W)=A_3(P)=0`$, so that (A31) would require $`A_4(W)/A_4(P)=0`$. But this cannot be since $`A_4(W)=0`$ would imply that $`|e(ϵ),f(ϵ)P_W`$, and therefore $`e(ϵ),f(ϵ|P|e(ϵ),f(ϵ)=0`$. $`\mathrm{}`$
Finally, we show in the next lemma that condition $`A_2(P)=0`$ is equivalent to having certain vector in the kernel of $`P`$. We will use the vector $`|\mathrm{\Psi }_{0,1}`$ defined in (A23).
Lemma A4: Given a positive operator $`P`$ and a set of vectors $`|e_0,f_0K(P)`$, $`|e_1|e_0`$ $`|f_1|f_0`$, and parameters $`\varphi _{e,f}`$, and $`\theta `$ then $`A_2(P)=0`$ iff $`|\mathrm{\Psi }_{0,1}K(P)`$.
Proof: Since $`P0`$ and $`|e_0,f_0K(P)`$ we have $`P_{1,1}^{0,0}=0`$. Then, we can write $`A_2(P)=\mathrm{\Psi }_{0,1}|P|\mathrm{\Psi }_{0,1}`$, with $`|\mathrm{\Psi }`$ is defined in (A23), from which it is obvious that $`A_2(P)=0`$ iff $`|\mathrm{\Psi }_{0,1}K(P)`$. $`\mathrm{}`$
### 4 Necessary and sufficient conditions for $`A_2(W)=0`$
The previous lemmas tell us that we cannot subtract a given operator $`P`$ provided we can find some vectors and parameters such that $`A_2(W)=0`$. The task of finding these vectors is difficult, in general. Here we will give a way to check whether these vectors exist. As before, we will denote by $`|e_0,f_0`$ a vector in $`P_W`$, and by $`|e_1`$ and $`|f_1`$ two vectors orthogonal to the first two. The quantity $`X(W)`$ defined in (A24) will play an important role in determining whether there exist vectors and parameters for which $`A_2(W)=0`$. In this subsection, we will always have to choose the phases $`\varphi _{e,f}`$ that minimize $`A_2(W)`$. That is
$$e^{i(\varphi _e\varphi _f\varphi _0)}=1,e^{i(\varphi _e+\varphi _f+\varphi _1)}=1.$$
(A33)
We will denote $`\stackrel{~}{A}_2(W)`$ the value of $`A_2(W)`$ for this particular choice of phases. We have
$`\stackrel{~}{A}_2(W)`$ $`=`$ $`\mathrm{cos}^2(\theta )W_{1,0}^{1,0}+\mathrm{sin}^2(\theta )W_{0,1}^{0,1}`$ (A35)
$`2\mathrm{sin}(\theta )\mathrm{cos}(\theta )\sqrt{W_{1,0}^{1,0}W_{0,1}^{0,1}X(W)},`$
where we have used (A24).
Let us start showing that $`X(W)`$ is positive. We will use this property later on to reexpress the condition $`A_2(W)=0`$ in terms of one that is simpler to check.
Lemma A5: $`X(W)0`$.
Proof: This follows from the fact that $`A_2(W)0`$ for all values of $`\varphi _{e,f}`$. In particular, $`\stackrel{~}{A}_2(W)0`$, which according to (A35) implies $`X(W)0`$. $`\mathrm{}`$
The next lemma shows that we just have to check whether $`X(W)=0`$ if we want to see if there exist parameters $`\varphi _{e,f}`$ and $`\theta `$ such that $`A_2(W)=0`$. This first condition is therefore much more useful than the last one.
Lemma A6: $`X(W)=0`$ iff there exist $`\varphi _{e,f}^0`$ and $`\theta ^0`$ such that $`A_2(W)=0`$.
Proof: (If) Given the phase $`\theta =\theta ^0`$ we have that $`0=A_2(W)\stackrel{~}{A}_2(W)`$. Thus, $`\stackrel{~}{A}_2(W)=0`$. According to (A35) we can have two cases: (a) $`\theta _00,\pi /2`$. In that case it is obvious that $`X(W)=0`$. (b) $`\theta _0=0,\pi /2`$. In the first (second) case we must have $`W_{1,0}^{1,0}=0`$ ($`W_{0,1}^{0,1}=0`$). But this implies that $`W_{0,0}^{1,1}=W_{0,1}^{1,0}=0`$ since otherwise we could always find some other value of $`\theta `$ such that $`\stackrel{~}{A}_2(W)<0`$. Then, $`X(W)=0`$. (Only if) We choose $`\varphi _{e,f}`$ as in (A33). For this value, according to (A35) we have
$$\stackrel{~}{A}_2(W)=\left[\mathrm{cos}(\theta )\sqrt{W_{1,0}^{1,0}}\mathrm{sin}(\theta )\sqrt{W_{0,1}^{0,1}}\right]^2,$$
(A36)
which can always be zero for some particular value of $`\theta `$. $`\mathrm{}`$
Note that according to the proof of Lemma A6, if $`W_{1,0}^{1,0}=0`$ then $`A_2(W)=0`$ only for $`\theta =0`$. But in that case one can easily check that the vector $`|e(\lambda ),f(\lambda )P_W`$ \[see (A 1)\] which cannot be. Similarly, we conclude that $`W_{1,0}^{1,0}0`$ if we want $`A_2(W)=0`$. Thus, from now one we will assume that both $`W_{1,0}^{1,0}`$ and $`W_{1,0}^{1,0}`$ are not zero.
### 5 Optimality test
Thus, we can now state the steps to check whether an EW, $`W`$, can be optimized or not. (1) For each $`|e_0,f_0P_W`$ we must check whether there exist $`|e_1|e_0`$ and $`|f_1|f_0`$ such that $`X(W)=0`$. Let us denote by $`|e_{0,1}^{(i)}`$ and $`|f_{0,1}^{(i)}`$ the set of vectors fulfilling that. (2) For each of these vectors, we have to find the corresponding values of $`\varphi _{e,f}^{(i)}`$ by using (A33) and of $`\theta ^{(i)}`$ by imposing that $`\stackrel{~}{A}_2(W)=0`$ in (A36). (3) Construct $`|\mathrm{\Psi }^{(i)}`$ according to (A23). (4) See whether the space spanned by $`P_W`$ and $`\{|\mathrm{\Psi }^{(i)}\}`$ is equal to $`H_AH_B`$. If it is, then $`W`$ is optimal. If it is not, we can always find some $`|\psi `$ orthogonal to that subspace that can be subtracted from $`W`$.
### 6 Necessary and sufficient conditions for $`X(W)=0`$
The hard part of the procedure outlined before to see whether and EW is optimal is the step (1), namely to find $`|e_1`$ and $`|f_1`$ such that $`X(W)=0`$. We start out by giving a necessary and sufficient condition for $`X(W)=0`$.
Lemma A7: Given $`|e_0,f_0P_W`$, and $`|e_1|e_0`$ and $`|f_1|f_0`$, then $`X(W)=0`$ iff
$`w_{0,0}^e|f_1`$ $`=`$ $`\sqrt{{\displaystyle \frac{W_{0,1}^{0,1}}{W_{1,0}^{1,0}}}}e^{i\varphi _f}(e^{i\varphi _e}w_{1,0}^e+e^{i\varphi _e}w_{0,1}^e)|f_0,`$ (A38)
$`w_{0,0}^f|e_1`$ $`=`$ $`\sqrt{{\displaystyle \frac{W_{0,1}^{0,1}}{W_{1,0}^{1,0}}}}e^{i\varphi _e}(e^{i\varphi _f}w_{1,0}^f+e^{i\varphi _f}w_{0,1}^f)|e_0,`$ (A39)
where $`\varphi _{e,f}`$ are given in (A33).
Proof: (If) We multiply by $`f_1|`$ Eq. (A38) and take the square of the absolute value of the result. We obtain
$`W_{1,0}^{1,0}W_{0,1}^{0,1}`$ $`=`$ $`|e^{i(\varphi _e+\varphi _f)}W_{1,1}^{0,0}+e^{i(\varphi _e\varphi _f)}W_{0,1}^{1,0}|^2`$ (A40)
$``$ $`(|W_{0,0}^{1,1}|+|W_{1,0}^{0,1}|)^2.`$ (A41)
Using Lemma A5 we conclude that $`X(W)=0`$. (Only if) Since $`X(W)=0`$ and according to Lemma A5 $`X(W)0`$, then $`X(W)`$ must be a minimum with respect to $`|e_1`$ and $`|f_1`$. Taking the derivatives of $`X(W)`$ with respect to these two vectors and imposing that they vanish, one obtains (A 6). $`\mathrm{}`$
Equations (A 6) are particularly useful if the dimension of one of the Hilbert spaces is 2. Without loss of generality, let us assume that $`dim(H_A)=2`$. In that case we can choose $`|e_1`$ as the one that is orthogonal to $`|e_0`$ (with an arbitrary choice of the global phase). The determination of $`\varphi _e`$ can be done as follows. Using (A 6) we write
$$\sqrt{\frac{W_{1,0}^{1,0}}{W_{0,1}^{0,1}}}e^{i\varphi _f}|f_1=\frac{1}{w_{0,0}^e}(e^{i\varphi _e}w_{1,0}^e+e^{i\varphi _e}w_{0,1}^e)|f_0$$
(A42)
where $`1/w_{0,0}^e`$ denotes the pseudo–inverse . We can use this expression to impose
$$W_{1,0}^{0,1}e^{i(\varphi _e\varphi _f)},W_{0,0}^{1,1}e^{i(\varphi _e+\varphi _f)}<0,$$
(A43)
i.e. they are negative real numbers. We obtain that
$$e^{i2\varphi _e}f_0|w_{1,0}^e\frac{1}{w_{0,0}^e}w_{1,0}^e|f_0<0,$$
(A44)
so that we determine $`\varphi _e`$. With these results, we can prove the following necessary and sufficient condition for $`X(W)=0`$ when $`dim(H_A)=2`$.
Lemma A8: If $`dim(H_A)=2`$, given $`|e_0,f_0P_W`$, then there exists $`|e_1,f_1`$ such that $`X(W)=0`$ iff
$`f_0|\left[w_{1,1}^ew_{0,1}^e{\displaystyle \frac{1}{w_{0,0}^e}}w_{1,0}^ew_{1,0}^e{\displaystyle \frac{1}{w_{0,0}^e}}w_{0,1}^e\right]|f_0=`$ (A45)
$`2\left|f_0|w_{0,1}^e{\displaystyle \frac{1}{w_{0,0}^e}}w_{0,1}^e|f_0\right|.`$ (A46)
Proof: (If) We define
$$|f_1=\frac{1}{w_{0,0}^e}(e^{i\varphi _e}w_{1,0}^e+e^{i\varphi _e}w_{0,1}^e)|f_0$$
(A47)
where $`\varphi _e`$ is determined by the condition (A44). Using this expression to calculate $`X(W)`$ one finds that indeed $`X(W)=0`$. (Only if) Using Lemma A7 we can write $`|f_1`$ as in (A42) so that the phases $`\varphi _{e,f}`$ ensure that (A43) is fulfilled. Substituting $`|f_1`$ in the equation $`X(W)=0`$ one finds (A45). $`\mathrm{}`$
In summary, for a given $`|e_0,f_0P_W`$, in order to find whether there exist $`|e_1,f_1`$ such that $`X(W)=0`$ we just have to check the condition (A45). If it is fulfilled, we can easily find $`|f_1`$ and the phases $`\varphi _{e,f}`$ using (A43) and (A42).
## B Canonical form of PPTES
The concept of “edge” PPTES seems to play a very special role in the characterization of PPTES. In particular, in view of the criterion given in Section VD, which is based on the fact that any density operator $`\rho `$ can be decomposed into a separable part and an “edge” PPTES (20). Among all the possible decompositions there might be one for which the trace of the separable part is maximal. When it exists, such a decomposition was termed positive partial transpose best separable approximation (PPT BSA) to $`\rho `$ . It extended the idea of BSA introduced in Refs. to the case of PPTES, which were based on the method of diminishing the range of $`\rho `$ by subtracting product vectors from its range, while keeping the remainder and, at the same time, its partial transpose, positive . In this Appendix we formalize the results regarding the existence and properties of the PPT BSA. In particular, the proofs presented in the quoted papers were restricted to the case in which there exist a finite, or at most, countable number of projectors on product vectors that can be subtracted from $`\rho `$. We will extend them below to continuous families of product vectors. The Appendix is written in a self-contained way, and can be read independently of the body of the paper.
We denote by $`\mathrm{\Gamma }_\rho `$ the set of projectors on product vectors $`\{|e_\alpha ,f_\alpha e_\alpha ,f_\alpha |\}`$ such that $`|e_\alpha ,f_\alpha R(\rho )`$ and $`|e_\alpha ,f_\alpha ^{}R(\rho ^{T_B})`$. In Ref. we showed that if $`\mathrm{\Gamma }_\rho `$ is finite then there exist an optimal decomposition (PPT BSA) $`\rho =(1p)\rho _{sep}+p\delta `$ where $`\delta `$ is an “edge” PPTES, and $`p`$ is minimal. Note that PPT BSA involves the state $`\delta `$ which violates the range criterion in a rather special way, i.e. with the additional requirement that $`\mathrm{\Gamma }_\rho `$ is a finite set. It can happen that there is an uncountable family of product vectors depending on continuous parameter that can be used for subtracting projectors. In the following we will show that in such case the above result is valid.
In order to consider the case of continuous families of product vectors we first prove the following:
Lemma B1: Let $`\rho `$ will be a PPTES defined on a Hilbert space $``$ , $`\mathrm{dim}<\mathrm{}`$. Then the set of product vectors $`\mathrm{\Gamma }_\rho `$ is compact.
Proof : Obviously $`\mathrm{\Gamma }_\rho `$ is a bounded set in finite-dimensional space, so it is enough to show that it is closed. Consider any sequence $`|g_n,h_n|\varphi `$, $`|g_n,h_nR(\rho )`$, $`|g_n,h_n^{}R(\rho ^{T_B})`$. The limit vector must: (i) respect the condition of orthogonality to $`K(\rho )`$ \[i. e. they must belong to $`R(\rho )`$\], (ii) belong to the sphere (i. e. set of all vectors $`|\varphi `$ with $`\varphi =1`$), (iii) finally, it must be a product state, because if it was entangled then its distance from the compact set of product pure states defined as $`\mathrm{min}_{|e,f}|\varphi |e,f`$ would be nonzero, which is obviously impossible. We conclude thus $`|\varphi =|g,hR(\rho )`$ for some $`|g,|h`$, which implies (up to irrelevant phase factors) that $`|g_n|g`$ and $`|h_n|h`$. We have (again up to irrelevant external phase factor) $`|g_n,h_n^{}|g,h^{}`$. The latter must belong to $`R(\rho ^{T_B})`$ as any element of the corresponding sequence is orthogonal to $`K(\rho ^{T_B})`$. $`\mathrm{}`$
Let us now prove the following general lemma, which is a generalization of one theorem from Ref. :
Lemma B2: Let the PPTES $`\rho `$ be defined on a finite dimensional Hilbert space. Consider the set $`\mathrm{\Sigma }_\rho `$ consisting of the trivial zero operator plus all unnormalized states $`\stackrel{~}{\rho }`$ ($`\mathrm{tr}\stackrel{~}{\rho }1`$) such that $`\stackrel{~}{\delta }\rho \stackrel{~}{\rho }`$ is positive and has positive partial transpose. Then, one can find $`\widehat{\rho }\mathrm{\Sigma }_\rho `$ such that with $`\mathrm{tr}(\widehat{\rho })1`$ is optimal in the sense that:
The trace of $`\widehat{\delta }\rho \widehat{\rho }`$ is minimal with respect to all separable $`\stackrel{~}{\rho }`$’s leading to positive partial transpose $`\stackrel{~}{\delta }`$’s.
The state $`\delta =\widehat{\delta }/\mathrm{tr}(\widehat{\delta })`$ is an “edge” PPTES.
Proof : To prove the existence of $`\widehat{\rho }\mathrm{\Sigma }_\rho `$ we just have to show that $`\mathrm{\Sigma }_\rho `$ is compact. This can be done by showing that $`\mathrm{\Sigma }_\rho `$ is a closed subset of another compact set, namely $`C=\mathrm{conv}\{\mathrm{\Gamma }_\rho `$ 0 $`\}`$. The latter set $`C`$ is compact as it is a convex hull of the compact set $`\{\mathrm{\Gamma }_\rho `$ 0 $`\}`$ in a finite dimensional space.
Note first that $`\mathrm{\Sigma }_\rho C`$. Indeed, by virtue of $`\stackrel{~}{\delta }0`$ any nonzero $`\stackrel{~}{\rho }`$ cannot have any vector in its range not belonging to $`R(\rho )`$. Analogously $`R(\stackrel{~}{\rho }^{T_B})R(\rho ^{T_B})`$. Hence, according to the properties of the ranges of density operators in general , $`\stackrel{~}{\rho }`$ must be a convex combination of vectors from $`\mathrm{\Gamma }_\rho `$, and as such it belongs to $`C`$. Let us show that $`\mathrm{\Sigma }_\rho `$ is closed. This follows immediately form the fact that $`\mathrm{\Sigma }_\rho `$ is a cross-section (performed over any projections $`P`$, $`Q`$) of the sets: $`\mathrm{\Sigma }_{\rho ,P}^1\{\stackrel{~}{\rho }:f_{P,\rho }(\stackrel{~}{\rho })\mathrm{tr}(P\rho P\stackrel{~}{\rho })0\}`$ and $`\mathrm{\Sigma }_{\rho ,Q}^2=\{\stackrel{~}{\rho }:g_{Q,\rho }(\stackrel{~}{\rho })\mathrm{tr}(Q^{T_B}\rho Q^{T_B}\stackrel{~}{\rho })0\}`$. Since the functions $`f_{P,\rho },g_{Q,\rho }`$ are continuous, all the sets participating in the cross section are closed. Now, the cross-section of closed sets is again a closed one.
Consider now the statement (ii). Since $`\delta ,\delta ^{T_B}0`$, we always have $`\delta =\beta P_{R(\delta )}+A`$ and $`\delta ^{T_B}=\beta ^{}P_{R(\delta ^{T_B})}+A^{}`$ with $`\beta ,\beta ^{}>0`$, some positive operators $`A,A^{}`$ (here, $`P_X`$ denotes a projector onto the subspace $`X`$). Then if, contrary to (ii), there were any $`|e,fR(\delta )`$ such that $`|e,f^{}R(\delta ^{T_B})`$, then the new operator $`\widehat{\rho }^{}=\widehat{\rho }+\gamma |e,fe,f|`$, $`\gamma =\mathrm{min}[\beta ,\beta ^{}]`$ would fulfill that $`\widehat{\delta }^{}=\rho \widehat{\rho }^{}`$ is a PPTES, and would contradict optimality with respect to (i). $`\mathrm{}`$
Let us remark that if we give up the condition regarding positivity of $`\stackrel{~}{\delta }^{T_B}`$, then we obtain a modified statement (ii) where the state $`\delta `$ has no product vectors in its range. This is nothing but the best separable approximation (BSA) of Ref. , extended here rigorously to the states $`\rho `$ having uncountable set of product vectors in $`R(\rho )`$.
From the Lemma B2 we obtain the following characterization of PPTES, which can be regarded to be among the main results of this appendix, since it provides a canonical form of PPTES:
Proposition : If the state $`\rho `$ is PPTES, then it is a convex combination
$$\rho =(1p)\rho _{sep}+p\delta $$
(B1)
of some normalized separable $`\rho _{sep}`$ and a normalized “edge” PPTES $`\delta `$. In the above decomposition the weight $`p`$ is minimal \[i. e. there does not exist a decomposition of type (B1) with a smaller $`p`$\].
The above proposition means, in particular, that the edge PPTES are responsible for PPT type entanglement.
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# Footprints of a Broad 𝜎(600) in Weak-Interaction Processes
## Appendix: $`\eta `$\- and $`\eta ^{}`$-Pole Contributions to $`K_L\gamma \gamma `$
If the $`K_L\gamma \gamma `$ amplitude of Fig. 3a is augmented by contributions from $`\eta `$ and $`\eta ^{}`$ poles, the $`\pi ^{}`$, $`\eta `$, and $`\eta ^{}`$ contributions to the amplitude are respectively given by
$$M_\pi ^{}=<2\gamma |\pi ^{}>\frac{1}{m_{K_L}^2m_\pi ^{}^2}<\pi ^{}|H_w^{pc}|K_L>,$$
(A.1)
$$M_\eta =<2\gamma |\eta >\frac{1}{m_{K_L}^2m_\eta ^2}<\eta |H_w^{pc}|K_L>,$$
(A.2)
$$M_\eta ^{}=<2\gamma |\eta ^{}>\frac{1}{m_{K_L}^2m_\eta ^{}^2}<\eta ^{}|H_w^{pc}|K_L>.$$
(A.3)
To find the relative contributions of these matrix elements, we first note that
$$<\eta |H_w|K_L>=cos\theta _p<\eta _8|H_w|K_L>sin\theta _p<\eta _0|H_w|K_L>,$$
(A.4)
$$<\eta ^{}|H_w|K_L>=sin\theta _p<\eta _8|H_w|K_L>+cos\theta _p<\eta _0|H_w|K_L>,$$
(A.5)
where the pseudoscalar mixing angle $`\theta _p=12.9^{}`$ . If the relative sizes of transitions from $`K_L`$ to nonstrange pseudoscalar-nonet states is scaled to the $`U(3)`$ structure constants \[i.e., $`(<\pi ^{}|H_w^{pc}|K_L>:`$ $`<\eta _8|H_w^{pc}|K_L>:`$ $`<\eta _0|H_w^{pc}|K_L>)=(d_{366}:d_{866}:d_{066})=(\frac{1}{2}:\frac{1}{2\sqrt{3}}:\sqrt{\frac{2}{3}})`$\], we then find that
$$<\eta |H_w|K_L>=0.198<\pi ^{}|H_w|K_L>,$$
(A.6)
$$<\eta ^{}|H_w|K_L>=1.72<\pi ^{}|H_w|K_L>.$$
(A.7)
Using the matrix elements
$`<2\gamma |\pi ^{}>=0.0250GeV^1,`$ (A.8)
$`<2\gamma |\eta >=0.0255GeV^1,`$ (A.9)
$`<2\gamma |\eta ^{}>=0.0335GeV^1`$ (A.10)
\[Levi-Civita covariants have been factored out of (A.8)\], we then find from (A.1-3) that the matrix elements for $`\pi ^{}`$, $`\eta `$, and $`\eta ^{}`$ pole contributions to $`K_L2\gamma `$ are respectively given by
$$M_\pi ^{}=(0.109GeV^3)<\pi ^{}|H_w^{pc}|K_L>,$$
(A.11)
$$M_\eta =(0.0975GeV^3)<\pi ^{}|H_w^{pc}|K_L>,$$
(A.12)
$$M_\eta ^{}=(+0.0861GeV^3)<\pi ^{}|H_w^{pc}|K_L>.$$
(A.13)
Consequently, there is a near cancellation of $`\eta `$ and $`\eta ^{}`$ pole contributions in the matrix-element sum:
$$M_\pi ^{}+M_\eta +M_\eta ^{}=(0.0976GeV^3)<\pi ^{}|H_w^{pc}|K_L>=(0.90)M_\pi ^{}.$$
(A.14)
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# QCD at Finite Isospin Density
\[
## Abstract
QCD at finite isospin chemical potential $`\mu _I`$ has no fermion sign problem and can be studied on the lattice. We solve this theory analytically in two limits: at low $`\mu _I`$, where chiral perturbation theory is applicable, and at asymptotically high $`\mu _I`$, where perturbative QCD works. At low isospin density the ground state is a pion condensate, whereas at high density it is a Fermi liquid with Cooper pairing. The pairs carry the same quantum numbers as the pion. This leads us to conjecture that the transition from hadron to quark matter is smooth, which passes several tests. Our results imply a nontrivial phase diagram in the space of temperature and chemical potentials of isospin and baryon number.
preprint: CU-TP-xxx
\]
Introduction.—Ample knowledge of QCD in the regime of finite temperature and baryon density is crucial for understanding a wide range of phenomena from heavy ion collisions to neutron stars and cosmology. First-principles lattice numerical Monte Carlo calculations provide a solid basis for our knowledge of the finite-temperature regime. However, the regime of finite baryon chemical potential $`\mu _B`$ is still inaccessible by Monte Carlo because present methods of evaluating the QCD partition function require taking a path integral with a measure which includes a complex fermion determinant. Ignoring the determinant (as in the popular quenched approximation) leads to qualitatively wrong answers for finite $`\mu _B`$ . Such a contrast to the case of $`\mu _B=0`$, where the quenched approximation proved useful, comes from the fact that the latter corresponds to an unphysical theory with pairs of quarks of opposite baryon charges (conjugate quarks) . This is one of the main reasons why our understanding of QCD at finite baryon density is still rudimentary. Many interesting phenomena, such as color superconductivity and color-flavor locking , occur at finite baryon density, beyond the reach of current lattice techniques.
To understand the regime of finite baryon density one would need to follow the transition from hadronic to quark degrees of freedom by increasing the density of a conserved charge (such as baryon number), i.e., without invoking the temperature. This is the motivation for us to turn to QCD at finite chemical potential $`\mu _I`$ of isospin (more precisely, of the third component, $`I_3`$). Nature provides us with nonzero $`\mu _I`$ systems in the form of isospin-asymmetric matter. These always contain both isospin density and baryon density. In any realistic setting $`\mu _I\mu _B`$. In this paper, however, we shall consider an idealization in which $`\mu _I`$ is nonzero while $`\mu _B=0`$. Such a system is unstable with respect to weak decays which do not conserve isospin. However, since we are interested in the dynamics of strong interaction alone, one can imagine that all relatively slow electroweak effects are turned off. Once this is done, we have a nontrivial regime which, as has been emphasized recently in , is accessible by present lattice Monte Carlo methods, while being, as we shall see, analytically tractable in various interesting limits. As a result, the system we consider has a potential to improve substantially our understanding of cold dense QCD. This regime carries many attractive traits of two-color QCD , but is realized in a physically relevant theory — QCD with three colors.
Positivity and QCD inequalities.—Since the Euclidean version of our theory has a real and positive fermion determinant, some rigorous results on the low-energy behavior can be obtained from QCD inequalities . In vacuum QCD, the latter rely on the following property of the Euclidean Dirac operator $`𝒟=\gamma (+iA)+m`$:
$$\gamma _5𝒟\gamma _5=𝒟^{}.$$
(1)
which, in particular, implies positivity $`det𝒟0`$. For the correlator of a generic meson $`M=\overline{\psi }\mathrm{\Gamma }\psi `$, we can write, using (1) and the Schwartz inequality:
$`M(x)M^{}(0)_{\psi ,A}=\mathrm{Tr}𝒮(x,0)\mathrm{\Gamma }𝒮(0,x)\overline{\mathrm{\Gamma }}_A=`$ (2)
$`\mathrm{Tr}𝒮(x,0)\mathrm{\Gamma }i\gamma _5𝒮^{}(x,0)i\gamma _5\overline{\mathrm{\Gamma }}_A\mathrm{Tr}𝒮(x,0)𝒮^{}(x,0)_A,`$ (3)
where $`𝒮𝒟^1`$ and $`\overline{\mathrm{\Gamma }}\gamma _0\mathrm{\Gamma }^{}\gamma _0`$. The inequality is saturated for mesons with $`\mathrm{\Gamma }=i\gamma _5\tau _i`$, since $`𝒟`$ commutes with isospin $`\tau _i`$, which means that the pseudoscalar correlators are larger, point-by-point, than all other $`I=1`$ meson correlators . As a consequence, one obtains an important restriction on the pattern of the symmetry breaking: for example, it cannot be driven by a condensate of $`\overline{\psi }\gamma _5\psi `$, which would give $`0^+`$ Goldstones.
At finite isospin density, $`\mu _I0`$, positivity still holds and certain inequalities can be derived (in contrast with the case of $`\mu _B0`$ when there is no positivity). Now $`𝒟=\gamma (+iA)+\frac{1}{2}\mu _I\gamma _0\tau _3+m`$, and Eq. (1) is not true anymore, since the operation on the right-hand side of (1) changes the relative sign of $`\mu _I`$. However, provided $`m_u=m_d`$, interchanging up and down quarks compensates for this sign change (the $`u`$ and $`d`$ quarks play the role of mutually conjugate quarks ), i.e,
$$\tau _1\gamma _5𝒟\gamma _5\tau _1=𝒟^{}.$$
(4)
Instead of isospin $`\tau _1`$ in (4) one can also use $`\tau _2`$ (but not $`\tau _3`$). Equation (4), in place of the now invalid Eq. (1), ensures that $`det𝒟0`$. Repeating the derivation of the QCD inequalities, by using (4) we find that the lightest meson, or the condensate, must be in channels $`\overline{\psi }i\gamma _5\tau _{1,2}\psi `$, i.e., a linear combination of $`\pi ^{}\overline{u}\gamma _5d`$ and $`\pi ^+\overline{d}\gamma _5u`$ states. Indeed, as shown below, in the two analytically tractable regimes of small and large $`\mu _I`$ the lightest mode is a massless Goldstone which is a linear combination of $`\overline{u}\gamma _5d`$ and $`\overline{d}\gamma _5u`$.
Small isospin densities.—When $`\mu _I`$ is small compared to the chiral scale (taken here to be $`m_\rho `$), we can use chiral perturbation theory. For zero quark mass and zero $`\mu _I`$ the pions are massless Goldstones of the spontaneously broken SU(2)$`{}_{L}{}^{}\times `$SU(2)<sub>R</sub> chiral symmetry. If the quarks have small equal masses, the symmetry is only SU(2)<sub>L+R</sub>. The low-energy dynamics is governed by the familiar chiral Lagrangian for the pion field $`\mathrm{\Sigma }`$ SU(2): $`=\frac{1}{4}f_\pi ^2\mathrm{Tr}[_\mu \mathrm{\Sigma }_\mu \mathrm{\Sigma }^{}2m_\pi ^2\mathrm{Re}\mathrm{\Sigma }]`$, which contains the pion decay constant $`f_\pi `$ and vacuum pion mass $`m_\pi `$ as phenomenological parameters. The isospin chemical potential further breaks SU(2)<sub>L+R</sub> down to U(1)<sub>L+R</sub>. Its effect can be included to leading order in $`\mu _I`$ without additional phenomenological parameters by promoting SU(2)$`{}_{L}{}^{}\times `$SU(2)<sub>R</sub> to a local gauge symmetry and viewing $`\mu _I`$ as the zeroth component of a gauge potential . Gauge invariance thus fixes the way $`\mu _I`$ enters the chiral Lagrangian:
$$_{\mathrm{eff}}=\frac{f_\pi ^2}{4}\mathrm{Tr}_\nu \mathrm{\Sigma }_\nu \mathrm{\Sigma }^{}\frac{m_\pi ^2f_\pi ^2}{2}\mathrm{ReTr}\mathrm{\Sigma },$$
(5)
where the covariant derivative is defined as
$$_0\mathrm{\Sigma }=_0\mathrm{\Sigma }\frac{\mu _I}{2}(\tau _3\mathrm{\Sigma }\mathrm{\Sigma }\tau _3).$$
(6)
By using (5), it is straightforward to determine vacuum alignment of $`\mathrm{\Sigma }`$ as a function of $`\mu _I`$ and the spectrum of excitations around the vacuum. We are interested in negative $`\mu _I`$, which favors neutrons over protons, as in neutron stars. The results are very similar to two-color QCD at finite baryon density :
(i) For $`|\mu _I|<m_\pi `$, the system is in the same ground state as at $`\mu _I=0`$: $`\overline{\mathrm{\Sigma }}=1`$. This is because the lowest lying pion state costs a positive energy $`m_\pi |\mu _I|`$ to excite, which is impossible at zero temperature.
(ii) When $`|\mu _I|`$ exceeds $`m_\pi `$ it is favorable to excite $`\pi ^{}`$ quanta, which form a Bose condensate. In the language of the effective theory, such a pion condensate is described by a tilt of the chiral condensate $`\overline{\mathrm{\Sigma }}`$,
$`\overline{\mathrm{\Sigma }}`$ $`=`$ $`\mathrm{cos}\alpha +i(\tau _1\mathrm{cos}\varphi +\tau _2\mathrm{sin}\varphi )\mathrm{sin}\alpha ,`$ (8)
$`\mathrm{cos}\alpha =m_\pi ^2/\mu _I^2.`$
The tilt angle $`\alpha `$ is determined by minimizing the vacuum energy. The energy is degenerate with respect to the angle $`\varphi `$, corresponding to the spontaneous breaking of the U(1)<sub>L+R</sub> symmetry generated by $`I_3`$ in the Lagrangian (5). The ground state is a pion superfluid, with one massless Goldstone mode. Since we start from a theory with three pions, there are two massive modes which can be identified with $`\pi _0`$ and a linear combination of $`\pi ^+`$ and $`\pi ^{}`$. At the condensation threshold, $`m_{\pi _0}=m_\pi `$ and the mass of the other mode is $`2m_\pi `$, while for $`|\mu _I|m_\pi `$ both masses approach $`|\mu _I|`$.
The isospin density is found by differentiating the ground state energy with respect to $`\mu _I`$ and is equal to:
$$n_I=f_\pi ^2\mu _I\mathrm{sin}^2\alpha =f_\pi ^2\mu _I\left(1\frac{m_\pi ^4}{\mu _I^4}\right),|\mu _I|>m_\pi .$$
(9)
For $`|\mu _I|`$ just above the condensation threshold, $`|\mu _I|m_\pi m_\pi `$, Eq. (9) reproduces the equation of state of the dilute nonrelativistic pion gas .
It is also possible to find baryon masses, i.e., the energy cost of introducing a single baryon into the system. The most interesting baryons are those with lowest energy and highest isospin, i.e. neutron $`n`$ and $`\mathrm{\Delta }^{}`$ isobar. There are two effects of $`\mu _I`$ on the baryon masses. The first comes from the isospin of the baryons, which effectively reduces the neutron mass by $`\frac{1}{2}|\mu _I|`$ and the $`\mathrm{\Delta }^{}`$ mass by $`\frac{3}{2}|\mu _I|`$. Alone, this effect would lead to the formation of baryon/antibaryon Fermi surfaces, manifested in nonvanishing zero-temperature baryon susceptibility $`\chi _Bn_B/\mu _B`$ when $`\mu _I>\frac{2}{3}m_\mathrm{\Delta }`$. However, long before that, another effect turns on: the $`\pi ^{}`$’s in the condensate tend to repel the baryons, lifting up their masses. These effects can be treated in the framework of the baryon chiral perturbation theory , giving
$$m_n=m_N\frac{|\mu _I|}{2}\mathrm{cos}\alpha ,m_\mathrm{\Delta }^{}=m_\mathrm{\Delta }\frac{3|\mu _I|}{2}\mathrm{cos}\alpha $$
(10)
in the approximation of nonrelativistic baryons. Equation (10) can be interpreted as follows: as a result of the rotation (8) of the chiral condensate, the nucleon mass eigenstate becomes a superposition of vacuum $`n`$ and $`p`$ states. The expectation value of the isospin in this state is proportional to $`\mathrm{cos}\alpha `$ appearing in (10). With $`\mathrm{cos}\alpha `$ given in Eq.(8), we see that the two mentioned effects cancel each other when $`m_\pi |\mu _I|m_\rho `$. Thus the baryon mass never drops to zero, and $`\chi _B=0`$ at zero temperature in the region of applicability of the chiral Lagrangian.
As one forces more pions into the condensate, the pions are packed closer and their interaction becomes stronger. When $`\mu _Im_\rho `$, the chiral perturbation theory breaks down. To find the equation of state in this regime, full QCD has to be employed. As we have seen, this can be done using present lattice techniques since the fermion sign problem is not present at finite $`\mu _I`$, similar to the two-color QCD .
Asymptotically high isospin densities.—In the opposite limit of very large isospin densities, or $`|\mu _I|m_\rho `$, the description in terms of quark degrees of freedom applies since the latter are weakly interacting due to asymptotic freedom. In our case of large negative $`\mu _I`$, or $`n_I`$, the ground state consists of $`d`$ quarks and $`\overline{u}`$ antiquarks which, neglecting the interaction, fill two Fermi spheres with equal radii $`|\mu _I|/2`$. Turning on the interaction between the fermions leads to the instability with respect to the formation and condensation of Cooper pairs, similar, to some extent, to the diquark pairing at high baryon density . In our case, $`\mu _I<0`$, the Cooper pair consists of a $`\overline{u}`$ and a $`d`$ in the color singlet channel. The order parameter has the same quantum numbers as the pion condensate at lower densities,
$$\overline{u}\gamma ^5d0.$$
(11)
Because of Cooper pairing, the fermion spectrum acquires a gap $`\mathrm{\Delta }`$ at the Fermi surface, where
$$\mathrm{\Delta }=b|\mu _I|g^5e^{c/g},c=3\pi ^2/2$$
(12)
where $`g`$ should be evaluated at the scale $`|\mu _I|`$. This behavior comes from the long-range magnetic interaction, as in the superconducting gap at large $`\mu _B`$ . The constant $`c`$ is smaller by a factor of $`\sqrt{2}`$ compared to the latter case due to the stronger one-gluon attraction in the singlet $`q\overline{q}`$ channel compared to the $`\overline{\mathrm{𝟑}}`$ diquark channel. Consequently, the gap (12) is exponentially larger than the diquark gap at comparable baryon chemical potentials. By using the methods of , one can estimate $`b10^4`$.
The perturbative one-gluon exchange responsible for pairing at large $`\mu _I`$ does not distinguish $`\overline{u}d`$ and $`\overline{u}\gamma _5d`$ channels: the attraction is the same in both. The $`\overline{u}\gamma _5d`$ channel is favored by the instanton-induced interactions, which explains the fact that the condensate is a pseudoscalar and breaks parity. This is consistent with our observation that QCD inequalities also constrain the $`I=1`$ condensate to be a pseudoscalar at any $`\mu _I`$.
Quark-hadron continuity.—Since the order parameter (11) has the same quantum numbers and breaks the same symmetry as the pion condensate in the low-density regime, it is plausible that there is no phase transition along the $`\mu _I`$ axis. In this case the Bose condensate of weakly interacting pions smoothly transforms into the superfluid state of $`\overline{u}d`$ Cooper pairs. The situation is very similar to that of strongly coupled superconductors with a “pseudogap” , and possibly of high-temperature superconductors . This also parallels the continuity between nuclear and quark matter in three-flavor QCD as conjectured by Schäfer and Wilczek . We hence conjecture that, in two-flavor QCD, one can move continuously from the hadron phase to the quark phase without encountering a phase transition. Since a first order deconfinement phase transition at intermediate isospin chemical potential cannot be rigorously ruled out (though it is unlikely, see below), this conjecture needs to be verified by lattice calculations.
A number of nontrivial arguments support the continuity hypothesis. One notices that all fermions have a gap at large $`|\mu _I|`$, which implies that $`\chi _B=0`$ at $`T=0`$. This is also true at small $`\mu _I`$. It is thus natural to expect that $`\chi _B`$ remains zero at $`T=0`$ for all $`\mu _I`$, which also suggests one way to check the continuity on the lattice.
Another argument comes from considering the limit of a large number of colors $`N_c`$. In finite-temperature QCD, the fact that the number of gluon degrees of freedom is $`𝒪(N_c^2)`$ while that of hadrons is $`𝒪(N_c^0)`$ hints at a first order confinement-deconfinement phase transition. At very large $`\mu _I`$ thermodynamic quantities such as the density of isospin $`n_I`$ are proportional to $`N_c`$. On the other hand, in the large $`N_c`$ limit the pion decay constant scales as $`f_\pi ^2=𝒪(N_c)`$, and according to Eq. (9) the isospin density in the pion gas is also proportional to $`N_c`$. Physically, the repulsion between pions becomes weaker as one goes to large $`N_c`$, thus more pions are stacked at a given chemical potential. As a result, the $`N_c`$ dependence of thermodynamic quantities is the same in the quark and the hadronic regimes.
Confinement.— At large $`\mu _I`$, gluons softer than $`\mathrm{\Delta }`$ are not screened by the Meissner or by the Debye effect : the condensate does not break gauge symmetry (in contrast to the color superconducting condensate ) and there are no low-lying color excitations to screen the electric field. Thus, the gluon sector below the $`\mathrm{\Delta }`$ scale is described by pure gluodynamics, which is confining. This means there are no quark excitations above the ground state: all particles and holes must be confined in colorless objects, mesons and baryons, just like in vacuum QCD. If there is no transition along the $`\mu _I`$ axis, we expect confinement at all values of $`\mu _I`$. Since the running strong coupling $`\alpha _s`$ at the scale of $`\mathrm{\Delta }`$ is small, the confinement scale $`\mathrm{\Lambda }_{\mathrm{QCD}}^{}`$ (which is, in general, different from $`\mathrm{\Lambda }_{\mathrm{QCD}}`$) is much less than $`\mathrm{\Delta }`$. At large $`\mu _I`$, we thus predict a temperature driven deconfinement phase transition at a temperature $`T_c^{}`$ of order $`\mathrm{\Lambda }_{\mathrm{QCD}}^{}`$, which is expected to be of first order as in pure gluodynamics. Since $`\mathrm{\Lambda }_{\mathrm{QCD}}^{}\mathrm{\Delta }`$ the hadronic spectrum is similar to that of a heavy quarkonium, with $`\mathrm{\Delta }`$ playing the role of the heavy quark mass.
The $`(T,\mu _I)`$ phase diagram.—By considering nonzero $`\mu _I`$, we make the phase diagram of QCD three dimensional: $`(T,\mu _B,\mu _I)`$. Two planes in this three-dimensonal space are of a special interest: the $`\mu _B=0`$ $`(T,\mu _I)`$ plane, which is completely accessible by present lattice techniques, and the $`T=0`$ $`(\mu _I,\mu _B)`$ plane, where the neutron star matter belongs. Two phenomena determine the $`(T,\mu _I)`$ phase plane (Fig.1): pion condensation and confinement.
At sufficiently high temperature the condensate (11) melts (solid line in Fig. 1). For large $`\mu _I`$, this critical temperature is proportional to the BCS gap (12). There are two phases which differ by symmetry: the high temperature phase, where the explicit flavor $`U(1)_{L+R}`$ symmetry is restored, and the low-temperature phase, where this symmetry is spontaneously broken. The phase transition is in the O(2) universality class . The critical temperature $`T_c`$ vanishes at $`\mu _I=m_\pi `$ and is an increasing function of $`\mu _I`$ in both regimes we studied: $`|\mu _I|m_\rho `$ and $`|\mu _I|\mathrm{\Lambda }_{\mathrm{QCD}}`$. Thus, it is likely that $`T_c(\mu _I)`$ is a monotonous function of $`\mu _I`$. In addition, at large $`\mu _I`$, there is a first order deconfinement phase transition at $`T_c^{}`$ much lower than $`T_c(\mu _I)`$. Since there is no phase transition at $`\mu _I=0`$ (for small $`m_{u,d}`$) or at $`T=0`$ (assuming quark-hadron continuity), this first-order line must end at some point $`A`$ on the $`(T,\mu _I)`$ plane (Fig. 1). The exact location of $`A`$ should be determined by lattice calculations; one of the possibilities is drawn in Fig. 1.
The $`(\mu _I,\mu _B)`$ phase diagram.—This phase diagram deserves a separate study. Here we shall only consider the regime $`|\mu _I|\mu _B`$ (the opposite limit $`\mu _B|\mu _I|`$ was considered in Ref. ). Finite $`\mu _B`$ provides a mismatch between $`\overline{u}`$ and $`d`$ Fermi spheres, which makes the superconducting state unfavorable at some value of $`\mu _B`$ of order $`\mathrm{\Delta }`$. It is known that the destruction of this state occurs through two phase transitions: one at $`\mu _B`$ slightly below $`\mathrm{\Delta }/\sqrt{2}`$ and another at $`\mu _B=0.754\mathrm{\Delta }`$. The ground state between the two phase transitions is the Fulde-Ferrell-Larkin-Ovchinnikov (FFLO) state , characterized by a spatially modulated superfluid order parameter $`\overline{u}\gamma _5d`$ with a wavenumber of order $`2\mu _B`$. The FFLO state has the same symmetries as the inhomogeneous pion condensation state which might form in electrically neutral nuclear matter at high densities . It is conceivable that the two phases are actually one, i.e., continuously connected on the $`(\mu _I,\mu _B)`$ phase diagram.
The authors thank L. McLerran, J. Kogut, R. Pisarski, and E. Shuryak for discussions, the DOE Institute for Nuclear Theory at the University of Washington for its hospitality, and K. Rajagopal for drawing their attention to Ref. .
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# 1 Introduction.
## 1 Introduction.
Physical theories formulated in different-than-usual spacetimes signatures have recently found increased attention. One of the reasons can be traced to the conjectured $`F`$-theory which supposedly lives in $`(2+10)`$ dimensions . The current interest in AdS theories motivated by the AdS/CFT correspondence furnishes another motivation. Two-time physics e.g. has started been explored by Bars and collaborators in a series of papers . For another motivation we can also recall that a fundamental theory is expected to explain not only the spacetime dimensionality, but even its signature (see ). Quite recently Hull and Hull-Khuri pointed out the existence of dualities relating different compactifications of theories formulated in different signatures. Such a result provides new insights to the whole question of spacetime signatures. Other papers (e.g.) have remarked the existence of space-time dualities.
Majorana-Weyl spacetimes (i.e. those supporting Majorana-Weyl spinors) are at the very core of the present knowledge of the unification via supersymmetry, being at the basis of ten-dimensional superstrings, superYang-Mills and supergravity theories (and perhaps the already mentioned $`F`$-theory). A well-established feature of Majorana-Weyl spacetimes is that they are endorsed of a rich structure. A legitimate question that could be asked is whether they are affected, and how, by space-time dualities. The answer can be stated as follows, all different Majorana-Weyl spacetimes which are possibly present in any given dimension are each-other related by duality transformations which are induced by the $`Spin(8)`$ triality automorphisms. The action of the triality automorphisms is quite non-trivial and has far richer consequences than the $`𝐙_2`$-duality (its most trivial representative) associated to the space-time $`(s,t)(t,s)`$ exchange discussed in . It corresponds to $`S_3`$, the six-element group of permutations of three letters, identified with the group of congruences of the triangle and generated by two reflections. The lowest-dimension in which the triality action is non-trivial is $`8`$ (not quite a coincidence), where the spacetimes $`(8+0)(4+4)(0+8)`$ are all interrelated. They correspond to the transverse coordinates of the $`(9+1)(5+5)(1+9)`$ spacetimes respectively, where the triality action can also be lifted. Triality relates as well the $`12`$-dimensional Majorana-Weyl spacetimes $`(10+2)(6+6)(2+10)`$, i.e. the potentially interesting cases for the $`F`$-theory. Triality allows explaining the presence of points (read theories) in the brane-molecule table of ref. , corresponding to the different versions of e.g. superstrings, $`11`$-dimensional supermembranes, fivebranes.
As a consequence of triality, supersymmetric theories formulated with Majora-na-Weyl spinors in a given dimension but with different signatures, are all dually mapped one into another. A three-language dictionary can be furnished with the exact translations among, e.g., the different versions of such supersymmetric theories.
It should be stressed the fact that, unlike , the dualities here discussed are already present for the uncompactified theories and in this respect look more fundamental The reason why the triality of the $`d=8`$-dimension plays a role even for Majorana-Weyl spacetimes in $`d>8`$ is in consequence of the representation properties of $`\mathrm{\Gamma }`$-matrices in higher dimensions.
## 2 The set of data for Majorana-Weyl supersymmetric theories.
At first we present the set of data needed to specify a supersymmetric theory involving Majorana-Weyl spinors. The notations here introduced follow and .
The most suitable basis is the Majorana-Weyl basis (MWR), where all spinors are either real or imaginary. In such a representation the following set of data underlines any given theory:
i) the vector fields (or, in the string/brane picture, the bosonic coordinates of the target $`x_m`$), specified by a vector index denoted by $`m`$;
ii) the spinor fields (or, in the string/brane picture, the fermionic coordinates of the target $`\psi _a`$, $`\chi _{\dot{a}}`$), specified by chiral and antichiral indices $`a`$, $`\dot{a}`$ respectively;
iii) the diagonal (pseudo-)orthogonal spacetime metric $`(g^1)^{mn}`$, $`g_{mn}`$ which we will assume to be flat;
iv) the $`𝒜`$ matrix, used to define barred spinors, coinciding with the $`\mathrm{\Gamma }^0`$-matrix in the Minkowski case; in a Majorana-Weyl basis is decomposed in an equal-size block diagonal form such as $`𝒜=A\stackrel{~}{A}`$, with structure of indices $`(A)_{a}^{}{}_{}{}^{b}`$ and $`(\stackrel{~}{A})_{\dot{a}}^{}{}_{}{}^{\dot{b}}`$ respectively;
v) the charge-conjugation matrix $`𝒞`$ which also appears in an equal-size block diagonal form $`𝒞=C^1\stackrel{~}{C}^1`$. It is invariant under bispinorial transformations and it can be promoted to be a metric in the space of chiral (and respectively antichiral) spinors, used to raise and lower spinorial indices. Indeed we can set $`(C^1)^{ab}`$, $`(C)_{ab}`$, and $`(\stackrel{~}{C}^1)^{\dot{a}\dot{b}}`$, $`(\stackrel{~}{C})_{\dot{a}\dot{b}}`$;
vi) the $`\mathrm{\Gamma }`$-matrices, which are decomposed in equal-size blocks, $`\sigma ^m`$’s the upper-right blocks and $`\stackrel{~}{\sigma }^m`$’s the lower-left blocks having structure of indices $`(\sigma ^m)_{a}^{}{}_{}{}^{\dot{b}}`$ and $`(\stackrel{~}{\sigma }^m)_{\dot{a}}^{}{}_{}{}^{b}`$ respectively;
vii) the $`\eta =\pm 1`$ sign, labeling the two inequivalent choices for $`𝒞`$.
The above structures are common in any theory involving Majorana-Weyl spinors. An explicit dictionary relating Majorana-Weyl spacetimes having the same dimensionality but different signature is presented in . The structures i)-vii) are related via triality transformations which close the $`S_3`$ permutation group. They constitute the “words” in the three-language dictionary.
Majorana-Weyl spacetimes exist for different signatures of a given dimension $`d`$ if $`d8`$. The special $`d=8`$ case is the fundamental one. Indeed, we are able to express higher-dimensional $`\mathrm{\Gamma }`$ matrices (and in consequence all the above-mentioned structures which define a Majorana-Weyl theory) in terms of lower-dimensional ones according to the recursive formula
$`\mathrm{\Gamma }_{d}^{}{}_{}{}^{i=1,\mathrm{},s+1}`$ $`=`$ $`\sigma _x\mathrm{𝟏}_L\mathrm{\Gamma }_{s}^{}{}_{}{}^{i=1,\mathrm{},s+1}`$
$`\mathrm{\Gamma }_{d}^{}{}_{}{}^{s+1+j=s+2,\mathrm{},d}`$ $`=`$ $`\sigma _y\mathrm{\Gamma }_{r}^{}{}_{}{}^{j=1,\mathrm{},r+1}\mathrm{𝟏}_R`$ (1)
where $`\mathrm{𝟏}_{L,R}`$ are the unit-matrices in the respective spaces, while $`\sigma _x=e_{12}+e_{21}`$ and $`\sigma _y=ie_{12}+ie_{21}`$ are the $`2`$-dimensional Pauli matrices. $`\mathrm{\Gamma }_{r}^{}{}_{}{}^{r+1}`$ corresponds to the “generalized $`\mathrm{\Gamma }^5`$-matrix” in $`r+1`$ dimensions. In the above formula the values $`r,s=0`$ are allowed. The corresponding $`\mathrm{\Gamma }_{0}^{}{}_{}{}^{1}`$ is just $`1`$.
With the help of this formula we are able to reduce the analysis of different-signatures Majorana-Weyl spacetimes to the $`8`$-dimensional case. In this particular dimension the three indices, vector ($`m`$), chiral ($`a`$) and antichiral ($`\dot{a}`$) take values $`m,a,\dot{a}\{1,\mathrm{},8\}`$.
The three Majorana-Weyl solutions, for signatures $`(4+4)`$, $`(8+0)`$, $`(0+8)`$ find a representation in a Majorana-Weyl basis with definite (anti-)symmetry property of the $`\mathrm{\Gamma }`$ matrices. In particular for the $`(4+4)`$-signature the $`(4_S+4_A)`$-representation (see ) of the $`\mathrm{\Gamma }`$-matrices has to be employed for both values of $`\eta `$ in order to provide a Majorana-Weyl basis. In the ($`t=8`$, $`s=0`$) signature the $`(8_S+0_A)`$-representation offers a MW basis for $`\eta =+1`$, while the $`(0_S+8_A)`$ offers it for $`\eta =1`$. The converse is true in the ($`t=0`$, $`s=8`$)-signature.
## 3 Trialities.
The $`S_3`$ outer automorphisms of the $`D_4`$ Lie algebra is responsible for the triality property among the $`8`$-dimensional vector, chiral and antichiral spinor representations of $`SO(8)`$ and $`SO(4,4)`$ which has been first discussed by Cartan . However, besides such Cartan’s V-C-A triality, other triality related properties follow as a consequence. For purpose of clarity it will be convenient to represent them symbolically with triangle diagrams.
A first extra-consequence of triality appears at the level of Majorana-Weyl type of representations for Clifford’s $`\mathrm{\Gamma }`$-matrices. Such representations can be defined as those where all $`\mathrm{\Gamma }`$’s matrices exhibit a well-defined (anti-)symmetry property. In dimension $`d=8`$ three different representations of this kind exist (they have been mentioned in the previous section). Such different eight-dimensional representations can be accomodated into the triangle diagram
$`\begin{array}{ccccc}& & (4_S+4_A)& & \\ & & & & \\ (8_S+0_A)& & & & (0_S+8_A)\end{array}`$ (5)
exhibiting the triality operating at the level of $`\mathrm{\Gamma }`$-matrices. The $`S_3`$ transformations relating the three above representations are realized by similarity transformations. They depend on the concrete choice of the $`\mathrm{\Gamma }`$-matrix representatives and will not been furnished here (see however ).
We have already recalled that such MW-representations are associated with the space-time signature, due to the fact that the introduction of a Majorana-Weyl basis for spinors requires the use of the corresponding Majorana-Weyl representation for Clifford’s $`\mathrm{\Gamma }`$ matrices. As a consequence the triality can be lifted to operate on the whole set of data introduced in the previous section; it can therefore be regarded as operating on the different space-times signatures which support Majorana-Weyl spinors in a given dimensionality, according to the triangles
$`\left(\begin{array}{ccccc}& & 5+5& & \\ & & & & \\ 9+1& & & & 1+9\end{array}\right)`$ $``$ $`\left(\begin{array}{ccccc}& & 4+4& & \\ & & & & \\ 8+0& & & & 0+8\end{array}\right)`$ (12)
The arrow has been inserted to recall that such triality can be lifted to higher dimensions or, conversely, that the $`8`$-dimensional spacetimes arise as transverse coordinates spaces in some physical theories (the most natural example is the $`10`$-dimensional superstring).
The Cartan’s V-C-A triality, schematically represented as
$`\left(\begin{array}{ccccc}& & V& & \\ & & & & \\ C& & & & A\end{array}\right)`$ (17)
and the signature triality can also be combined and symbolically represented by a sort of fractal-like double-triality diagram as follows
$`\begin{array}{ccccc}& & \begin{array}{ccccc}& & V& & \\ & & \mathrm{𝟒}+\mathrm{𝟒}& & \\ C& & & & A\end{array}& & \\ & & & & \\ & & & & \\ & & & & \\ & & & & \\ & & & & \\ \begin{array}{ccccc}& & V& & \\ & & \mathrm{𝟖}+\mathrm{𝟎}& & \\ C& & & & A\end{array}& & & & \begin{array}{ccccc}& & V& & \\ & & \mathrm{𝟎}+\mathrm{𝟖}& & \\ C& & & & A\end{array}\end{array}`$ (34)
The bigger triangle illustrates the signature triality, while the smaller triangles visualize the trialities for vectors, chiral and antichiral spinors which can be accomodated in each space-time.
It is worth stressing the fact that the arising of the $`S_3`$ permutation group as a signature-duality group for Majorana-Weyl spacetimes in a given dimension is not a completely straightforward consequence of the existence of Majorana-Weyl spacetimes in three different signatures. Some extra-requirements have to be fulfilled in order to reach this result. As an example we just mention that a necessary condition for the presence of $`S_3`$ requires that each given couple of the three different spacetimes must differ by an even number of signatures (in this point is discussed in full detail); the flipping of an odd number of signatures, like the Wick rotation from Minkowski to the Euclidean space, cannot be achieved with a $`𝐙_2`$ group when spinors are involved. An example is provided by the fact that the change of signature e.g. from $`(++)()`$ can be realized on $`\mathrm{\Gamma }`$-matrices through similarity trasformations expressed in terms of the $`\sigma _y`$ Pauli matrix $`\sigma _y=ie_{12}+ie_{21}`$, through
$`\sigma _y\mathrm{𝟏}_2\sigma _{y}^{}{}_{}{}^{T}`$ $`=`$ $`\mathrm{𝟏}_2`$ (36)
Of course $`\sigma _y`$ satisfies $`\sigma _{y}^{}{}_{}{}^{2}=\mathrm{𝟏}_2`$ and therefore it closes a $`𝐙_2`$ group. On the contrary, a standard Wick rotations from the Minkowski to the Euclidean space leads to a $`𝐙_4`$ group when represented on $`\mathrm{\Gamma }`$ matrices.
Similarity transformations realized by $`\sigma _y`$ Pauli matrices are among the building blocks for constructing the $`S_3`$ duality transformations for different signature Majorana-Weyl spacetimes. The formulas will not be reproduced here (they are furnished in , together with the details of the construction).
Let us conclude this section by mentioning that triality can be seen not only as a source of duality-mappings, but as an invariance property. In the original Cartan’s formulation this is seen as follows. At first a group $`𝒢`$ of invariance is introduced as the group of linear homogeneous transformations acting on the $`8\times 3=24`$ dimensional space leaving invariant, separately, the bilinears $`_V`$, $`_C`$, $`_A`$ for vectors, chiral and antichiral spinors respctively (the spinors are assumed commuting in this case) plus a trilinear term $`𝒯`$. Next, the triality group $`𝒢_{Tr}`$ is defined by relaxing one condition, as the group of linear homogeneous transformations leaving invariant $`𝒯`$ and the total bilinear $`_{Sum}`$,
$`_{Sum}`$ $`=`$ $`_V+_C+_A`$ (37)
It can be proven that $`𝒢_{Tr}`$ is given by the semidirect product of $`𝒢`$ and $`S_3`$:
$`𝒢_{Tr}`$ $`=`$ $`𝒢_SS_3`$
This feature can be extended to the other aspects of triality. It follows the possibility to look at formulations of higher dimensional supersymmetric theories presenting an $`S_3`$ group of symmetry under the exchange of space and time coordinates.
It should be mentioned that the higher-dimensional supersymmetry strongly restricts the class of finite groups which can provide “unification between space and time” or, otherwise stated, symmetry under time-versus-space coordinates exchange. In the bosonic case such class of groups is quite large, while if we consider e.g. the $`10`$-dimensional supersymmetric case only three possibilities are left, namely i) the identity $`\mathrm{𝟏}`$, corresponding to a theory formulated in the single spacetime $`(5,5)`$, ii) the $`𝐙_2`$ group for a theory which is formulated by using two spacetimes copies $`(1,9)`$ and $`(9,1)`$, iii) the $`S_3`$ group; whose corresponding “space-time unified” theory requires the introduction of the whole set of three $`10`$-dimensional Majorana-Weyl spacetimes $`(1,9)`$, $`(5,5)`$, $`(9,1)`$.
## 4 Conclusions.
In this paper we have shown that the triality automorphisms of $`Spin(8)`$, besides its consequences on the representation properties of the $`8`$-dimensional vectors, chiral and antichiral spinors (the usual Cartan’s notion of triality), can be realized on classes of $`\mathrm{\Gamma }`$-matrices representations which furnish a Majorana-Weyl basis for Majorana-Weyl spinors. Next, triality transformations can be lifted to connect spacetimes supporting Majorana-Weyl spinors sharing the same dimensionality, but different signatures. Recursive formulas for $`\mathrm{\Gamma }`$-matrix representations allow to extend the $`8`$-dimensional properties to higher-dimensional cases as well. Dualities induced by triality are found connecting even-dimensional Majorana-Weyl spacetimes (and odd-dimensional Majorana ones). The presence of different formulations of e.g. brane theories, as shown in the brane-scan molecule table of ref. arises as a consequence.
Indeed higher dimensional supersymmetric theories admits formulations in different signatures which are all interrelated by triality induced transformations.
Besides this action of triality as a source of duality mappings between different versions of supersymmetric theories, triality can provide a setting to discuss formulation of theories invariant under space-versus-time coordinates exchange. This would amount to investigate the formulation of supersymmetric theories exhibiting a manifest $`S_3`$-invariance under signature-triality transformations.
The range of possible applications for the methods and the ideas here discussed is vast. Let us just mention that are currently investigated the web of dualities connecting the six different versions of the $`12`$-dimensional Majorana-Weyl spacetimes which should support the $`F`$-theory (the number $`6=3\times 2`$ is due to the three different signatures of Majorana-Weyl spacetimes and the two values of the $`\eta `$ sign), with the $`6`$ versions of the $`11`$-dimensional Majorana spacetimes (for the $`M`$-theory) in $`(10+1)`$, $`(9+2)`$, $`(6+5)`$, $`(5+6)`$, $`(2+9)`$, $`(1+10)`$ signatures and with the different (again $`3\times 2`$) versions of the $`10`$-dimensional Majorana-Weyl spacetimes.
Acknowledgments
The talk here presented is mainly based on a work with M. A. De Andrade and M. Rojas, who I am very pleased to acknowledge for the fruitful collaboration.
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# Optical and evaporative cooling of cesium atoms in the gravito-optical surface trap
## I Introduction
Optical dipole traps based on far-detuned laser light have become very popular as versatile tools for the storage of ultracold atomic gases and can be employed for a great variety of experiments, e.g., on quantum phenomena, precision measurements, ultracold collisions and quantum gases gri00 . Optical dipole traps have opened up fascinating new experimental possibilities not offered by other trapping methods.
A very important advantage compared with magnetic traps is the fact that a dipole trap can store atoms in any sub-state or mixtures of sub-states of the electronic ground state. Because of the unusual scattering properties of cesium the use of an optical dipole trap may be the only way to achieve Bose-Einstein condensation of this interesting atomic species. For Cs the quantum-mechanical scattering is resonantly enhanced arn97 ; hop00 , and binary collisions flipping the spin state lead to severe loss from magnetic traps soe98 ; arl98 ; gue98a ; gue98b . The latter effect, which has been explained by resonant scattering in combination with a second-order spin-orbit coupling kok98 ; leo98 , has so far prevented the attainment of a quantum-degenerate gas of cesium atoms.
Inelastic two-body collisions are energetically suppressed in the absolute internal ground state of cesium, which is the high-field seeking state $`F=3,m_F=3`$. Atoms in this state cannot be trapped magnetically but in an optical trap. In recent experiments, Vuletic et al. have discovered a low-field Feshbach resonance vul99a , which promises an easy experimental way to tune the $`s`$-wave scattering length in a wide range. With magnetic fields between 17 G and 30 G a positive scattering length should allow for a stable Bose-Einstein condensate.
In our experiments to explore the quantum gas properties of cesium, we use the gravito-optical surface trap (GOST) ovc97 . This optical dipole trap allows one to confine a large sample of atoms in an almost conservative environment with very efficient precooling by a Sisyphus-type mechanism. In a second stage, evaporative cooling is implemented to further increase the atomic phase-space density. In this article, we summarize the basic properties of the GOST (Sec. II) and discuss the limitations of optical cooling at high densities (Sec. III). We then report our first evaporative cooling results (Sec. IV) and discuss the prospects of future experiments (Sec. V).
## II Gravito-optical surface trap
### II.1 Trapping potential and general properties
A schematic overview of the geometry of the trap is given in figure 1. The GOST is an “optical mug”, whose bottom consists of an evanescent-wave (EW) atom mirror generated by total internal reflection of a blue-detuned laser beam from the surface of a prism, while the walls are formed by an intense hollow beam (HB) which passes vertically through the prism surface.
The steep exponential decay of the EW intensity along the vertical direction and the sharp focussing of the hollow beam lead to large intensity gradients and thus in combination with the blue detuning of both light fields to a strong repulsive dipole force. We exploit this fact to efficiently keep the atoms in the dark inner region of the trap where the unwanted effect of heating through scattering of trapping light photons is suppressed. In addition to that, the concept also features a large trapping volume which allows for a transfer of a large number of atoms into the GOST. Due to the accurate focussing of the hollow laser beam into a ring-shaped intensity profile man98 , the shape of the potential is box-like along the horizontal directions whereas the combination of gravity and the repulsive wall of the EW leads to a wedge-shaped potential vertically.
In thermal equilibrium this geometry allows for a simple description of the ensemble. The vertical density profile is given by the “barometric” equation
$$n(z)=n_0\mathrm{exp}(z/z_0)$$
(1)
with $`n_0`$ being the peak number density and $`z_0=k_BT/mg`$ the $`1/e`$-height of the sample. $`m`$ denotes the mass of the cesium atom, $`g`$ the gravitational acceleration, $`z`$ the vertical coordinate and $`T`$ the temperature. Assuming a homogeneous horizontal density distribution within the hollow beam radius $`r_{HB}`$ this expression can be integrated to express the peak density $`n_0`$ in terms of the atom number $`N`$ and the temperature $`T`$,
$$n_0=\frac{mg}{\pi r_{HB}^2k_B}\frac{N}{T}.$$
(2)
The mean density $`\overline{n}`$ defined as
$$\overline{n}\frac{1}{N}d^3rn^2(𝐫)$$
(3)
is exactly half of the peak density ($`\overline{n}=n_0/2)`$. In the GOST potential the scaling of density with temperature ($`nT^1`$) is somewhat less as compared to a 3D harmonic oscillator where $`nT^{3/2}`$.
The mean energy per atom is $`\overline{E}=5/2k_BT`$, of which $`3/2k_BT`$ are kinetic energy and one $`k_BT`$ potential energy in the vertical gravitational field, as can be obtained using the Virial theorem. In the horizontal degree of freedom the box form of the confining potential leads to negligible potential energy.
While the horizontal extension of the sample is determined by the diameter of the hollow beam (2 $`r_{HB}`$), the height is proportional to the ensemble temperature with a constant of proportionality of 6.4 $`\mu `$m/$`\mu `$K. At typical equilibrium temperatures of a few $`\mu `$K the $`1/e`$-height is on the order of a few ten $`\mu `$m thus leading to a highly anisotropic “pancake-shape” of the ensemble.
### II.2 Experimental realization
The experimental constituents of the GOST are the evanescent-wave diode laser, a titanium:sapphire laser to create the hollow beam and an additional diode laser to provide the light for the repumping beam. For EW and repumping beam we use laser diodes (SDL-5712-H1, distributed Bragg reflector), which yield up to 100 mW of radiation at the cesium D<sub>2</sub>-line at a wavelength of 852 nm.
The EW laser is focussed to provide a round spot on the surface of the prism with a $`1/e^2`$-radius of 540 $`\mu `$m. The angle of incidence is $`\theta =45.6^{}`$ which is $`2^{}`$ above the critical angle of $`43.6^{}`$. This leads to a $`1/e^2`$ decay length of $`\mathrm{\Lambda }=(\lambda /2\pi )(n^2\mathrm{sin}^2\theta 1)^{1/2}500`$ nm, where $`n=1.45`$ is the refractive index of the fused-silica prism.
The EW detuning $`\delta _{EW}`$ is between a few GHz during optical Sisyphus cooling and up to 300 GHz in the process of evaporative cooling. Since atoms are predominantly kept in the lower hyperfine ground state ($`F`$=3) in the dipole trap, all detunings given throughout this paper refer to the transition between this state and the center of the <sup>2</sup>P<sub>3/2</sub> excited-state manifold.
At $`\delta _{EW}/2\pi =3`$ GHz and a power of 45 mW the EW creates a repulsive optical potential with a height of $`1`$ mK. The potential barrier is reduced by typically a factor of two by the attractive van-der-Waals interaction between the atoms and the dielectric surface lan96a .
The hollow beam is generated using an axicon optics man98 to create a ring-shaped focus of an inner and outer 1/e-radius of $`r_{HB}=260\mu `$m and $`r_{HB}+\mathrm{\Delta }r_{HB}=290\mu `$m, respectively. It has a power of 350 mW and its detuning is in the range between $`0.3`$ nm and $`2`$ nm. The HB provides a potential barrier on the order of $`100\mu `$K height. At this detuning the HB potential is almost conservative as the photon scattering rates can be estimated to be on the order of a few photons up to a few ten photons per second.
The repumping beam needed for the optical Sisyphus cooling soe95 is resonant with the $`F=4F^{}=4`$ hyperfine transition of the $`D_2`$-line and has an intensity on the order of 1 $`\mu `$W/cm<sup>2</sup>. It is shone on the trapping region from above (see fig. 1).
### II.3 Loading of the trap
The loading scheme of the GOST goes along standard ways starting from an effusive atomic beam which is decelerated by a Zeeman slower inkaDr . A magneto-optical trap (MOT) collects about $`3\times 10^8`$ atoms under ultrahigh vacuum conditions ($``$10<sup>-11</sup> mbar). The sample is then cooled to $`10\mu `$K and spatially compressed using a polarization gradient cooling scheme, in which the MOT laser detuning is increased within 50 ms from 3 $`\mathrm{\Gamma }`$ to 14 $`\mathrm{\Gamma }`$; here $`\mathrm{\Gamma }/2\pi =5.3`$MHz denotes the natural linewidth. After the atomic cloud is shifted from its loading position (3 mm above the surface) to a release position (0.5 mm) close to the evanescent wave using magnetic offset fields, the MOT-laser beams are shuttered and the atoms fall onto the EW-light sheet. As soon as the near-resonant light fields of the MOT are switched off, the repumping beam of the GOST is turned on to optically pump the atoms into the lower hyperfine state ($`F`$=3), and the Sisyphus cooling in the GOST starts. Initially up to $`2\times 10^7`$ atoms are transfered into the dipole trap and undergo optical cooling.
The number of atoms $`N`$ remaining in the GOST after a variable storage time is measured by recapturing them into the MOT and taking a fluorescence picture using a slow-scan CCD camera. The integrated fluorescence signal, calibrated with a more accurate absorption image of the MOT, allows us to determine $`N`$ with an estimated uncertainty of less than 50%.
## III Optical cooling at high densities
Here we discuss the efficient EW Sisyphus cooling mechanism and its role as optical precooling stage for evaporative cooling in the GOST. We present our observations of density-dependent effects that limit EW cooling, namely an excess temperature and trap loss due to ultracold collisions in the presence of blue-detuned light.
### III.1 Evanescent-wave Sisyphus cooling
The optical cooling process is based on inelastic reflections of the repeatedly bouncing cesium atoms from the evanescent wave soe95 ; des96 . In the great majority of the bounces the atoms are coherently reflected in the lower hyperfine state ($`F`$=3) without any dissipation of kinetic energy or heating. However, occasionally the absorption of an EW photon takes place and the subsequent spontaneous decay will either have the atom end up in the $`F`$=3- or the $`F`$=4 hyperfine state. A decay into the upper hyperfine state occurs with a branching ratio of $`q`$ = 0.25 and due to the reduced dipole force on this state will lead to a weaker repulsion of the atom from the EW and thereby to a damping of the vertical motion. After leaving the EW the atom is optically pumped back into the lower hyperfine state. Simple considerations on the balance of cooling through this process and heating due to photon scattering lead to an expression for the equilibrium temperature:
$$T=\left(\frac{1}{q}+\frac{1}{q_\mathrm{r}}\right)\left(1+\frac{\delta _{\mathrm{EW}}}{\delta _{\mathrm{HFS}}}\right)T_{\mathrm{rec}}.$$
(4)
The first term in brackets represents the average number of photons scattered per cooling reflection. $`q_r=5/12`$ denotes the branching ratio of the decay of the $`F`$’=4 excited state into the $`F`$=3 ground state. $`\delta _{\mathrm{EW}}`$ is the detuning of the evanescent wave and $`\delta _{\mathrm{HFS}}=2\pi \times 9.2`$ GHz the hyperfine splitting of the ground state. $`T_{\mathrm{rec}}=200`$nK is the recoil temperature of cesium. The experimental parameters yield a cooling limit of slightly less than 2 $`\mu `$K.
These considerations are consistent with observed temperatures of 2.0 $`\mu `$K in small samples ($`N10^5`$ atoms) at large hollow beam detunings (see figure 2). Dense samples with $`N10^5`$ atoms show an excess temperature as will be discussed in the following section.
The time scale on which the cooling process takes place is on the order of one second – considerably longer than in a MOT. Theoretical considerations on the cooling dynamics soe95 predict an exponential drop of the vertical temperature in the beginning of the cooling process with a rate constant
$$\beta =\frac{q}{3}\frac{\delta _{HFS}}{\delta _{EW}}\frac{mg\mathrm{\Lambda }}{\mathrm{}(\delta _{EW}+\delta _{HFS})}\mathrm{\Gamma }.$$
(5)
At $`\delta _{EW}=2\pi \times 3`$ GHz we find $`\beta =1.1`$ s<sup>-1</sup>.
Cooling of the horizontal motion is facilitated through mixing of the vertical and the horizontal degrees of freedom by either diffuse reflection of atoms from the EW ovc97 ; lan96b or in the case of dense samples by elastic collisions. Right after the transfer from the MOT into the GOST, the potential energy from the fall is gradually converted into thermal energy and leads to a very hot ($`100\mu `$K) sample after thermalization within $`0.5`$s. The few seconds it takes for the ensemble to reach the equilibrium temperature are consistent with the calculated cooling rate.
To measure temperatures we use a release-and-recapture method which is accomplished by turning off the EW-potential for a short duration (few milliseconds) and measuring the remaining fraction of atoms as a function of this release time. A fit of a theoretical model to the data which is based on a Boltzmann distribution inkaDr , yields $`T`$ as the only fit parameter. Under the conditions of the experiments reported here the sample is almost thermalized at any time, so that separate measurements of the horizontal temperature were not routinely performed.
### III.2 Excess temperature in dense samples
Measurements of $`T`$ in large atomic samples yield significantly higher temperatures ($``$$`10\mu `$K) compared to what we found in earlier experiments on small samples ovc97 . It turns out that the equilibrium temperature of EW Sisyphus cooling depends critically on the number of trapped atoms. Figure 3 shows the strong dependence of $`T`$ on atom number for typical operating conditions of the GOST. The data is reasonably described by a linear dependence
$$T=T_0+aN,$$
(6)
where $`T_0=4.5\mu `$K is the limit temperature achieved in small samples and $`a=1.5\mu `$K$`/(10^6`$ atoms) the slope. The limit temperature $`T_0`$ can be reduced close to the EW Sisyphus cooling limit (see eq. 4) by increasing the detuning of the hollow beam (see figure 2), which is easily explained by heating due to photon scattering. However at large detunings it is not possible to transfer large numbers of atoms into the GOST so that one has to find a compromise between appreciable atom number and low temperature.
The excess temperature as observed in large samples indicates the presence of multiple photon scattering. Similar observations are made in experiments on “gray molasses” where one also keeps the atoms predominantly in the $`F=3`$ ground state boi96 . We find that the slope $`a`$ decreases with increasing HB detuning. It is however not obvious that this excess heating solely depends on the hollow beam. The repumping beam of EW Sisyphus cooling is likely to contribute to this effect.
In a new experimental setup, in which the HB optics was substantially improved to provide aberration-free focussing ryc00 , we are now able to trap $`10^7`$ atoms at a HB detuning of 1 nm and achieve a temperature of $`10\mu `$K. This is better by a roughly a factor of two as compared to the conditions of the reported experiments. The improved setup was used for the evaporation experiments in Sec. IV.
### III.3 Binary collisions in a blue-detuned light field
Lifetimes of samples of more than $`10^6`$ atoms were found to be on the order of 5 to 10 seconds and clearly indicated the presence of a strong nonexponential contribution to the decay (see figure 4). We model the decay assuming the presence of a two-body loss process in the standard loss rate equation wei99
$$\dot{N}=\alpha N\beta \overline{n}N.$$
(7)
To solve this equation one has to consider the complete dependence of the mean density $`\overline{n}`$ on the atom number $`N`$, which is also influenced by the observed excess temperature. Since cooling and thermalization takes place on a time scale much shorter than the decay one can assume stationary conditions for $`T(N)`$ according to eq. 6. Using this and equation 2 one obtains
$$\overline{n}=\frac{mg}{2\pi r_{HB}^2k_B}\frac{N}{T_0+aN}.$$
(8)
Solving the differential equation one can almost perfectly fit the experimental data and extract the quadratic loss coefficient $`\beta `$ as a fit parameter. Figure 4 demonstrates the good agreement between the model and the measurement. We find $`\beta `$ to be on the order of $`10^{12}`$ cm<sup>3</sup>/s, which is about an order of magnitude less as compared to radiative escape in a MOT wei99 .
A possible mechanism for the observed loss is based on the excitation of a pair of colliding cesium atoms into a repulsive molecular state bal94 ; hof96 ; bur96 ; vul99b . The resonant dipole-dipole interaction splits the excited molecular state into an attractive and a repulsive branch. While the attractive branch gives rise to loss processes like radiative escape and photoassociation in red-detuned light, the repulsive branch becomes relevant in blue-detuned light. When a cesium pair is excited into a repulsive molecular state at the classical Condon point it is accelerated along the potential curve and obtains a kinetic energy equivalent to the detuning of the exciting light field with respect to the atomic transition. For a dipole trap the resulting energy gain is in general much larger than the trap depth and leads to an ejection of both colliding atoms from the trap. According to a simple semi-classical model the corresponding rate coefficient $`\beta `$ should scale with the laser detuning and intensity as $`I/\delta ^2`$ as long as the intensity is low enough to avoid optical shielding effects bal94 .
In a first set of measurements to investigate this loss process we obtained decay curves at different detunings of the EW-laser field and extracted the rate coefficient $`\beta `$. In a gravito-optical trap the mean EW dipole potential experienced by the bouncing atoms is given by $`\overline{U}_{dip}=mg\mathrm{\Lambda }/2`$, independent of intensity and detuning. Therefore the mean intensity $`\overline{I}`$ to which an atom is exposed is proportional to the detuning $`\delta _{EW}`$. Thus the expected scaling of the rate coefficient is $`\beta \overline{I}/\delta _{EW}^21/\delta _{EW}`$. Measurements in the detuning range between 1 GHz and 7 GHz ham99 indeed showed the expected $`1/\delta _{EW}`$-behaviour. This confirms that collisions in the evanescent wave can explain the observed loss.
As a practical consequence of this fact we can minimize trap loss during optical cooling by switching the EW detuning to higher values shortly after loading. Only in the initial cooling phase low detuning is required to provide a sufficiently large potential barrier. Therefore, in experiments aiming at high densities, the detuning is routinely switched from 3 GHz to about 7 GHz after 0.5 s.
To prove that excitation into repulsive molecular states leads to trap loss we use an additional blue-detuned “catalysis” laser to induce collisional trap loss. In this way this process can be studied independently of the parameters of the trapping potential. The rate coefficient $`\beta _{\mathrm{cat}}`$ was measured for different values of detuning and intensity of the catalysis laser. In additional temperature measurements we ensured that heating due to photon scattering from the catalysis laser was negligible for all combinations of intensity and detuning. Data was collected at different catalysis laser intensities. The results for $`\beta _{\mathrm{cat}}`$ shown in figure 5 agree with the expected scaling behaviour $`I/\delta ^2`$.
As an important conclusion of these results, trap loss by binary collisions involving repulsive excited molecular states plays a significant role in optical cooling. Nevertheless, a dense sample of many atoms at high densities can prepared. If the evanescent wave is then operated at very large detunings, a conservative trapping potential is obtained with very low photon scattering and strongly suppressed trap loss, i.e. a very good starting point for evaporative cooling.
## IV evaporative cooling
The GOST offers favorable conditions to implement forced evaporative cooling ket96 with the prospect to attain quantum-degeneracy of cesium or a two-dimensional quantum gas. In contrast to red-detuned dipole traps used for evaporation experiments ada95 ; eng00 , the spatial compression of the cold sample essentially results from gravity and is thus not affected when the optical potentials are ramped down. Moreover, many more atoms are initially loaded into the GOST as compared to typical red-detuned traps. Here we describe our first experiments demonstrating the feasibility of efficient evaporation in the GOST.
In our evaporation experiments we operate the trap at a HB detuning of $`1`$ nm and keep the EW parameters as in the measurements described above. In two seconds of optical cooling we prepare $`N=10^7`$ atoms at a temperature of $`T=10\mu `$K, and a peak density of $`n_0=6\times 10^{11}`$cm<sup>-3</sup>. For the unpolarized sample in the seven-fold degenerate $`F=3`$ ground state this corresponds to a peak phase-space density of $`D=n_0\lambda _{dB}^3/710^5`$ where $`\lambda _{dB}=h/\sqrt{2\pi mk_BT}`$ is the thermal de-Broglie wavelength. Elastic collisions take place at a rate on the order of $`50`$ s<sup>-1</sup> and, considering the resonant scattering of cesium arn97 , lead to a thermalization time of about 200 ms.
To implement forced evaporation we lower the EW potential by ramping up the EW detuning. This simultaneously reduces heating due to photon scattering and suppresses loss through the collision mechanism described above. Within $`4.5`$ seconds the EW detuning is increased exponentially from initially 7 GHz up to 250 GHz. This is accomplished by rapid mode-hop free temperature tuning of the EW diode laser. In the last $`2.5`$ seconds of the evaporation ramp the intensity of the hollow beam is reduced from $`350`$mW to $`11`$mW in order to reduce possible heating by residual light in the dark center of the hollow beam. The contribution of the HB potential ramp to the evaporation remains very small.
The experimental results are shown in figure 6. About 1 s after starting the exponential ramp, the temperature begins to drop \[filled squares in (a)\]. At the end of the ramp, it has reached $`T300`$ nK. This decrease of $`T`$ by about 1.5 orders of magnitude is accompanied by a decrease of the particle number $`N`$ \[open triangles in (a)\] from $`10^7`$ down to $``$3$`\times 10^4`$, i.e. about 2.5 orders of magnitude.
Although the number density $`n_0N/T`$ \[open circles in (b)\] decreases by about one order of magnitude, the phase-space density $`Dn_0T^{3/2}NT^{5/2}`$ \[filled squares in (b)\] shows a substantial increase by 1.5 orders of magnitude. At the end of the ramp, we obtain a phase-space density of $``$$`3\times 10^4`$.
In the regime of resonant elastic scattering ($`T1\mu `$K hop00 ), the relevant cross section scales $`T^1`$. In the GOST potential, this leads to a scaling of the elastic scattering rate and the thermal relaxation rate $`NT^{3/2}`$. Therefore the elastic scattering rate is almost constant for the conditions of our experiments. However, no runaway regime is reached.
An obvious problem in our evaporation scheme is that, for the applied exponential ramp, it takes about one second until the EW potential barrier becomes low enough to start the evaporation. Up to this point already about 50% of the particles are lost, presumably by the collisional mechanism investigated in Sec. III.3. After the corresponding initial loss of phase-space density the later evaporation then leads to a gain of almost two orders of magnitude. This clearly shows that the potential of evaporative cooling in the GOST is much larger than we could demonstrate in our first experiments discussed here.
Although the reported evaporation results are still quite preliminary they already show that efficient evaporation is attainable in the GOST by ramping down the optical trapping potential. Substantial improvements can be expected by an optimized evaporation ramp of the EW in combination with a simultaneous evaporation through the hollow-beam potential.
## V Conclusions and outlook
We have investigated optical and evaporative cooling in the gravito-optical surface trap. At high densities, optical cooling by inelastic reflections from the evanescent-wave bottom of the trap was found to be limited by an excess temperature, which we interpret as a result of multiple photon scattering. In addition, a loss process is induced by the blue-detuned trap light. Colliding atoms are excited into a repulsive molecular state which is followed by an energy release into the relative motion. Nevertheless, very good starting conditions are obtained for evaporative cooling. By reducing the trapping potentials we have cooled the sample down to a temperature of 300 nK and obtained a phase-space of $`3\times 10^4`$.
So far all our evaporation experiments have been performed with unpolarized atoms in the seven-fold degenerate $`F=3`$ ground state. Dramatic improvements can be expected by polarizing the atoms into the absolute ground state $`F=3,m_F=3`$. The phase-space density can be increased by a factor of seven, and by producing identical bosons the elastic collsion rate goes up by nearly a factor of two. Moreover, Feshbach tuning can be applied to modify the $`s`$-wave scattering length vul99a . We have already performed first experiments, in which we have demonstrated that atoms the GOST can be optically pumped into the state of interest polarizing .
The spatial compression of the atomic gas can be enhanced by adding a far red-detuned laser beam to the GOST. If such a beam propagates vertically in the center of the hollow beam it provides an additional potential well for transverse confinement. As an example, a 500-mW beam from a compact Nd:YAG laser focused to a 1/e-diameter of 0.1 mm already provides a dipole potential $`U_{\mathrm{red}}`$ of $``$3$`\mu `$K depth. In the GOST such an additional well can be loaded by elastic collisions, which in thermal equilibrium would result in a peak density enhancement of $`\mathrm{exp}(|U_{\mathrm{red}}|/k_BT)`$ along with a corresponding increase in phase-space density sta98 .
Another interesting option with an additional far red-detuned laser beam is to create a double-EW trap as suggested in Ref. ovc91 . The combination of a repulsive blue-detuned EW in combination with an attractive red-detuned EW of much larger decay length would allow to create a wavelength-size potential well very close to the dielectric surface. In such a scheme, a situation can be realized where only one vertical bound state exist. The situation would then be similar to atomic hydrogen on liquid helium, for which a two-dimensional quantum gas has been reported saf98 . The realization of such a system with alkali atoms could provide much more insight into the physical behavior of such a 2D gas.
In addition to such experiments on quantum gas properties, the GOST also represents a very promising source of ultracold atoms for experiments related to, e.g., atom interferometry szr96 , atom-surface interactions cot98 ; mar00 , and quantum chaos sai98 .
## Acknowledgments
This work was supported by the Deutsche Forschungsgemeinschaft in the frame of the Gerhard-Hess-Programm. We thank D. Schwalm for continuous support and encouragement. One of us (V.D.) acknowledges a fellowship by the Konrad-Adenauer-Stiftung.
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# L483: a protostar in transition from Class 0 to Class I
## 1 Introduction
Class 0 objects are the youngest stellar objects known (André et al. And93 (1993), And99 (1999)). They commonly power bipolar outflows with extreme characteristics like a very high degree of collimation and evidence for shock processing of molecular gas even in cases of very low stellar luminosity (see Bachiller & Tafalla Bac99 (1999) for a recent review). In order for these outflows to evolve into the more quiescent (“standard”) outflows associated with Class I sources, rapid changes in outflow morphology and kinematics have to occur in the few $`10^4`$ yr that Class 0 lasts (André et al. And93 (1993)). These changes are most likely associated with changes in the source itself, which is undergoing its major phase of assembling via gravitational infall (e.g., Bontemps et al. Bon96 (1996), Mardones et al. Mar97 (1997)). Understanding how these first evolutionary changes of the stellar and outflow life occur is a major challenge to star formation studies, and it requires the simultaneous analysis of Class 0 objects, their outflows, and their dense gas environments. Here we present a molecular line study of the L483 core and its outflow powered by IRAS 18148-0440 (IRAS 18148 hereafter), a system that we find at the end of its Class 0 stage, starting its transition to become a Class I object.
The source IRAS 18148 in L483, first identified as an embedded object by Parker (1988a ), is one of the reddest low-mass sources known (Ladd et al. 1991a , 1991b ), and is located toward the Aquila Rift, at a most likely distance of 200 pc (Dame & Thaddeus Dam85 (1985)). Ladd et al. (1991a ) and Fuller et al. (Ful95 (1995)) estimate a source bolometric temperature (in the sense of Myers & Ladd Mye93 (1993)) of 50-60 K, and using the flux compilation by Fuller et al. (Ful95 (1995)) (their Fig. 4), we estimate a $`L_{\mathrm{smm}}`$/$`L_{\mathrm{bol}}0.9`$ % (also, Fuller et al. Ful95 (1995) fit the spectral energy distribution with a single-temperature dust model at 40 K, and from their 1.1mm flux, we estimate $`L_{\mathrm{bol}}`$/$`L_{1.1\mathrm{mm}}<2.5\times 10^4`$). These numbers suggest that IRAS 18148 is a Class 0 object (André et al. And93 (1993), And99 (1999), Chen et al. Che95 (1995)), as already proposed by Fuller et al. (Ful95 (1995)) and Fuller & Wooten (Ful00 (2000)), although it is less extreme than the prototype Class 0 source VLA1623 (André et al. And90 (1990), André et al. And93 (1993)). IRAS 18148 has a luminosity of about 10 L and drives a well-collimated bipolar CO outflow (Parker et al. 1988b , Par91 (1991), Fuller et al. Ful95 (1995), Bontemps et al. Bon96 (1996), Hatchell et al. Hat99 (1999)), and is associated with a variable H<sub>2</sub>O maser (Xiang & Turner Xia95 (1995)) and shocked H<sub>2</sub> emission (Fuller et al. Ful95 (1995), Buckle et al. Buc99 (1999)). NIR imaging of the source vicinity shows a well-defined, parabolic reflection nebula, which is optically invisible and coincides with the blue lobe of the CO outflow (Hodapp Hod94 (1994), Fuller et al. Ful95 (1995)). Ammonia observations by Goodman et al. (Goo93 (1993)), Fuller & Myers (Ful93 (1993)), Anglada et al. (Ang97 (1997)), and Fuller & Wootten (Ful00 (2000)) reveal that the L483 core is centrally concentrated, has a strong velocity gradient across it, and a gas kinetic temperature of about 10 K. H<sub>2</sub>CO and CS spectra toward the central source present strong self absorption with lines having brighter blue peak, a signature of infall motions (Myers et al. Mye95 (1995), Mardones et al. Mar97 (1997)).
The combination in L483 of Class 0 characteristics, like a low $`T_{\mathrm{bol}}`$ and infall asymmetry, together with the presence of a bright NIR nebula, indicative of partial core disruption, makes this source an interesting object to study the early evolution of a very young stellar object. To carry out such a study, we have observed L483 in tracers sensitive to different aspects early stellar life, like the outflow (CO, section 3.1), the dense core and possible chemical outflow anomalies (CH<sub>3</sub>OH, section 3.2), and infall and shock chemistry (H<sub>2</sub>CO, 3.3). From the combination of these observations, we propose that the central source in L483 has already started its transition toward Class I, and that the outflow has lost part of the chemical richness characteristic of Class 0 flows (section 4).
## 2 Observations
We observed L483 with the IRAM 30m telescope during several sessions in 1994 September, 1995 September, November, and 1996 June. Different receiver configurations were used to map the core in <sup>12</sup>CO(2–1) \[230.53799 GHz\], H<sub>2</sub>CO(2<sub>12</sub>–1<sub>11</sub>) \[140.839518 GHz\], and CH<sub>3</sub>OH($`2_k`$$`1_k`$) \[96.741420 GHz\], and to observe selected positions in <sup>13</sup>CO(1–0), C<sup>18</sup>O(1–0), C<sup>17</sup>O(1–0), H$`{}_{}{}^{13}{}_{2}{}^{}`$CO(2<sub>12</sub>–1<sub>11</sub>), H<sub>2</sub>CO(3<sub>12</sub>–2<sub>11</sub>), and SiO(2–1). Most observations were done in position switching mode (PSW) in order to obtain flat baselines. After searching for a clean off position, we settled with ($`600^{\prime \prime }`$, $`300`$) with respect to our map center ($`\alpha _{1950}=18^\mathrm{h}14^\mathrm{m}50\stackrel{s}{.}6`$, $`\delta _{1950}=4\mathrm{°}40^{}49\stackrel{}{.}0,`$ position of IRAS 18148-0440). This position lies outside the optical obscuration associated with L483 and seems free from dense gas tracer emission, although it has a weak <sup>12</sup>CO(2–1) line at a level of 2 K between $`V_{\mathrm{LSR}}=69`$ km s<sup>-1</sup> (outside the L483 range). A frequency switched (FSW) <sup>12</sup>CO(2–1) spectrum of this position was taken so it could be added to the data if needed, and a test was made by adding the off spectrum to a PSW spectrum from the origin and comparing the result with a FSW spectrum; the two were indistinguishable. The data shown in this paper, except when indicated, correspond to PSW observations without addition of the off position.
The backend was an autocorrelator split into different windows with resolutions ranging from 0.1 km s<sup>-1</sup> for <sup>12</sup>CO(2–1) to 0.03 km s<sup>-1</sup> for CH<sub>3</sub>OH($`2_k`$$`1_k`$). The telescope pointing was corrected frequently by observing continuum sources and is expected to be accurate within $`3^{\prime \prime }`$. The $`T_\mathrm{A}^{}`$ scale of the telescope was converted into $`T_{\mathrm{mb}}`$ using the main beam efficiencies recommended by Wild (Wil95 (1995)). The full width at half maximum of the telescope beam varies linearly with wavelength, and ranges from $`25^{\prime \prime }`$ for CH<sub>3</sub>OH($`2_k`$$`1_k`$) to $`11^{\prime \prime }`$ for <sup>12</sup>CO(2–1). Spectra were taken with $`10^{\prime \prime }`$ spacing, which is slightly less than one beam at the highest frequency.
## 3 Results
### 3.1 CO data
#### 3.1.1 Ambient CO emission
In this section we derive the basic parameters of the gas along the line of sight towards the L483 core center using a series of CO isotopomer lines. As Fig. 1 shows, the <sup>12</sup>CO(2–1) line is heavily self absorbed between $`V_{LSR}`$ 4.5 and 6.5 km s<sup>-1</sup>, with a $`T_{\mathrm{mb}}`$ of 4-5 K Given that the gas between these velocities is optically thick, its excitation temperature $`T_{\mathrm{ex}}=`$ has to be 8.5-9.5 K, a value close to the kinetic temperature of the dense gas (10 K, from NH<sub>3</sub> and HC<sub>3</sub>N, Fuller & Myers Ful93 (1993), Anglada et al. Ang97 (1997)). The gas responsible for the <sup>12</sup>CO(2–1) absorption represents low density gas in the outer layers of the cloud, far from the dense core, so the above result suggests that all the ambient gas along the line of sight has an almost constant temperature of 10 K (note that the outflow gas is warmer, see below).
In contrast with the self absorbed <sup>12</sup>CO(2–1) emission, C<sup>18</sup>O(1–0) and C<sup>17</sup>O(1–0) are optically thin. This is not evident because of the non-Gaussian shape of the spectra (Fig. 1), but can be proved using the following simple model: we create a C<sup>17</sup>O(1–0) spectrum by adding three replicas of the C<sup>18</sup>O(1–0) line each shifted in velocity by the proper hyperfine splitting and weighted by the optically thin relative intensity. The result, shown in lighter shade below the C<sup>17</sup>O(1–0) spectrum, matches very well the observations, indicating that both C<sup>18</sup>O(1–0) and C<sup>17</sup>O(1–0) are thin. The non Gaussian shape of these lines has therefore to result from velocity structure along the line of sight (see section 3.2 for further details).
With the excitation temperature of 10 K and the fact that C<sup>17</sup>O(1–0) is thin, we can derive the core central H<sub>2</sub> column density. We integrate the C<sup>17</sup>O(1–0) emission in velocity and assume local thermodynamic equilibrium, estimating a N(C<sup>17</sup>O) of 1.5 $`\times 10^{15}`$ cm<sup>-2</sup>. For a standard C<sup>17</sup>O abundance of $`4.7\times 10^8`$ (Frerking et al. Fre82 (1982), Wilson & Rood Wil94 (1994)), this value implies an H<sub>2</sub> column density of 3 $`\times 10^{22}`$ cm<sup>-2</sup>.
To finish this section, we note the presence in Fig. 1 of CO emission outside the ambient cloud range ($`V_{LSR}=4.5`$-6.5 km s<sup>-1</sup>). Part of this emission comes from outflow gas (discussed in section 3.1.2), but other part must come from additional clouds along the line of sight. This is the case for two features at $`V_{LSR}=3`$ and 8 km s<sup>-1</sup>, mostly seen in <sup>13</sup>CO, as they are spread over an area larger than 6 arcminutes. It must also be the case for some contribution between 2.5 and 4.5 km s<sup>-1</sup>, which is very prominent in the <sup>12</sup>CO spectra outside the outflow range (see (0,$`40`$) spectrum in Fig. 1). As we will see below, this low velocity emission limits our study of the outflow gas, to which we now turn our attention.
#### 3.1.2 The outflow
Fig. 2 presents our CO(2–1) map of the L483 outflow (see also Fig. 9). As previous maps (Parker et al. Par91 (1991), Fuller et al. Ful95 (1995), Bontemps et al. (Bon96 (1996)), and Hatchell et al. Hat99 (1999)), it shows that the accelerated CO emission extends E-W with IRAS 18148 at the center. The red CO lies to the east of IRAS 18148 while the blue CO lies to the west, although there is some red gas near the western tip of the blue lobe. This anomalous red gas coincides with a region of shocked H<sub>2</sub> emission (Fuller et al. Ful95 (1995)), and probably represents outflow material with an enhanced turbulent component because of the shock.
A comparison of the CO outflow with the K (2.1 $`\mu `$m) image of Hodapp (Hod94 (1994)) (Fig. 2) shows that the blue CO delineates very closely the reflection nebula and presents a relative maximum toward its center. In contrast, the red (eastern) CO is weak and highly collimated near IRAS 18148 and reaches a maximum 0.04 pc away from it. This east-west asymmetry of the CO outflow suggests that a similar asymmetry occurs in the underlying nebula, which would therefore have an intrinsically brighter western side.
As the <sup>12</sup>CO(2–1) spectrum in Fig. 1 shows, the outflow gas is warmer than the ambient core, or otherwise its emission would not be brighter than the self absorption by 10 K gas. The strongest CO wings imply excitation temperatures of 20 K (twice the ambient kinetic temperature), and this is a lower limit because the CO emission may not be thick and thermalized. In fact, Hatchell et al. (Hat99 (1999)) have argued, from a comparison of CO(4–3) and CO(2–1) lines, that temperatures up to 50 K could be present in the CO outflow. Given that we do not find brightness temperatures larger than 20 K, we take this value as a lower limit for our further calculations of the outflow energetics (see below) and molecular abundances (sections 3.2 and 3.3).
With the above value for the CO(2–1) $`T_{ex}`$, we estimate the outflow energetics. This requires special care because the outflow is rather slow and the CO spectra are contaminated at low velocities, so in the Appendix we present a method to deal with the contaminating emission. Of course, we cannot correct for the outflow gas hidden by the self absorption, so all CO emission in the central 3 km s<sup>-1</sup> (about $`V_{\mathrm{LSR}}=5.5`$ km s<sup>-1</sup>) will be ignored, and our estimate will represent a lower limit. As the highest outflow velocities, we take $`V_{LSR}`$ $`3`$ and 14 km s<sup>-1</sup>, because no faster outflow emission is found in the CO(2–1) velocity maps. Assuming a CO abundance of $`8.5\times 10^5`$ (Frerking et al. Fre82 (1982)), we derive an outflow mass of 0.01 M, a momentum of 0.03 M km s<sup>-1</sup>, and a kinetic energy of $`2\times 10^{42}`$ erg. From a total outflow length of $`160^{\prime \prime }`$ and a total velocity extent of 15 km s<sup>-1</sup>, we derive a kinematical time of $`10^4`$ yr. These values are in reasonable agreement with those from from Parker et al. (Par91 (1991)).
### 3.2 CH<sub>3</sub>OH and SiO data
Figure 3 presents a CH<sub>3</sub>OH($`2_k`$$`1_k`$) spectrum toward the core center with labels indicating the different $`k`$ components. Among the E-type lines, the $`k`$=1 component is not detected, while the the integrated intensity ratio between the $`2_1`$$`1_1`$ and $`2_0`$$`1_0`$ lines is approximately 7, very similar to the ratio found by Turner (Tur98 (1998)) for TMC-1 and L183. The relatively high intensity of the lines (about 2.5 K at the peak) implies an excitation temperature of at least 5.5 K, while an LTE rotation diagram analysis (e.g., Menten et al. Men88 (1988)) indicates temperatures of 4 K at most. This suggests that non LTE conditions may apply, and that for the E-type CH<sub>3</sub>OH, the lower lying $`k`$=$`1`$ ladder has a higher $`T_{\mathrm{ex}}`$ than the higher $`k`$=0 ladder. The excitation conditions are probably constant over the core, as an average over all (75) positions outside the central $`20^{\prime \prime }`$ gives a spectrum with the same intensity ratios among all the E and A components as in the spectrum shown in Fig. 3; we will use below this apparently constant conditions to infer a lower limit to the density gradient in the core. From the integrated intensities in the central spectrum, we estimate a CH<sub>3</sub>OH column density of about $`7\times 10^{13}`$ cm<sup>-2</sup>, which together with our estimated H<sub>2</sub> column density implies a CH<sub>3</sub>OH abundance of $`2\times 10^9`$, which is very close to the abundances estimated for other dark clouds (Friberg et al. Fri88 (1988), Bachiller et al. Bac95 (1995), Turner Tur98 (1998)).
The spatial distribution of the CH<sub>3</sub>OH($`2_k`$$`1_k`$) integrated intensity, shown in Fig. 4a, is rather round and centrally concentrated toward IRAS 18148. This is also the case with other dense gas tracers like HC<sub>3</sub>N, NH<sub>3</sub>, and N<sub>2</sub>H<sup>+</sup> (Fuller & Myers Ful93 (1993), Anglada et al. Ang97 (1997), Caselli et al. Cas00 (2000)), and is probably a sign of the extreme youth of the embedded star, which has not had time to perturb the bulk of the parental core (despite carving the nebula). The central concentration of the emission reinforces our interpretation that the emission is mostly optically thin, as otherwise we would expect a flat distribution (our unpublished NH<sub>3</sub> data indicate a constant gas kinetic temperature across the core).
Fuller et al. (Ful95 (1995)) have inferred an r<sup>-2</sup> density gradient within $`30^{\prime \prime }`$ of the IRAS source based on the intensity contrast between the two sides of the IR reflection nebula, but their calculation depends critically on assuming that the nebula is symmetric, which we have seen is probably not the case. More recently, Fuller & Wootten (2000) have proposed that the r<sup>-2</sup> profile continues to larger radii ($`100^{\prime \prime }`$) using a model for the NH$`{}_{3}{}^{}(1,1)`$ emission, although this result could be affected by interferometer missing flux. Here we use the CH<sub>3</sub>OH emission to derive an independent estimate, assuming that this emission is thin and the excitation temperature is constant. In Fig. 4b, we show a radial average of the emission, together with the results from our model for three density power laws: $`r^1`$, $`r^{1.5}`$, and $`r^2`$ (results convolved with a $`25^{\prime \prime }`$ Gaussian beam). As the figure shows, the $`r^2`$ density law is too steep, and the observations are better fit between $`R=15^{\prime \prime }`$ and $`100^{\prime \prime }`$ using a $`r^1`$ profile with probably some steepening for $`R>50^{\prime \prime }`$.
In the optically thin limit, the brightness radial profiles of Fig. 4 are proportional to column density radial profiles, so we can use them to estimate the core mass by normalizing them to the central H<sub>2</sub> column density and integrating them radially. In this way, we derive a core mass between 5 M ($`r^{1.5}`$ density profile) and 10 M ($`r^1`$ density profile). These values are in good agreement with the ammonia result from Anglada et al. Ang97 (1997), and also agree with a virial estimate using the average CH<sub>3</sub>OH line width over the core (0.64 km s<sup>-1</sup>), which gives 7 and 8 M for $`r^{1.5}`$ and $`r^1`$ density profiles, respectively.
Although the CH<sub>3</sub>OH lines are relatively narrow compared with other dense gas tracers like H<sub>2</sub>CO, they are systematically asymmetric and change velocity across the core. This is illustrated in Fig. 4c with a position-velocity diagram along the east-west axis of the core, which shows that the line center velocity shifts abruptly by about 0.3 km s<sup>-1</sup> near the position of IRAS 18148 (origin of offset coordinates); lines toward $`\mathrm{\Delta }\alpha 20^{\prime \prime }`$ are brighter, bluer, and narrower than toward the east. These two velocity components coincide with the C<sup>18</sup>O (and C<sup>17</sup>O) peaks we have found in section 3.1.1, and and can also be seen in our N<sub>2</sub>H<sup>+</sup>(1–0) data. Fuller & Myers (Ful93 (1993)), with low resolution observations, have reported that the spectra from L483 seem to have two velocity components, while Goodman et al. Goo93 (1993) found a systematic velocity gradient across the core, all in the same direction as the velocity change we find. The origin of this behavior is not clear, but we notice that the sense of the velocity change agrees with the sense of the bipolar outflow. It is therefore possible that it arises from the acceleration of dense gas by the outflow, since the position velocity diagram is similar to that of L1228 in C<sub>3</sub>H<sub>2</sub> (Tafalla & Myers Taf97 (1997)), where outflow acceleration is the cause of a similar velocity shift.
Independently of the origin of the velocity shift, it is clear that the CH<sub>3</sub>OH spectra in L483 do not show the prominent high-velocity wings seen in some low mass outflows like NGC 1333-IRAS 2 (Sandell et al. San94 (1994)) and L1157 (Bachiller at al. Bac95 (1995)). The wings in these systems arise from large abundance enhancements of CH<sub>3</sub>OH, which in the best studied case of L1157, amounts to a factor of 400 (Bachiller et al. Bac95 (1995), Bachiller & Pérez-Gutiérrez Bac97 (1997)). For L483, the non detection of CH<sub>3</sub>OH at the velocities with CO outflow implies that any possible CH<sub>3</sub>OH abundance enhancement cannot be larger than around 10, which is at least 40 times lower than in L1157. Higher signal-to-noise data can probably lower this limit by a significant amount.
Another tracer with prominent wings toward certain outflows is SiO (e.g., Bachiller & Pérez-Gutiérrez Bac97 (1997)), and three L483 positions were observed in the 2–1 line of this molecule (origin, $`30^{\prime \prime }`$ E, and $`30^{\prime \prime }`$ W). Our non detections (with limits of the order of 0.1 K km s<sup>-1</sup>) imply SiO abundances lower than $`10^{11}`$ and $`8\times 10^{10}`$ for the ambient gas and the outflow, respectively (numbers derived using an large velocity gradient analysis). Although the non detections make impossible to estimate an abundance enhancement in the outflow, comparing our outflow limit with the SiO abundance found in L1157 ($`10^7`$, Mikami et al. Mik92 (1992), Bachiller & Pérez-Gutiérrez Bac97 (1997)), we estimate that any possible enhancement in L483 is at least 100 times smaller than in L1157.
### 3.3 H<sub>2</sub>CO Data
Our other dense gas tracer, H<sub>2</sub>CO, presents very different lines than CH<sub>3</sub>OH. Fig. 5 shows a series of H<sub>2</sub>CO(2<sub>12</sub>–1<sub>11</sub>) spectra for different core positions illustrating the variety of profiles. At ambient velocities ($`V_{\mathrm{LSR}}5.5`$ km/s), the lines are strongly self absorbed over most of the core, as can be checked by comparing the double-peaked main isotope with the single-peaked rare H$`{}_{}{}^{13}{}_{2}{}^{}`$CO line toward the central position (Fig. 5 top panel). This self absorption is slightly red shifted with respect to the emitting gas, so the resulting spectrum has a brighter blue peak, a fact already noticed by Myers et al. (Mye95 (1995)) and Mardones et al. (Mar97 (1997)), who have presented preliminary versions of the H<sub>2</sub>CO spectrum towards the core center. Such red-shifted self absorptions are characteristic signatures of inward motions (e.g, Leung & Brown Leu77 (1977)), and their presence in L483 makes this core one of the best infall candidates known. A detailed study of the spatial distribution of the self absorption and its interpretation in terms of infall motions will be presented elsewhere (Mardones et al. Mar00 (2000), see also Mardones Mar98 (1998)), so here we limit ourselves to briefly comment on this feature.
Double-peaked H<sub>2</sub>CO spectra abound in the core and are more prominent along the outflow axis, although a map of the spectral asymmetry parameter $`\delta v`$ (defined as the difference between the thick and thin line peaks normalized to the thin line width, see Mardones et al. Mar97 (1997)) shows that the blue asymmetry is stronger perpendicular to the flow. Overall, despite the presence of high-velocity blue and red wings in the H<sub>2</sub>CO lines, the ambient self absorption is mostly red shifted and an average spectrum over the central core is clearly asymmetric in the sense of infall. To estimate the global infall rate in L483, we first determine the infall radius from the extent of the H<sub>2</sub>CO spectra with brighter blue peak, which we measure from the data as 0.02 pc. Then, we use the simple 2-layer model of Myers et al. (Mye96 (1996)), which allows to derive an infall velocity from the contrast between the blue and red peaks of a self absorbed line profile knowing the intrinsic line width from a thin tracer (see their Eq. 9). To do this, we take the average H<sub>2</sub>CO spectrum inside the infall radius, and use as thin tracer the average CH<sub>3</sub>OH spectrum over the same area. In this way, we derive a mean infall speed of 0.02 km s<sup>-1</sup>, which is clearly subsonic. Finally, we derive a mean density inside the infall radius using the power-law density models derived in our CH<sub>3</sub>OH analysis, which give a mean density of $`3.3\times 10^5`$ cm<sup>-3</sup> for a $`r^1`$ profile and $`2.8\times 10^5`$ cm<sup>-3</sup> for a $`r^{1.5}`$ profile. Averaging the above densities to $`3\times 10^5`$ cm<sup>-3</sup>, we derive a mass infall rate of $`2\times 10^6`$ M yr<sup>-1</sup> (see Mardones et al. Mar00 (2000) for further details).
At velocities outside the ambient range, the H<sub>2</sub>CO spectra present strong wings that change with position following the distribution of accelerated gas in the bipolar outflow. The spectra at ($`30^{\prime \prime }`$, 0) and ($`10^{\prime \prime }`$, 0) in Fig. 5 illustrate this effect, which implies that the bipolar outflow is accelerating part of the dense core gas, probably shocking it and altering its chemical composition (see below). To better illustrate the dense-gas acceleration traced by H<sub>2</sub>CO, we present in Figure 6 velocity maps for the blue, ambient, and red regimes. The blue and red maps agree very well with equivalent outflow maps from CO (see Fig. 2), especially for the blue lobe. There, both emissions present two maxima, one associated with the reflection nebula, and the other with the region of strong H<sub>2</sub> emission at the end of the lobe. Even more, both emissions present an “anomalous” red peak towards the position of bright H<sub>2</sub> emission, again suggesting this gas is shock related. These similarities suggest that CO and H<sub>2</sub>CO, despite their different dipole moments, are tracing the same (or very closely connected) gas.
As a further illustration of the importance of the outflow acceleration in the H<sub>2</sub>CO emission, we show in Fig. 7 a position-velocity diagram along the the Right Ascension axis (i.e., approximately parallel to the outflow). At ambient velocities, the strong self absorption can be easily seen in the form of two bright emission peaks at about 5.2 and 6 km s<sup>-1</sup>. At high velocities, the outflow wings form two triangular-shaped extensions which illustrate how the emission terminal velocity increases almost linearly with distance from the outflow source (at $`\mathrm{\Delta }\alpha =0`$); such linear increases are common in the CO emission from outflows (e.g., Meyers-Rice & Lada Mey91 (1991)). We can also see in the position-velocity diagram that towards the west ($`\mathrm{\Delta }\alpha 50^{\prime \prime }`$), both the wing and the ambient emission drop simultaneously. This suggests that the high velocity gas is breaking through the dense core, something also suggested by the presence of bright H<sub>2</sub> emission. Towards the east, the ambient gas is more extended than the outflow wing, so it seems the outflow is bounded in this direction.
We finish studying the H<sub>2</sub>CO abundance in the outflow and comparing it to the abundance in the ambient cloud. As the H<sub>2</sub>CO(2<sub>12</sub>–1<sub>11</sub>) line is self absorbed at ambient velocities, we use in this range the H$`{}_{}{}^{13}{}_{2}{}^{}`$CO isotopomer, which is very likely thin (see Fig. 5). In this way, using an LTE analysis with $`T_{\mathrm{ex}}=510`$ K and a <sup>12</sup>C/<sup>13</sup>C ratio of 77 (Wilson & Rood Wil94 (1994)), we derive a H<sub>2</sub>CO ambient column density of $`5\times 10^{13}`$ cm<sup>-2</sup> for the central position (a similar analysis using H<sub>2</sub>CO(3<sub>12</sub>–2<sub>11</sub>) gives the same number). This value implies an ambient abundance of $`1.5\times 10^9`$. For the outflow regime, we use the H<sub>2</sub>CO(2<sub>12</sub>–1<sub>11</sub>) line and compare it with the outflow CO emission. Applying the same LTE analysis as before, we estimate outflow abundances at ($`10^{\prime \prime }`$, 0) (blue) and ($`30^{\prime \prime }`$, 0) (red) of about $`3\times 10^8`$, suggesting an H<sub>2</sub>CO abundance enhancement of a factor of 20. This number is slightly smaller that the factor of 60-80 found in L1157 (Bachiller & Pérez-Gutiérrez Bac97 (1997)), but is significantly larger than our limit for the CH<sub>3</sub>OH enhancement. Thus, in contrast with L1157 for which CH<sub>3</sub>OH is enhanced by almost an order of magnitude more than H<sub>2</sub>CO, L483 is richer in H<sub>2</sub>CO.
## 4 Evolutionary status of IRAS 18148-0440
As mentioned in the Introduction, IRAS 18148 belongs to Class 0 (André et al. And93 (1993), André et al. And99 (1999)) due to its red spectral energy distribution, although it is less extreme than objects like VLA1623. IRAS 18148 has additional Class 0 characteristics like the presence of infall asymmetry in line spectra (Mardones et al. Mar97 (1997)). After analyzing its molecular environment in previous sections, we now compare this object and its outflow with other Class 0 sources, and to do this, we present in Table 1 a summary of properties of different Class 0 sources with similar luminosity ($`10L_{}`$) for which enough molecular data are available.
As Table 1 shows, L483 is very different kinematically from L1448-C, a source characterized by an extremely collimated and fast bipolar CO outflow with discrete high velocity components (“bullets”) (Bachiller et al. Bac90 (1990), see also Barsony et al. Bar98 (1998), O’Linger et al. Oli99 (1999), Eislöffel Eis00 (2000), and the similar object VLA1623 found by André et al. And90 (1990)). Despite our limited knowledge of the chemical composition of the L1448-C outflow, Table 1 shows that L483 also differs from L1448-C in its chemical enhancements, again in the sense of L1448-C being more extreme.
More similarities are found between L483 and L1157 and BHR 71 (L1157: Umemoto et al. Ume92 (1992), Mikami et al. Mik92 (1992), Bachiller et al. Bac95 (1995), Tafalla & Bachiller Taf95 (1995), Zhang et al. Zha95 (1995), Avery & Chiao Ave96 (1996), Bachiller & Pérez-Gutiérrez Bac97 (1997) Gueth et al. Gue98 (1998), Umemoto et al. Ume99 (1999), Zhang et al. Zha00 (2000); BHR 71: Bourke et al. Bou97 (1997), Garay et al. Gar98 (1998)). These outflows have comparable kinematics in the sense of having lower velocities and lacking “bullets” (also their collimation is similar), but differ by the amount of the chemical enhancement in the sense that the numbers for L483 are systematically smaller. Although the uncertainties in the chemical estimates are rather large, the differences shown in Table 1 are too extreme to be due to observational error. This means that among “chemically active” outflows, L483 is a weak case.
Bachiller & Tafalla Bac99 (1999) have argued that outflows evolve with time from having “bullets” to not having them, and from being chemically active to not being so (see Bontemps et al. Bon96 (1996) for evidence that outflow momentum flux systematically decreases with outflow age). If this evolution scenario is correct, IRAS 18148 in L483 would represent a rather late stage of a Class 0 source, and the objects in Table 1 would be ordered by increasing age from top to bottom. Unfortunately, no independent stellar clock exists yet to order by age the different Class 0 objects, and as Table 1 shows, the bolometric temperature $`T_{\mathrm{bol}}`$ cannot distinguish between extremely young sources, probably due to the lack of mid and far IR photometry. There is a reason, however, to suspect that IRAS 18148 is more evolved than other Class 0 sources, and that the above evolution scenario is correct: L483 has a bright NIR nebula (Hodapp Hod94 (1994), Fuller et al. Ful95 (1995)) in contrast with L1448-C and L1157 (Bally et al. Bal93 (1993), Hodapp Hod94 (1994), Davis & Eislöffel Dav95 (1995)) (but note BHR 71 also has a NIR nebula, Bourke et al. Bou97 (1997)). The presence of such a nebula suggests that the L483 outflow has been accelerating ambient material for longer than the outflows from the other objects in Table 1, and that IRAS 18148 is more advanced in its transition to become a visible object (e.g., Shu et al. 1987). Other Class 0 object with a NIR nebula is L1527 (Eiroa et al. Eir94 (1994)), so if the above scenario is correct, one would expect to find in this object a weak abundance enhancement of the molecules shown in Table 1. Observations of this object should be done to test this point.
The status of L483 as a somewhat evolved Class 0 object is also consistent with its relatively low value $`L_{\mathrm{smm}}`$/$`L_{\mathrm{bol}}`$ ($`0.9`$ %, compare with the 10 % of VLA1623, André et al. And93 (1993)) and with the work of Bontemps et al. Bon96 (1996). Applying the factor of 10 correction these authors apply to their CO data (see their Eq. 2), we estimate for L483 a momentum flux á la Bontemps et al. of $`3\times 10^5`$ M km s<sup>-1</sup> yr<sup>-1</sup>. This value is almost a factor of 2 lower than the mean momentum flux of Class 0 objects, but still 8 times larger than the average number for Class I sources. L483, again, appears as a Class 0 object already evolving toward Class I.
If the Class 0 encompasses sources as diverse as those powering the L1448-C and L483 outflows, outflow evolution should occur very rapidly during the star’s first few $`10^4`$ yr (the expected duration of the Class 0 stage, see André et al. And93 (1993)). It is possible that this evolution is driven by a rapid decrease in the infall/accretion rate on the central object (Bontemps et al. Bon96 (1996), Tomisaka Tom96 (1996), Henriksen et al. Hen97 (1997)), but although Class 0 objects do have stronger infall signatures than Class I sources, there is no clear infall trend among Class 0 objects themselves (Mardones et al. Mar97 (1997)). Further study of transition objects like L483 is needed to understand these earliest changes of stellar life, and as this work has shown, the combination of the chemistry and kinematics of the outflow may hold the key to that understanding.
## 5 Summary
We have observed the L483 core and outflow in different mm molecular transitions and made full maps in CO(2–1), CH<sub>3</sub>OH(2<sub>k</sub>–1<sub>k</sub>), and H<sub>2</sub>CO(2<sub>21</sub>–1<sub>11</sub>). With these data, we have studied the outflow, the core, and their relation with the IR cometary nebula around IRAS 18148. The main conclusions of our work are as follows:
1. The <sup>12</sup>CO emission at ambient velocities is extremely thick with the brightness temperature expected for gas at 9 K, the temperature previously estimated for the core gas. Outside the ambient regime the <sup>12</sup>CO lines present bright wings indicating outflow material warmer than the ambient gas by at least a factor of 2. A simple model for the C<sup>17</sup>O(1–0) emission towards the core center shows that this line is optically thin and non Gaussian due to the presence of two velocity components. From the integrated C<sup>17</sup>O(1–0) emission we estimate a central H<sub>2</sub> column density of $`3\times 10^{22}`$ cm<sup>-3</sup> in the inner $`20^{\prime \prime }`$.
2. The CO outflow emission is compact and slow, with a total length of 0.15 pc and a kinematical age of $`10^4`$ yr. Lower limits to the outflow mass, momentum, and energy are 0.01 M, 0.03 M km s<sup>-1</sup>, and $`2\times 10^{42}`$ erg, respectively. The CO outflow is asymmetric, with a blue lobe having a bright spot coinciding with the NIR nebula and the red lobe being weaker near IRAS 18148 and having a relative maximum $`45^{\prime \prime }`$ from the source. This asymmetry suggests that the reflection nebula around the IRAS source may also be asymmetric and have a more prominent blue side.
3. The CH<sub>3</sub>OH emission traces a dense core with no appreciable outflow wing contribution, although there is a shift in the line velocity along the direction of the outflow. The CH<sub>3</sub>OH emission is centrally concentrated on the IRAS position, and in the central $`200^{\prime \prime }`$ (0.1 pc) it decreases radially in a manner intermediate between what would be expected for optically thin emission with density power laws of $`r^1`$ and $`r^{1.5}`$. The estimated mass in this region is 5-10 M. No evidence for CH<sub>3</sub>OH or SiO abundance enhancement is found in the outflow.
4. The H<sub>2</sub>CO(2<sub>12</sub>–1<sub>11</sub>) emission is self absorbed at ambient velocities, and presents spectra with brighter blue peak, characteristic of inward motions, toward the central $`40^{\prime \prime }`$. With a simple model, we estimate an average infall speed of 0.02 km s<sup>-1</sup> and an infall rate of $`2\times 10^6`$ M yr<sup>-1</sup>. At high velocities, the H<sub>2</sub>CO(2<sub>12</sub>–1<sub>11</sub>) line presents bright wings in the same sense as the CO outflow wings, indicative of outflow acceleration. Comparing the H<sub>2</sub>CO and CO wing intensities we find that the H<sub>2</sub>CO abundance in the outflow regime is enhanced with respect to the ambient regime by a factor of 20.
5. The combination of CO, CH<sub>3</sub>OH, and H<sub>2</sub>CO data shows that the L483 outflow is less extreme than other outflows from Class 0 objects, like L1448-C and L1157, although it has some of their characteristics, such as gas heating and some abundance enhancement. We therefore suggest that the Class 0 source at the center of the L483 outflow is more evolved than other Class 0 sources, and it is in its transition to become a Class I object.
###### Acknowledgements.
This research has made use of the Simbad data base, operated at CDS, Strasbourg, France, and NASA’s Astrophysics Data System Abstract Service. The Digitized Sky Survey was produced at the Space Telescope Science Institute under US Government grant NAG W-2166. MT and RB acknowledge partial support from the Spanish DGESIC grant PB96-104, PCM acknowledges support from NASA Origins of Solar Systems grant NAG5-6266, and DM acknowledges partial support from grant FONDECYT 1990632.
## Appendix A Separation between outflow and ambient cloud emission
Estimating the outflow energetics for L483 is complicated because its low velocity and the presence of an extended component. In this appendix we discuss the method we have applied to correct for contamination by background gas seen in emission. Unfortunately, the part of the outflow emission absorbed by foreground gas cannot be recovered by any simple means.
The reason it is possible to correct for the extended emission in L483 is because this emission seems constant over the flow. This can be seen in Figure 8, where we present a map of spectra with the spectrum from ($`40^{\prime \prime }`$, $`40^{\prime \prime }`$) (a representative ambient position) superposed in lighter shade over each map position. As the figure shows, positions without outflow wings have spectra with the same shape, suggesting that the ambient cloud superposed to the outflow emission contributes everywhere with a similar spectrum. If the outflow emission is optically thin (as suggested by the lack of <sup>13</sup>CO(1–0) wings), this extended component can be subtracted out, leaving outflow-only emission as a residual. The origin of the extended component is in part ambient cloud background to the outflow (the foreground part causes the absorption and cannot be corrected for) and in part background emission from unrelated gas from the Aquila Rift, as discussed in section 3.1.1. The fact that its spectrum is rather flat topped suggests that this emission is partly thick with a kinetic temperature around 10 K, like that of the ambient L483 emission (but less extreme given its weaker <sup>13</sup>CO(1–0) emission). Given these characteristics, the gas that does not appear in absorption is most likely background to the outflow and therefore susceptible to correction.
To avoid adding noise in the process, we have used as background emission the average of all non outflow positions, and we have subtracted this spectrum to each observed position. The result seems to contain outflow emission only, as illustrated by the total integrated emission map of Figure 9, which is very similar to the outflow map in Fig. 2 (where no background subtraction was applied). To these background-subtracted spectra we have applied the standard energetics analysis (cf. Margulis & Lada Mar85 (1985)), and have ignored any contribution in the range $`V_{\mathrm{LSR}}=4`$-7 km s<sup>-1</sup>, as these velocities are contaminated by self absorption. For being forced to ignore these very low velocities, our estimates are necessarily lower limits to the real outflow parameters.
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# 1 Introduction
## 1 Introduction
The interest in hard exclusive reactions has recently been renewed in the context of skewed parton distributions (SPDs) . The SPDs, defined as hadronic matrix elements of bilocal products of quark or gluon field operators, are hybrid objects which combine properties of form factors and ordinary parton distributions. In fact reduction formulas reveal the close connection of these quantities. It has been shown that, at large photon virtuality, $`Q^2`$, and small momentum transfer, deeply virtual Compton scattering (DVCS) and deeply virtual electroproduction of mesons (DVEM) factorize into hard photon-parton scattering and SPDs describing the soft coupling between partons and hadrons. DVEM is dominated by longitudinally polarized photons for $`Q^2\mathrm{}`$; the cross section for transversally polarized photons is suppressed by $`1/Q^2`$. Complementary to the large $`Q^2`$ region is the large momentum transfer region (small $`Q^2`$). In this kinematical region Compton scattering off protons factorizes into a hard parton-level subprocess and a soft proton matrix element that is described by new form factors . These form factors represent $`1/x`$-moments of SPDs at large momentum transfer. Based on light-cone wave function overlaps as a model for the SPDs, detailed predictions for cross sections and polarization observables for real and virtual Compton scattering have been achieved in Refs. .
Here, in this work, we are going to apply the soft mechanism proposed in Refs. to electroproduction of flavor neutral pseudoscalar ($`P=`$ $`\pi ^0`$, $`\eta `$, $`\eta ^{}`$) and longitudinally polarized vector ($`V=`$ $`\rho ^0`$, $`\omega `$, $`\varphi `$) mesons. We will show that all arguments given in Ref. in order to establish factorization of Compton scattering, apply here too. I.e. provided the virtualities of the partons and their intrinsic transverse momenta, defined with respect to their parent proton’s momentum, are restricted by the proton’s wave function, the dominant contribution to electroproduction is generated from the handbag-type diagram shown in Fig. 1. It factorizes into meson electroproduction off partons and soft proton matrix elements described by the same type of form factors as appear in Compton scattering. It is shown in Ref. that, at large momentum transfer, there is one parton with a large virtuality that couples to the meson and forces the exchange of at least one hard gluon. We, therefore, follow the concept used in the calculation of DVEM and treat meson electroproduction off partons to leading-twist, lowest order perturbative QCD. A purely soft mechanism for large momentum transfer electroproduction of mesons, i.e. a soft overlap of the three light-cone wave functions for the hadrons involved is not possible . It is to be stressed however that the soft mechanism is not dominant for asymptotically large momentum transfer. In this limit the hard perturbative mechanism, for which all partons participate in the hard process, provides the leading contribution and the soft one merely represents a power correction. In this respect factorization of the soft mechanism is not on the same footing as the one, say, for DVEM, where the factorising diagrams are dominant for asymptotically large photon virtuality, and where factorization can be proven to hold in all orders of perturbation theory. The soft mechanism applies to photoproduction of mesons as well. However, the contributions from the hadronic component of the photon seem to dominate these processes for values of energy and momentum transfer accessible in current experiments.
It is also important to realize that the soft mechanism is complementary to the perturbative one, and both the contributions have to be taken into account in principle. However, recent developments, initiated by the CLEO measurements of the $`\pi \gamma `$ transition form factor and its theoretical analysis, e.g. , revealed that soft contributions play an important role in hard exclusive reactions at experimentally accessible momentum transfer which is of the order of a few GeV. Indeed, in the case of the electromagnetic form factor of the proton, the perturbative contribution has been shown to be small as compared to experiment provided the end-point regions, where one of the parton momentum fractions tends to zero, and where perturbative QCD is not applicable , are sufficiently suppressed. This can be achieved by employing the modified perturbative approach in which the transverse degrees of freedom and Sudakov suppressions are taken into account.
The soft contribution to large momentum transfer Compton scattering evaluated along the same lines as for the electromagnetic form factors, is in reasonable agreement with experiment . The perturbative contribution, on the other hand, has only been calculated to leading-twist accuracy and is way below the Compton data unless strongly asymmetric, i.e. end-point concentrated distribution amplitudes are used. These give, however, results for which the bulk of the contribution is accumulated in the soft end-point regions where the assumptions of leading-twist perturbative calculations break down. Even if asymmetric distribution amplitudes are utilized one obtains a perturbative contribution to Compton scattering that likely amounts to less than $`10\%`$ of the data for momentum transfer in the region of a few GeV ; the onset of the perturbative regime is expected to be above 10 GeV. The calculation of the leading-twist perturbative contribution to photoproduction of mesons has been attempted by Farrar et al. . The results are at drastic variance with experiment and need verification since the method for the numerical integrations used by Farrar et al. is questionable and known to fail in Compton scattering. On account of experience with electromagnetic form factors and Compton scattering, we will assume that the soft contribution to electroproduction of mesons are much larger than the perturbative ones for momentum transfers of the order of a few GeV and that the onset of the perturbative regime is beyond 10 GeV. There is still another contribution to electroproduction: it has two active partons, the photon couples to one of them while, by insertion of a hard gluon, the other one generates the vector meson. This contribution has the topology of the so-called cat’s ears diagrams. It has been discussed in Ref. that, in the large momentum transfer region, large virtualities or intrinsic transverse momenta occur in these diagrams forcing the exchange of additional hard gluons. It is reasonable to assume that the magnitude of the cat’s ears contribution is between the soft and the perturbative ones.
The paper is organized as follows: In Sect. 2 we will present the derivation of the soft mechanism. Next we will discuss the necessary phenomenological input that parameterizes the soft hadronic matrix elements (Sect. 3). In Sect. 4 we will comment on the case of photoproduction and then present our results for electroproduction of mesons (Sect. 5). In Sect. 6 we present our summary.
## 2 The soft mechanism
We are interested in electroproduction of mesons in the kinematical region where the Mandelstam variables $`s=(p+q)^2`$, $`t=\mathrm{\Delta }^2`$ and $`u=(pq^{})^2`$ are large on a hadronic scale, $`\mathrm{\Lambda }`$, of order 1 GeV. $`Q^2`$ is not considered as a large scale. Therefore the limit $`Q^20`$, the case of photoproduction, is included in the following. The calculation of soft contributions to the process of interest can be performed in full analogy to the case of Compton scattering; all its steps can be adopted straightforwardly. We can, therefore, restrict ourselves to an outline of the calculation; for details we refer to . The process amplitude is evaluated from the handbag-type diagram shown in Fig. 1 where also the four-momenta are defined (as usual $`Q^2=q^2`$). We work in a symmetric frame where the transverse momenta of the incoming and outgoing protons are treated in a symmetric way (see Fig. 2)
$$p=[p^+,\frac{m^2t/4}{2p^+},\frac{1}{2}𝚫_{}],p^{}=[p^+,\frac{m^2t/4}{2p^+},\frac{1}{2}𝚫_{}],$$
(1)
where $`m`$ is the mass of the proton. The plus and minus light-cone components of the momentum transfer are zero in this frame $`(\mathrm{\Delta }^+=\mathrm{\Delta }^{}=0)`$ and therefore $`t=𝚫_{}^2`$. The chief advantage of the symmetric frame is that the skewedness parameter, defined by
$$\zeta =\frac{\mathrm{\Delta }^+}{p^+}=1\frac{p^{}_{}{}^{}+}{p^+},$$
(2)
is zero. In order to specify the frame fully we further impose $`p_3+q_3=0`$. It coincides with c.m. frame for photoproduction with the 3-axis along $`𝐩+𝐩^{}`$.
The parton momenta are denoted by $`k_i`$ and $`k_i^{}`$. They are characterized by the usual momentum fractions
$$x_i=k_i^+/p^+,x_i^{}=k_i^+/p^+,$$
(3)
and the transverse components $`𝐤_{}_i`$ and $`𝐤_{}^{}_i`$. Because of $`\zeta =0`$ in the frame we are working, $`x_i=x_i^{}`$. The arguments of the light-cone wave functions are given by the momentum fractions and the intrinsic transverse parton momenta, i.e. the transverse components of the parton momenta in a frame where the transverse momentum of the parent proton is zero. By performing appropriate (transverse) boosts one finds for the light-cone wave function arguments of the incoming hadron
$$\stackrel{~}{x}_i=x_i,\stackrel{~}{𝐤}_i=𝐤_{}{}_{i}{}^{}+x_i𝚫_{}/2.$$
(4)
The arguments of the light-cone wave function of the scattered proton are $`\widehat{x}^{}=x_i^{}`$ and $`\widehat{𝐤}_i^{}=𝐤_{}^{}{}_{i}{}^{}x_i^{}𝚫_{}/2`$. For the sake of notation, we henceforth drop the subscripts for the active partons, i.e. for those participating in the subprocess that mediates the photon-meson transition (see Fig. 1).
The crucial hypothesis in the soft physics approach is now that the soft proton wave functions, i.e. the full wave functions with their perturbative tails removed from them, are dominated by parton virtualities in the range $`|k_i^2|,|k_i^2|\stackrel{<}{}\mathrm{\Lambda }^2`$ and by intrinsic transverse parton momenta satisfying $`\stackrel{~}{𝐤}{}_{i}{}^{2}/x_i,\widehat{𝐤}^2{}_{i}{}^{}/x_i^{}\stackrel{<}{}\mathrm{\Lambda }^2`$. <sup>1</sup><sup>1</sup>1 A restriction to intrinsic transverse momenta $`\stackrel{~}{𝐤}{}_{i}{}^{2}\stackrel{<}{}\mathrm{\Lambda }^2`$ instead of $`\stackrel{~}{𝐤}{}_{i}{}^{2}/x_i\stackrel{<}{}\mathrm{\Lambda }^2`$ fails as is shown in . At least one of the parton virtualities would be of order $`\mathrm{\Lambda }\sqrt{t}`$ and not $`\mathrm{\Lambda }^2`$. With the help of this hypothesis one can show that the subprocess Mandelstam variables, $`\widehat{s}=(k+q)^2`$ and $`\widehat{u}=(k^{}q)^2`$, are respectively equal to $`s`$ and $`u`$ up to corrections of order $`\mathrm{\Lambda }^2(t\pm Q^2)/t`$ provided $`s`$ and $`u`$ are large on a hadronic scale. This implies that the poles at $`\widehat{s}0`$ and $`\widehat{u}0`$ appearing in the lowest order Feynman graphs that contribute to the subprocess $`\gamma ^{}qMq`$ (see Fig. 3) are avoided and, hence, the pole contributions can be neglected.
The physical situation is that of a hard parton-level subprocess ($`\widehat{s},t,\widehat{u}\mathrm{\Lambda }^2`$) and the soft emission and reabsorption of a parton by the proton described by a soft proton matrix element. Hence, we can write the helicity amplitude for the process $`\gamma ^{}pMp`$ as
$`_{\mu ^{}\nu ^{},\mu \nu }^{M(q)}`$ $`=`$ $`{\displaystyle \underset{a}{}}ee_aB_a^M{\displaystyle d^4k\theta (k^+)\frac{d^4z}{(2\pi )^4}e^{ikz}}`$ (5)
$`\times `$ $`[p^{}\nu ^{}\left|T\overline{\psi }_{a\alpha }(0)\psi _{a\beta }(z)\right|p\nu H_{\mu ^{}\mu }^{M(q)\alpha \beta }(k^{},k)`$
$`+`$ $`p^{}\nu ^{}\left|T\overline{\psi }_{a\alpha }(z)\psi _{a\beta }(0)\right|p\nu H_{\mu ^{}\mu }^{M(q)\alpha \beta }(k,k^{})],`$
where $`H_{\mu ^{}\mu }^{M(q)}`$ is the tree-level expression for the hard scattering kernel. $`\mu `$ and $`\mu ^{}`$ respectively denote the helicities of the photon and the meson, $`\nu `$ and $`\nu ^{}`$ those of the protons. For the sake of legibility we label explicit helicities only by their signs, e.g. we write $`+`$, $``$ instead of $`+1/2`$, $`1/2`$ for fermions. The helicities are defined in the $`\gamma ^{}p`$ c.m. frame which is convenient for phenomenological applications and facilitates comparison with other results. On the other hand, the symmetric frame is adapted to discuss the reaction mechanism. The sum runs over quark flavors $`a`$, $`e_a`$ being the electric charge of quark $`a`$ in units of the positron charge $`e`$ and $`B_a^M`$ denotes the meson’s flavor wave function. The first term in (5) corresponds to the case where the incoming parton in the subprocess is a quark, the second term corresponds to an incoming antiquark. For the production of flavor neutral vector mesons gluons have to be considered as active partons too. We will discuss this contribution separately below.
Since the subprocess is dominated by a large scale, we can approximate the momenta $`k,k^{}`$ of the active partons in the subprocess as being on-shell, collinear with their parent hadrons
$$k[k^+,\frac{t}{8k^+},\frac{1}{2}𝚫_{}],k^{}[k^+,\frac{t}{8k^+},\frac{1}{2}𝚫_{}].$$
(6)
The integration over $`k^{}`$ and $`𝐤_{}`$ in (5) can then be performed explicitly leaving an integral $`𝑑k^+𝑑z^{}`$ and forcing the relative distance of fields in the matrix elements on the light cone, $`z\overline{z}=[0,z^{},\mathrm{𝟎}_{}]`$. After this the time ordering of the fields can be dropped .
The proton matrix element can be viewed as the amplitude for a proton with momentum $`p`$ emitting the active parton with momentum $`k`$ and a number of on-shell spectators times the corresponding conjugated amplitude for $`p^{},k^{}`$ summed over all spectator configurations, see Fig. 4. This corresponds to inserting a complete set of intermediate states between quark and antiquark fields in (5). Realizing that at the proton-parton vertices one has large plus components but, on account of the central hypothesis of small parton virtualities and small intrinsic transverse momenta, $`\stackrel{~}{𝐤}_i^2/x_i,\widehat{𝐤}_i^2/x_i^{}\stackrel{<}{}\mathrm{\Lambda }^2`$, one cannot form large kinematical invariants. With this feature of the soft mechanism at hand one can replace the products of fields in (5) by
$`\overline{\psi }{}_{\alpha }{}^{}(0)\psi _\beta (\overline{z})`$ $``$ $`\left({\displaystyle \frac{1}{2k^+}}\right)^2{\displaystyle \underset{\lambda ,\lambda ^{}}{}}\left(\overline{\psi }(0)\gamma ^+u(k^{},\lambda ^{})\right)`$ (7)
$`\times `$ $`\left(\overline{u}(k,\lambda )\gamma ^+\psi (\overline{z})\right)\overline{u}_\alpha (k^{},\lambda ^{})u_\beta (k,\lambda ),`$
where $`\lambda `$ and $`\lambda ^{}`$ denote the helicities of the active partons and $`u`$ their on-shell spinors. An analogous replacement is possible for the product $`\overline{\psi }_{a\alpha }(\overline{z})\psi _{a\beta }(0)`$. In this case antiquark spinors, $`v`$, appear. Due to this replacement the hard scattering kernels in (5) are multiplied with the spinors for on-shell (anti)quarks
$$_{\mu ^{}\lambda ^{},\mu \lambda }^{M(q)}=\overline{u}(k^{},\lambda ^{})H_{\mu ^{},\mu }^{M(q)}(k^{},k)u(k,\lambda ),$$
(8)
which guarantees electromagnetic gauge invariance of our result. The charge conjugation properties of Dirac matrices and spinors relate the subprocess amplitudes involving antiquarks to the quark amplitudes
$$\overline{v}(k,\lambda )H_{\mu ^{}\mu }^{M(q)}(k,k^{})v(k^{},\lambda ^{})=\kappa _M_{\mu ^{}\lambda ^{},\mu \lambda }^{M(q)},$$
(9)
where $`\kappa _V=1`$ for vector mesons and $`\kappa _P=+1`$ for pseudoscalar ones. The replacement (7) reveals that the plus components of the non-local currents dominate the proton matrix element and that the operators in the matrix elements are in fact the same as those of the leading-twist parton distributions occurring in deep-inelastic lepton-nucleon scattering, DVCS or DVEM. This is a nontrivial dynamical feature of large momentum transfer Compton scattering and electroproduction of mesons, given that, in contrast to the deeply virtual reactions, not only the plus components of the parton momenta but also their minus and transverse components are large now.
As mentioned above we follow the concept used in the calculation of DVEM and treat the formation of the helicity zero mesons ($`\mu ^{}=0`$) to leading-twist, lowest order perturbative QCD (cf. Fig. 3). In combination with the disregard of quark masses this formation mechanism leads to conservation of quark helicity in the subprocess, $`\lambda ^{}=\lambda `$. This feature and properties of massless spinors allow to simplify the expression (5) further, and to arrive at
$`_{0\nu ^{},\mu \nu }^{M(q)}`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{\lambda }{}}{\displaystyle \underset{a}{}}ee_aB_a^M{\displaystyle \frac{dk^+}{k^+}\theta (k^+)\frac{dz^{}}{2\pi }e^{ik^+z^{}}_{0\lambda ,\mu \lambda }^{M(q)}}`$ (10)
$`\times `$ $`[p^{}\nu ^{}|\overline{\psi }{}_{a}{}^{}(0)\gamma ^+\psi _a(\overline{z})+\kappa _M\overline{\psi }{}_{a}{}^{}(\overline{z})\gamma ^+\psi _a(0)|p\nu `$
$`+`$ $`\lambda p^{}\nu ^{}|\overline{\psi }{}_{a}{}^{}(0)\gamma ^+\gamma _5\psi _a(\overline{z})\kappa _M\overline{\psi }{}_{a}{}^{}(\overline{z})\gamma ^+\gamma _5\psi _a(0)|p\nu ].`$
Following , we take $`k^+=p^+`$, i.e. the light-cone fractions $`x=x^{}=1`$ in the hard scattering which is in line with the requirement to have no hard parton directly coupling to the protons. Admittedly, the global factor $`1/k^+`$ in (10) cannot be plainly associated with either the hard scattering or the soft matrix element. We therefore choose to keep $`k^+=xp^+`$ for this factor. We can now pull out the hard scattering amplitude from the integrals and use a form factor decomposition for the integrated proton matrix element
$`{\displaystyle _0^1}{\displaystyle \frac{dx}{x}}p^+{\displaystyle \frac{dz^{}}{2\pi }e^{ixp^+z^{}}p^{}\nu ^{}\left|\overline{\psi }_a(0)\gamma ^+\psi _a(\overline{z})+\kappa _M\overline{\psi }_a(\overline{z})\gamma ^+\psi _a(0)\right|p\nu }`$
$`=R_V^{Ma}(t)\overline{u}(p^{},\nu ^{})\gamma ^+u(p,\nu )+R_T^{Ma}(t){\displaystyle \frac{i}{2m}}\overline{u}(p^{},\nu ^{})\sigma ^{+\beta }\mathrm{\Delta }_\beta u(p,\nu ),`$
$`{\displaystyle _0^1}{\displaystyle \frac{dx}{x}}p^+{\displaystyle \frac{dz^{}}{2\pi }e^{ixp^+z^{}}p^{}\nu ^{}\left|\overline{\psi }_a(0)\gamma ^+\gamma _5\psi _a(\overline{z})\kappa _M\overline{\psi }_a(\overline{z})\gamma ^+\gamma _5\psi _a(0)\right|p\nu }`$
$`=R_A^{Ma}(t)\overline{u}(p^{},\nu ^{})\gamma ^+\gamma _5u(p,\nu ).`$ (11)
$`R_V^{Ma}`$, $`R_T^{Ma}`$ and $`R_A^{Ma}`$ are new form factors, depending on the type of the meson, $`V`$ or $`P`$, and on the flavor of the active quark. As the definition (11) reveals they are $`1/x`$-moments of SPDs at zero skewedness. The link-operator needed to render the definition of the SPDs gauge invariant, is not displayed in (11), i.e. we assume the use of a light-cone gauge combined with an appropriate choice for the integration path which reduces the link operator to unity. Due to time reversal invariance the form factors are real functions. The form factor $`R_T^{Ma}`$ is controlled by higher-twist dynamics and is expected to be suppressed by $`m^2/t`$ as compared with $`R_V^{Ma}`$ . Since the calculation of the soft contributions is only accurate up to corrections in $`\mathrm{\Lambda }^2/t`$, $`R_T^{Ma}`$ is to be omitted for consistency. Hence, we can only calculate the amplitudes conserving the proton helicity. Explicitly they read
$`_{0+,\mu +}^{M(q)}(s,t)`$ $`=`$ $`{\displaystyle \frac{e}{2}}\{_{0+,\mu +}^{M(q)}(s,t)[R_V^M(t)+R_A^M(t)]`$ (12)
$`+_{0,\mu }^{M(q)}(s,t)[R_V^M(t)R_A^M(t)]\},`$
where the form factors specific to the process $`\gamma ^{}pMp`$ are defined as
$$R_{V,A}^M(t)=\underset{a}{}e_aB_a^MR_{V,A}^{Ma}(t).$$
(13)
From parity invariance one has $`_{0\nu ,\mu \nu }^{M(q)}=\kappa _M(1)^\mu _{0\nu ,\mu \nu }^{M(q)}`$ and an analogous equation for the parton-level amplitudes $`_{0\lambda ^{},\mu \lambda }^{M(q)}`$. The amplitudes for longitudinally polarized photons simplify as a consequence of parity invariance: the vector form factor, $`R_V^M`$, contributes only in the case of vector meson production while the axial vector form factor, $`R_A^M`$, contributes in the case of pseudoscalar mesons. This is analogous to DVEM. For transversally polarized photons, on the other hand, both form factors contribute.
Let us now turn to the calculation of the parton-level amplitudes. The mesons are described by their valence Fock components and, for a longitudinally polarized vector meson, we write the corresponding matrix element in the usual way as
$`V,q^{}\left|\overline{\psi }(x)\gamma _\mu \psi (y)\right|0=q_\mu ^{}f_V{\displaystyle _0^1}𝑑\tau \varphi _V(\tau )e^{iq^{}(\tau x+\overline{\tau }y)},`$ (14)
where the proportionality between the meson’s polarization vector and its momentum, $`q^{}`$, for longitudinally polarized vector mesons is employed. For pseudoscalar mesons we have:
$`P,q^{}\left|\overline{\psi }(x)\gamma _5\gamma _\mu \psi (y)\right|0=iq_\mu ^{}f_P{\displaystyle _0^1}𝑑\tau \varphi _P(\tau )e^{iq^{}(\tau x+\overline{\tau }y)}.`$ (15)
The meson masses are ignored. $`\tau `$ is the fraction of the meson’s momentum the valence quark in the meson carries. The momentum fraction of the antiquark is $`\overline{\tau }=1\tau `$. $`f_M`$ is the meson’s decay constant and $`\varphi _M`$ its distribution amplitude which is normalized as
$$_0^1𝑑\tau \varphi _M(\tau )=1.$$
(16)
The definitions (14) and (15) are equivalent to the other frequently used ones $`(\text{/}q^{}/\sqrt{2})f_V\varphi _V/(2\sqrt{2N_c})`$ and $`(\text{/}q^{}\gamma _5/\sqrt{2})f_P\varphi _P/(2\sqrt{2N_c})`$. The color factor $`1/\sqrt{N_c}`$ (where $`N_c`$ denotes the number of colors) is not displayed in (14) and (15); it is taken into account in the parton-level amplitudes explicitly.
Working out the Feynman graphs shown in Fig. 3, one finds for the parton-level amplitudes
$$_{0+,\mu +}^{M(q)}=\mathrm{\hspace{0.17em}2}\pi \alpha _s(\mu _R)f_M\frac{C_F}{N_c}_0^1𝑑\tau \varphi _M(\tau )f_\mu ^{(q)}(\tau ),$$
(17)
where
$`f_+^{(q)}(\tau )`$ $`=`$ $`{\displaystyle \frac{\sqrt{2t}}{s+Q^2}}\left\{{\displaystyle \frac{(s+Q^2)(\tau s+Q^2)\overline{\tau }uQ^2}{\overline{\tau }s(\tau t\overline{\tau }Q^2)}}+{\displaystyle \frac{(s+Q^2)(\tau sQ^2)\overline{\tau }uQ^2}{\tau u(\overline{\tau }t\tau Q^2)}}\right\},`$
$`f_{}^{(q)}(\tau )`$ $`=`$ $`\overline{\tau }{\displaystyle \frac{\sqrt{2t}}{s+Q^2}}\left\{{\displaystyle \frac{u}{\overline{\tau }(\tau t\overline{\tau }Q^2)}}+{\displaystyle \frac{s}{\tau (\overline{\tau }t\tau Q^2)}}\right\},`$
$`f_0^{(q)}(\tau )`$ $`=`$ $`{\displaystyle \frac{2Q\sqrt{su}}{s+Q^2}}\left\{{\displaystyle \frac{u}{s(\tau t\overline{\tau }Q^2)}}+{\displaystyle \frac{s+Q^2+\overline{\tau }u}{\tau u(\overline{\tau }t\tau Q^2)}}\right\}.`$ (18)
$`C_F=(N_c^21)/(2N_c)`$ is the usual SU(3) color factor. Parity invariance fixes the amplitudes with negative quark helicities. For the scale of the parton-level amplitudes we choose $`\mu _R=s/4`$ which is roughly the average of the gluon and quark virtualities in the hard process.
In principle the amplitudes (17) hold for all values of $`t`$ and $`Q^2`$ provided the internal quark and gluon virtualities are sufficiently large. In the limit of either $`Q^20`$ or $`t0`$ the amplitudes (17) simplify strongly. In the case of photoproduction we find
$`_{0+,++}^{M(q)}`$ $`=`$ $`2\pi \alpha _s(\mu _R)f_M{\displaystyle \frac{C_F}{N_c}}1/\tau _M{\displaystyle \frac{\sqrt{2t}}{u}},`$
$`_{0+,+}^{M(q)}`$ $`=`$ $`2\pi \alpha _s(\mu _R)f_M{\displaystyle \frac{C_F}{N_c}}1/\tau _M{\displaystyle \frac{\sqrt{2t}}{s}},`$ (19)
and $`_{0+,0+}^{M(q)}=0`$ ($`Q`$ for $`Q^20`$). In deriving (19) we made use of the symmetry of the distribution amplitude for the mesons of interest under the interchange $`\tau \overline{\tau }`$, $`\varphi _M(\tau )=\varphi _M(\overline{\tau })`$. One observes that only the moment
$$1/\tau _M=_0^1𝑑\tau \frac{\varphi _M(\tau )}{\tau }$$
(20)
contributes. In the limit of large $`Q^2`$ and small $`t`$, the case of DVEM, the amplitude for longitudinally polarized photons, $`_{0+,0+}^{M(q)}`$, also becomes proportional to the $`1/\tau `$ moment, cf. for instance , while terms $`1/\tau ^2`$, $`1/\overline{\tau }^2`$ in the other two amplitudes signal the break-down of factorization for transversally photons . Inserting the parton-level amplitudes (17) into (12), one obtains the final expressions for the helicity amplitudes.
For flavor-neutral vector meson there is a complication which we now have to discuss, namely gluons have to be considered as active partons as well. Again this situation is similar to DVEM . In the kinematical region of $`\mathrm{\Lambda }^2Q^2s`$, characteristic of the HERA experiments, the gluon contribution even dominates . We start the calculation of the gluon contribution from an expression similar to (5)
$`^{V(g)}`$ $`=`$ $`{\displaystyle \underset{a}{}}ee_aB_a^V{\displaystyle d^4k\theta (k^+)\frac{d^4z}{(2\pi )^4}e^{ikz}}`$ (21)
$`\times `$ $`p^{}|TA^{\rho b}(0)A^{\rho ^{}b^{}}(z)|pH_{\rho \rho ^{}bb^{}}^{V(g)}(k^{},k),`$
where $`k`$ and $`k^{}`$ denote the momenta of the on-shell gluons in the symmetric frame, see (6). $`A_{\rho b}`$ is the gluon field with color $`b`$. For the sake of legibility we do not display helicity labels in this equation. The proton matrix element is only non-zero if $`b=b^{}`$. With this in mind we omit color labels in the following for convenience.
Now, we have to repeat all steps of the derivation of the quark contribution. As there, the use of the approximation (6) forces the relative distance of the fields in the proton matrix elements to the light cone, $`z\overline{z}`$. The important point is now the use of light-cone gauge, $`nA=0`$ (where $`n=[0,1,\mathrm{𝟎}_{}]`$) which allows to express the gluon field by an integral over the field strength tensor $`G_{\nu \mu }`$
$$A_\nu (\overline{z};n)=n^\mu _0^{\mathrm{}}𝑑\sigma e^{\epsilon \sigma }G_{\nu \mu }(\overline{z}+\sigma n).$$
(22)
(The limit $`\epsilon 0`$ is understood.) With the help of arguments similar to those leading to (7) we can replace the products of fields appearing in (21) by
$`A^\rho (0)A^\rho ^{}(\overline{z})`$ $`=`$ $`{\displaystyle \underset{\lambda ,\lambda ^{}=\pm 1}{}}ϵ^\rho (k,\lambda )ϵ^\rho ^{}(k^{},\lambda ^{}){\displaystyle 𝑑\sigma 𝑑\sigma ^{}e^{\epsilon \sigma \epsilon ^{}\sigma ^{}}}`$ (23)
$`\times `$ $`G_{\mu +}(\sigma ^{}n)G_{\mu ^{}+}(\overline{z}+\sigma n)ϵ^\mu (k,\lambda )ϵ^\mu ^{}(k^{},\lambda ^{}).`$
In the symmetric frame the polarization vectors of the on-shell gluons read
$$ϵ(k,\lambda )=[0,ϵ^{},ϵ_{}(\lambda )],ϵ(k^{},\lambda ^{})=[0,ϵ^{}{}_{}{}^{},ϵ_{}(\lambda ^{})],$$
(24)
where $`ϵ_{}(\pm 1)=(1,\pm i)/\sqrt{2}`$. The minus components need not be specified since they do not contribute in light-cone gauge.
The hard scattering kernels appearing in (21) are contracted by the first set of polarization vectors in (23) which leads to gauge invariant parton-level amplitudes
$$_{0\lambda ^{},\mu \lambda }^{V(g)}=ϵ^\rho ^{}(k^{},\lambda ^{})H_{\rho ^{}\rho }^{V(g)}(k^{},k)ϵ^\rho (k,\lambda )$$
(25)
Gluon helicity flip ($`\lambda =\lambda ^{}`$) is suppressed in the proton matrix element at large $`t`$ since two units of orbital angular momentum are required in order to avoid helicity flips of the proton ($`\nu =\nu ^{}`$). Thus, matrix elements involving helicity flips of the gluons are suppressed at least as $`m^2/t`$ and will be omitted. This argument is of importance only for longitudinally polarized photons because, for $`\mu =\pm 1`$, the parton-level amplitudes $`_{0\lambda ,\mu \lambda }`$ come out to zero in any case. For $`\lambda =\lambda ^{}`$ one may decompose the last part of (23) into an unpolarized and a polarized gluon contribution
$`G_{\mu +}(\sigma ^{}n)G_{\mu ^{}+}(\overline{z}+\sigma n)ϵ^\mu (k,\lambda )ϵ^\mu ^{}(k^{},\lambda )=`$ (26)
$`{\displaystyle \frac{1}{2}}G_{\mu +}(\sigma ^{}n)G_{\mu ^{}+}(\overline{z}+\sigma n)\left[g_{}^{\mu \mu ^{}}+\lambda 𝒫^{\mu \mu ^{}}\right],`$
where $`g_{}^{11}=g_{}^{22}=1`$ and $`𝒫^{12}=𝒫^{21}=i`$ while all other components of these tensors are zero. As for the case of quarks, see (11), we introduce a form factor decomposition for the proton matrix elements of the field strength tensors
$`{\displaystyle _0^1}𝑑xp^+{\displaystyle \frac{dz^{}}{2\pi }e^{ixp^+z^{}}_0^{\mathrm{}}𝑑\sigma 𝑑\sigma ^{}e^{\epsilon \sigma \epsilon ^{}\sigma ^{}}}`$ (27)
$`\times p^{},\nu ^{}|G_{\mu +}(\sigma ^{}n)G_{\mu ^{}+}(\overline{z}+\sigma n)|p,\nu g_{}^{\mu \mu ^{}}`$
$`={\displaystyle \frac{\overline{u}(p^{},\nu ^{})\gamma _+u(p,\nu )}{2p_+}}R_V^g(t)+{\displaystyle \frac{i}{2m}}{\displaystyle \frac{\overline{u}(p^{},\nu ^{})\sigma _{+\nu }\mathrm{\Delta }^\nu u(p,\nu )}{2p_+}}R_T^g(t).`$
and
$`{\displaystyle _0^1}𝑑xp^+{\displaystyle \frac{dz^{}}{2\pi }e^{ixp^+z^{}}_0^{\mathrm{}}𝑑\sigma 𝑑\sigma ^{}e^{\epsilon \sigma \epsilon ^{}\sigma ^{}}}`$ (28)
$`\times p^{},\nu ^{}|G_{\mu +}(\sigma ^{}n)G_{\mu ^{}+}(\overline{z}+\sigma n)|p,\nu 𝒫^{\mu \mu ^{}}`$
$`={\displaystyle \frac{\overline{u}(p^{},\nu ^{})\gamma _+\gamma _5u(p,\nu )}{2p_+}}R_A^g(t).`$
The form factors are related to SPDs at zero skewedness ($`_{\zeta =0}^g`$, $`𝒦_{\zeta =0}^g`$ and $`𝒢_{\zeta =0}^g`$), e.g.
$$R_V^g(t)=_0^1\frac{dx}{x^2}_{\zeta =0}^g(x,t).$$
(29)
Since the forward limits of $`_\zeta ^g`$ and $`𝒢_\zeta ^g`$ are defined in such a way that
$$xg(x)=_{\zeta =0}^g(x,t=0),x\mathrm{\Delta }g(x)=𝒢_{\zeta =0}^g(x,t=0),$$
(30)
one may still call these form factors $`1/x`$-moments of SPDs. Neglecting, as in the case of quarks, $`R_T^g`$ we finally arrive at the helicity amplitudes for the gluon contribution
$`_{0+,\mu +}^{V(g)}(s,t)`$ $`=`$ $`{\displaystyle \frac{e}{2}}[_{01,\mu 1}^{V(g)}(s,t)(R_V^{V(g)}(t)+R_A^{V(g)}(t))`$ (31)
$`+`$ $`_{01,\mu 1}^{V(g)}(s,t)(R_V^{V(g)}(t)R_A^{V(g)}(t))],`$
where
$$R_{V,A}^{Vg}=\underset{a}{}e_aB_a^VR_{V,A}^g.$$
(32)
Because of parity invariance $`_{01,\mu 1}^{V(g)}=(1)^\mu _{01,\mu 1}^{V(g)}`$ the form factor $`R_A^{V(g)}`$ does not contribute to the amplitudes for longitudinal photons.
The parton-level amplitudes, to be evaluated from the six lowest order Feynman graphs shown in Fig. 5, read
$$_{0+,\mu +}^{V(g)}=\frac{2\pi \alpha _s(\mu _R)}{N_c}f_V_0^1𝑑\tau \varphi _V(\tau )f_\mu ^{(g)}(\tau ),$$
(33)
where
$`f_+^{(g)}(\tau )`$ $`=`$ $`\sqrt{{\displaystyle \frac{t}{2su}}}{\displaystyle \frac{Q^2}{s+Q^2}}{\displaystyle \frac{1}{\tau \overline{\tau }}}{\displaystyle \frac{tQ^2(s+Q^2)^24\tau \overline{\tau }su}{(\overline{\tau }t\tau Q^2)(\tau t\overline{\tau }Q^2)}},`$
$`f_{}^{(g)}(\tau )`$ $`=`$ $`\sqrt{{\displaystyle \frac{ut}{2s}}}{\displaystyle \frac{s}{s+Q^2}}{\displaystyle \frac{14\tau \overline{\tau }}{\tau \overline{\tau }}}{\displaystyle \frac{Q^2}{(\overline{\tau }t\tau Q^2)(\tau t\overline{\tau }Q^2)}},`$
$`f_0^{(g)}(\tau )`$ $`=`$ $`{\displaystyle \frac{Q}{s+Q^2}}{\displaystyle \frac{1}{\tau \overline{\tau }}}{\displaystyle \frac{2\tau \overline{\tau }[Q^2(s+Q^2)t(sQ^2)]tQ^2}{(\overline{\tau }t\tau Q^2)(\tau t\overline{\tau }Q^2)}}.`$ (34)
Since the gluon amplitudes only contribute to photoproduction of flavor neutral vector mesons whose associated distribution amplitudes are symmetric under the interchange $`\tau \overline{\tau }`$, we display only the $`\tau \overline{\tau }`$ symmetric part of the amplitudes in (34). The quark amplitudes (18), on the other hand, contribute to the production of flavored mesons too. We, therefore, show the full quark amplitudes although we do not discuss these cases here.
As one may see from (34), the gluon amplitudes vanish in the limit $`Q^20`$. Finite $`\gamma gVg`$ amplitudes may be obtained if meson masses and/or transverse momenta are taken into account (cf. for instance ).
The amplitudes (21) have to be added to those given in (12) for vector mesons:
$$_{0\nu ^{},\mu \nu }^V=_{0\nu ^{},\mu \nu }^{V(q)}+_{0\nu ^{},\mu \nu }^{V(g)}.$$
(35)
## 3 The form factors and the meson distribution amplitudes
Before we present numerical results for the observables of electroproduction of mesons, we have to model the new form factors. In Eq. (13) the general composition of these new form factors is presented in terms of the individual flavor contributions. The explicit flavor structure of the form factors for various mesons is given in Tab. 1. In contrast to Compton scattering where the sum runs over all flavors, here the sum is over the valence quarks of the produced mesons. In other words, the meson selects its valence quarks from the proton. The physical situation is thus similar to DVEM in this respect. For vector mesons, also the flavor factors associated with the gluonic form factors (see (31), (32)) are listed in the table.
$`\omega \varphi `$ mixing is ignored since the corresponding mixing angle is very small. $`\eta \eta ^{}`$ mixing, on the other hand, is taken into account. Following , we work in the quark flavor basis and write
$`\eta `$ $`=`$ $`\mathrm{cos}\varphi _P\eta _q\mathrm{sin}\varphi _P\eta _s,`$
$`\eta ^{}`$ $`=`$ $`\mathrm{sin}\varphi _P\eta _q+\mathrm{cos}\varphi _P\eta _s.`$ (36)
$`\eta _q`$ is a state built from $`u`$ and $`d`$ quarks only while $`\eta _s`$ is a $`s\overline{s}`$ state. The parameters of that $`\eta \eta ^{}`$ mixing scheme, the mixing angle, $`\varphi _P`$, and the decay constants, $`f_q`$ and $`f_s`$, of the basis states are determined in on exploiting the divergencies of the axial vector currents which embody the axial vector anomaly. For the mixing angle a value of $`39.2^{}`$ is found in .
As an inspection of Eqs. (10) and (11) reveals the form factors $`R_{V,A}^{Va}`$ are exactly the same as those appearing in Compton scattering . Thus, in principle, from a combined analysis of data on Compton scattering and production cross section for various mesons, one may extract information on the form factors for individual flavors from experiment. This allows to test the soft mechanism independent of a specific model for the form factors.
For a numerical estimate of the form factors we use the model proposed in Ref. . In a frame where $`\mathrm{\Delta }^+=0`$ the SPDs and, hence, the form factors can be represented as overlaps of light-cone wave functions summed over all Fock states in close analogy to the familiar Drell-Yan formula . A detailed discussion of that overlap representation is given in Refs. . Each $`N`$-particle Fock state is described by a number of terms, each with its own momentum space wave function $`\mathrm{\Psi }_{N\beta }`$, where $`\beta `$ labels different spin-flavor combinations of the $`N`$ partons. Assuming a single Gaussian $`\stackrel{~}{k}_i`$-dependence of the soft Fock state wave functions
$$\mathrm{\Psi }_{N\beta }(x_i,\stackrel{~}{k}_i)\mathrm{exp}[a_N^2\underset{i=1}{\overset{N}{}}\frac{\stackrel{~}{k}_i^2}{x_i}],$$
(37)
one can explicitly carry out the momentum integration in the overlap formula. The ansatz (37) satisfies various theoretical requirements and is in line with our central hypothesis that the soft hadronic wave functions are dominated by transverse momenta with $`\stackrel{~}{k}_i^2/x_i\mathrm{\Lambda }^2`$, necessary to achieve the factorization of the electroproduction amplitudes into soft and hard parts <sup>2</sup><sup>2</sup>2 The wave function (37), perhaps multiplied by a polynomial in the $`x_i`$, is not continuous in the end-points $`x_i=\stackrel{~}{k}_i=0`$ ($`i=1,2`$ or 3). It can, however, be shown that the overlaps evaluated from such wave functions, are infrared stable, i.e. they are not dominated by contributions from regions of very small $`x_i`$ and $`\stackrel{~}{k}_i`$.. The results of the transverse momentum integration for the vector and axial vector form factors are respectively related with the Fock state contributions to the unpolarized ($`q_a(x)`$, $`g(x)`$) and polarized ($`\mathrm{\Delta }q_a(x)`$, $`\mathrm{\Delta }g(x)`$) parton distributions. For simplicity one may assume a common transverse size parameter $`a_N=\widehat{a}`$ for all Fock states which seems to be a reasonable approximation since, for large -t, the main contribution to the overlap integral is only due to a limited number of Fock states . This simplification immediately allows one to sum over the Fock states without specifying the $`x_i`$-dependence of the wave functions. One then arrives at the following model for the form factors for individual flavors ($`a=u,d,s`$):
$`R_V^{Va}(t)`$ $`=`$ $`{\displaystyle _0^1}{\displaystyle \frac{dx}{x}}\mathrm{exp}\left[{\displaystyle \frac{1}{2}}\widehat{a}^2t{\displaystyle \frac{1x}{x}}\right]\{q_a(x)+\overline{q}_a(x)\},`$
$`R_V^{Pa}(t)`$ $`=`$ $`{\displaystyle _0^1}{\displaystyle \frac{dx}{x}}\mathrm{exp}\left[{\displaystyle \frac{1}{2}}\widehat{a}^2t{\displaystyle \frac{1x}{x}}\right]\{q_a(x)\overline{q}_a(x)\},`$
$`R_V^g(t)`$ $`=`$ $`{\displaystyle _0^1}{\displaystyle \frac{dx}{x}}\mathrm{exp}\left[{\displaystyle \frac{1}{2}}\widehat{a}^2t{\displaystyle \frac{1x}{x}}\right]g(x).`$ (38)
The corresponding axial vector form factors are obtained from (38) by replacing the unpolarized parton distributions $`q_a`$, $`g`$ with the polarized ones, $`\mathrm{\Delta }q_a`$, $`\mathrm{\Delta }g`$.
As shown in an evaluation of these form factors from the parton distributions of Glück et al. (GRV) (taken at a scale of 1 GeV) and with $`\widehat{a}1\mathrm{GeV}^1`$, leads to results for Compton scattering in fair agreement with experiment. In order to improve the model (38) the lowest three Fock states were modeled explicitly in assuming specific distribution amplitudes, e.g.
$$\varphi _{123}(x_i)=60x_1x_2x_3(1+3x_1),$$
(39)
for the valence Fock state . The form factors (38) for the quarks are then evaluated from these three lowest Fock states (with $`a_3=a_4=a_5=0.75\mathrm{GeV}^1`$) and the contribution from all higher Fock states is estimated by setting ($`a_N=1.3a_3`$ for $`N>5`$)
$$\underset{N>5}{}q_a^{(N)}(x)=q_a(x)\underset{N=3,4,5}{}q_a^{(N)}(x).$$
(40)
The $`q_a(x)`$ are taken from the GRV parameterization and the three lowest Fock state contribution $`q_a^{(N=3,4,5)}(x)`$ are evaluated from the light-cone wave functions. This model provides a good fit to Compton scattering and to the proton form factor $`F_1`$ (by expressions similar to (38) ). In Fig. 6 numerical results for the vector and axial vector form factors are shown. The simplifying model assumption of unpolarized gluons and sea quarks has the consequence of a zero gluon form factor $`R_A^g`$. Most of the strange form factors are therefore zero too:
$$R_V^{Ps}=R_A^{Vs}=R_A^{Ps}=0.$$
(41)
Only $`R_V^{Vs}`$ is non-zero, even though very small. The form factors for $`u`$ quarks are largest.
They approximately behave as $`1/t^2`$ in the momentum transfer region from about 5 to $`15\mathrm{GeV}^2`$; for the other flavors and for the gluon that range is shifted to somewhat smaller values of $`t`$. With increasing $`t`$ the form factors for $`u`$ and $`d`$ quarks gradually turn into the soft physics asymptotic $`1/t^4`$, while the other form factors decrease faster. The leading powers of $`1/t`$ in the asymptotic behavior of the form factors (38) follow from the $`x_i`$-dependence of the model wave functions at the end points (see for instance Eq. (39)) . In the region where the form factors drop as $`1/t^4`$ or faster, which is above 100 GeV<sup>2</sup>, the perturbative contribution will take the lead.
The other soft physics information required in our approach is that of the form of the meson’s distribution amplitude. From analyzes of the pion-photon transition form factor (see for instance ) it became evident that the pion’s distribution amplitude (its formal definition is given in (15)) is close to the asymptotic form
$$\varphi _{\mathrm{AS}}(\tau )=6\tau (1\tau ).$$
(42)
This result is supported by the instanton model and by recent QCD sum rule studies . The analyzes of the $`\eta `$\- and $`\eta ^{}`$-photon transition form factors revealed that the $`\eta _q`$ distribution amplitude is close to the form (42), too . Although the transition form factor data are compatible with the asymptotic distribution amplitude for the $`\eta _s`$ as well, a somewhat narrower one cannot be excluded. For vector mesons no phenomenological information is available but QCD sum rules taught us that the distribution amplitudes for longitudinally polarized vector mesons can also be approximated by (42). In order to keep matters simple we therefore choose the form (42) for all mesons. We expect that the uncertainties in the predicted production cross sections due this choice do not exceed $`1015\%`$. Associated with the distribution amplitude (42) is a value of 3 for the $`1/\tau `$-moment (20).
For the meson decay constants we use the values
$`f_\pi `$ $`=`$ $`132\mathrm{MeV},f_\rho =216\mathrm{MeV},`$
$`f_\omega `$ $`=`$ $`195\mathrm{MeV},f_\varphi =237\mathrm{MeV},`$ (43)
and for the decay constants of the states $`\eta _q`$ and $`\eta _s`$
$$f_q=141\mathrm{MeV},f_s=177\mathrm{MeV}.$$
(44)
## 4 Photoproduction of mesons
Using (12) and (19), we obtain for the photoproduction cross section of pseudoscalar mesons
$`{\displaystyle \frac{d\sigma }{dt}}^P`$ $`=`$ $`{\displaystyle \frac{1}{2}}\alpha _{\mathrm{em}}\left[\pi \alpha _s(\mu _R)f_P1/\tau _P{\displaystyle \frac{C_F}{N_C}}\right]^2`$ (45)
$`\times `$ $`{\displaystyle \frac{t}{u^2s^4}}\left\{(su)^2(R_V^P(t))^2+t^2(R_A^P(t))^2\right\}.`$
The corresponding expressions for photoproduction of vector mesons are a bit more complicated due to the occurrence of the gluonic contribution. In Fig. 7 we show, as an typical example, the soft physics contribution to the large momentum transfer photoproduction cross section of uncharged pions as evaluated from the form factors discussed in Sect. 3. We see that at fixed scattering angle the cross section approximately exhibits the $`s^7`$-scaling as predicted by dimensional counting . This scaling behavior holds in the soft physics approach as long as the form factors $`R_{V,A}^M`$ behave as $`1/t^2`$ (see Fig. 6). Similar results are found for the production of the other mesons.
As compared to experiment (at $`s10\mathrm{GeV}^2`$) the soft physics contributions are too small by orders of magnitude. This can easily be understood by evaluating the ratio of $`\pi ^0`$ production and Compton cross section
$$\frac{d\sigma (\gamma p\pi ^0p)}{d\sigma (\gamma p\gamma p)}=\frac{t}{s}\frac{\alpha _s^2(\mu _R)}{\alpha _{\mathrm{em}}}\frac{f_\pi ^2<1/\tau >_\pi ^2}{s}c_{soft},$$
(46)
where $`c_{soft}`$, a ratio of form factors and kinematical factors, is of order 1. The ratio of the two cross sections is therefore about $`2\mathrm{GeV}^2/s`$, i.e. much smaller than unity in contradiction to experiment , where the ratio is about 50 (at $`s10\mathrm{GeV}^2`$, and a $`\gamma ^{}p`$ c.m. scattering angle, $`\theta `$, of $`90^{}`$). Mainly responsible for the small ratio (46) is the perturbative formation of the meson which only probes small quark-antiquark separations in the meson. The amplitude is, therefore, proportional to the meson’s decay constant which, for dimensional reasons, is to be scaled by $`\sqrt{s}`$. The $`1/\tau _\pi `$ moment, appearing as a consequence of the perturbative meson formation, cannot compensate the small ratio $`f_\pi /\sqrt{s}`$.
The ratio (46) also holds in perturbative calculations, in the pure quark picture as well as in the diquark model, a variant of the standard perturbative approach in which diquarks are considered as quasi-elementary constituents of baryons . The factor $`c_{pert}`$ may, however, be larger than unity. Although there is no obvious enhancement in any of the many Feynman graphs contributing to the perturbative amplitude, the graphs may conspire in such a way that a large value of $`c_{pert}`$ is built up. In order to see whether or not this is the case, an explicit and reliable calculation of meson production within the perturbative approach is called for.
The observation of a relatively large photoproduction cross section in experiment is in line with the power law behavior in $`s`$ at fixed scattering angle. $`s^n`$-fits to the present data in the range $`6.5\mathrm{GeV}^2\stackrel{<}{}s\stackrel{<}{}12\mathrm{GeV}^2`$ and $`50^{}<\theta <130^{}`$ provide :
$`\gamma p\pi ^0p:`$ $`n=8.0\pm 0.1,`$
$`\gamma p\gamma p:`$ $`n=6.1\pm 0.3.`$
For the $`\gamma p(\rho ^0+\omega )p`$ data the statistics does not allow a meaningful determination of the power $`n`$. However, $`n`$ seems to be larger than 7, rather compatible with 8. For Compton scattering the power is compatible with dimensional counting, while for $`\pi ^0`$ production, and possibly for the sum of $`\rho ^0`$ and $`\omega `$ production, the value of $`n`$ rather equals that one observed in elastic $`\pi p`$ scattering ($`n8`$; data are averaged over resonance-like structures) <sup>3</sup><sup>3</sup>3 We recall that a $`s^8`$-scaling of the fixed-angle meson baryon cross section is easily accounted for by the soft physics approach in the relevant region of energy .. Admittedly, the photoproduction data are rather poor and need confirmation. Data on photoproduction of $`\varphi `$ mesons will become available from the TJlab soon which will perhaps allow a determination ot the power $`n`$ for that reaction.
Both the observations in the large momentum transfer photoproduction data, the large powers of $`s`$ at fixed scattering angle and the large cross sections, indicates that another dynamical mechanism is at work here. It is tempting to assign it to the hadronic component of the photon. This proposition is supported by a vector meson dominance (VMD) estimate of the photoproduction cross section. Combined with quark model ideas VMD, for instance, relates photoproduction of $`\rho ^0`$-mesons to elastic pion-nucleon scattering
$$\frac{d\sigma }{dt}(\gamma p\rho ^0p)=\alpha _{\mathrm{em}}\frac{\pi f_\rho ^2}{m_\rho ^2}\left[\frac{d\sigma }{dt}(\pi ^+p\pi ^+p)+\frac{d\sigma }{dt}(\pi ^{}p\pi ^{}p)\right].$$
(47)
This relation is satisfied by experiment within a factor of 2-3 . With respect to the uncertainties arising from possible spin effects and the poor quality of data this may be considered as fair agreement. Thus, it seems that photoproduction of $`\rho ^0`$ and $`\pi ^0`$ – and likely of other mesons – is indeed controlled by the hadronic component of the photon. In this case one would expect the produced vector mesons to be polarized transversally rather than longitudinally. Since the fixed-angle energy dependencies of the contributions from the hadronic component of the photon ($`s^8`$) and from the soft mechanism ($`s^7`$) are so close, much higher energies are needed before the soft contribution (and/or the perturbative one) will control photoproduction of mesons. We, therefore, refrain from presenting more predictions from the soft physics approach for photoproduction of mesons. We stress that our approach to photoproduction requires high energies, large momentum transfer and small values of $`|\mathrm{cos}\theta |`$. If $`s/t1`$ Pomeron exchange becomes dominant, see for instance .
One may wonder whether Compton scattering is also dominated by the hadronic component of the photon. However, the analogous VMD estimate, with both the photons replaced by vector mesons, provides values for the Compton cross section that are about an order of magnitude below experiment at $`s10\mathrm{GeV}^2`$. Moreover, the Compton cross section exhibits $`s^6`$-scaling and not an $`s^8`$ one as would be the case if the hadronic component of the photon dominates. Thus, the simplest elementary process, elastic scattering of point-like photons from quarks, dominates Compton scattering off protons already at rather low energies .
## 5 Electroproduction of mesons
Let us now discuss our results for large momentum transfer electroproduction. As is well-known the cross section for $`epepM`$ can be decomposed as follows
$`{\displaystyle \frac{d^4\sigma ^M}{dsdQ^2dtd\phi }}`$ $`=`$ $`{\displaystyle \frac{\alpha _{\mathrm{em}}s}{16\pi ^2E_L^2m^2Q^2(1ϵ)}}`$ (48)
$`\times `$ $`\left({\displaystyle \frac{d\sigma _T^M}{dt}}+ϵ{\displaystyle \frac{d\sigma _L^M}{dt}}+2ϵ\mathrm{cos}2\phi {\displaystyle \frac{d\sigma _{TT}^M}{dt}}+\sqrt{2ϵ(1+ϵ)}\mathrm{cos}\phi {\displaystyle \frac{d\sigma _{LT}^M}{dt}}\right),`$
where $`\phi `$ denotes the azimuthal angle between the hadronic and leptonic scattering planes. $`E_L`$ is the energy of the incoming electron in the laboratory frame and $`ϵ`$ is the ratio of longitudinal to transverse photon flux. Details of the kinematics can be found for instance in Ref. . The partial cross sections in (48) read:
(i) The cross sections for transverse photons (reducing to the unpolarized cross section for photoproduction of mesons, i.e. for $`Q^2=0`$) and for longitudinal photons,
$`{\displaystyle \frac{d\sigma _T^M}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{32\pi s(s+Q^2)}}{\displaystyle \underset{\nu ^{},\nu }{}}|_{0\nu ^{},+\nu }^M|^2,`$
$`{\displaystyle \frac{d\sigma _L^M}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{32\pi s(s+Q^2)}}{\displaystyle \underset{\nu ^{},\nu }{}}|_{0\nu ^{},\mathrm{\hspace{0.17em}0}\nu }^M|^2.`$ (49)
(ii) The transverse-transverse and longitudinal-transverse interference terms
$`{\displaystyle \frac{d\sigma _{TT}^M}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{64\pi s(s+Q^2)}}\mathrm{Re}{\displaystyle \underset{\nu ^{},\nu }{}}_{0\nu ^{},+\nu }^M_{0\nu ^{},\nu }^M,`$
$`{\displaystyle \frac{d\sigma _{LT}^M}{dt}}`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}}{64\pi s(s+Q^2)}}\mathrm{Re}{\displaystyle \underset{\nu ^{},\nu }{}}_{0\nu ^{},\mathrm{\hspace{0.17em}0}\nu }^M\left[_{0\nu ^{},+\nu }^M_{0\nu ^{},\nu }^M\right].`$ (50)
In Fig. 8 we show, as a typical example, the $`\rho ^0`$ production cross section for longitudinally and transversally polarized photon as a function of $`Q^2`$ at a scattering angle of $`90^{}`$. Except for small $`Q^2`$ ($`Q^2/s<<1`$) the longitudinal cross section
$`{\displaystyle \frac{d\sigma _L^M}{dt}}`$ $`=`$ $`{\displaystyle \frac{\alpha _{\mathrm{em}}}{4N_c^2}}{\displaystyle \frac{[\pi \alpha _s(\mu _R)f_MC_F]^2}{s(s+Q^2)}}\{{\displaystyle _0^1}d\tau \varphi _M(\tau )`$ (51)
$`\times `$ $`[(1\kappa _M)[f_0^{(q)}(\tau )R_V^M(t)+{\displaystyle \frac{1}{C_F}}f_0^{(g)}(\tau )R_V^{Mg}(t)]`$
$`+`$ $`(1+\kappa _M)f_0^{(q)}(\tau )R_A^M(t)]\}^2`$
dominates, i.e. the cross section that conserves $`s`$-channel helicity. Again we have similarity to DVEM. Also similar is the fact that the longitudinal cross section for the production of vector mesons is associated with the vector form factors, $`R_V^V`$, while in the case of pseudoscalar mesons it is connected with the axial vector ones, $`R_A^P`$. In contrast to the limiting cases of either $`t0`$ or $`Q^20`$ the explicit form of the mesons distribution amplitude is required in the evaluation of the large momentum transfer electroproduction cross section. Tuning the ratio $`Q^2/t`$ details of the distribution amplitude can be explored. Even if the form factors behave $`1/t^2`$ strictly, $`d\sigma /dt`$ does not exhibit $`s^7`$-scaling. Despite of this we keep multiplying the cross section by $`s^7`$ since this compensates most of the energy dependence.
The hadronic component of the photon vanishes with increasing $`Q^2`$ rapidly, approximately as $`m_V^2/(Q^2+m_V^2)`$ in the amplitude. This is for instance, clearly visible in the integrated cross section for $`\rho ^0`$ production: while, at low $`Q^2(\stackrel{<}{}2\mathrm{GeV}^2)`$, its energy dependence is very similar to that of total cross sections for elastic hadron-hadron scattering, it is much steeper at large $`Q^2`$ . The steep rise of the $`\rho ^0`$ cross section is correlated with the behaviour of the gluon SPD for $`0<x^{}x1`$ which should reflect the strong increase of the gluon distribution for $`x0`$ . By virtue of the rapidly decreasing hadronic component of the photon and the strong rise of the longitudinal cross section (see Fig. 8) with increasing $`Q^2`$, we expect the soft physics approach to be applicable for photon virtualities larger than about $`23\mathrm{GeV}^2`$ provided $`s`$, $`t`$ and $`u`$ are large.
Results for the partial cross sections for electroproduction of $`\rho ^0`$-mesons are shown in Fig. 9 (at $`s=20\mathrm{GeV}^2`$ and $`40\mathrm{GeV}^2`$). We see that the longitudinal cross section is dominant in the region of small $`|\mathrm{cos}\theta |`$. For larger values of $`|\mathrm{cos}\theta |`$ the transverse cross section as well as the longitudinal-transverse interference become sizeable.
In order to demonstrate the relative magnitude of the production cross sections for various mesons we display predictions for the scaled longitudinal cross sections at $`s=20\mathrm{GeV}^2`$ and $`Q^2=3\mathrm{GeV}^2`$ in Fig. 10. The gluonic contributions to the cross sections for vector mesons are generally small in the kinematical region of interest. This is obvious from the relative strength of the quark and gluon form factors, see Fig. 6. The exceptional case is the $`\varphi `$ meson. The only quark form factor contributing, $`R_V^{Vs}`$, is very small and the gluonic contribution therefore dominates the production of $`\varphi `$ mesons. Since the form factor $`R_A^{Ps}`$ is zero in the model (see Eq. (41)) the ratio of longitudinal cross sections for the production of $`\eta ^{}`$ and $`\eta `$ mesons is given by the square of the tangent of the pseudoscalar meson mixing angle, see (36).
Finally we comment on the accuracy of our calculation. As we said repeatedly, large momentum transfer and photon virtualities larger than about $`23\mathrm{GeV}^2`$ are required for the dominance of the soft mechanism. However, the momentum transfer should not be too large since then the pure perturbative contribution will become important. The onset of the perturbative regime is expected to be beyond $`t100\mathrm{GeV}^2`$ . Even in the soft regime one has to be aware of corrections. For instance, the lowest-order, leading twist perturbative formation of the mesons may be subject to substantial corrections of perturbative and/or soft origin. These may give rise to a $`\kappa `$-factor in the normalization of the electroproduction cross section as is known from the Drell-Yan process and, according to Martin et al. , is required in DVEM at least at small $`Q^2/s`$.
## 6 Summary
In the present work electroproduction of flavor neutral pseudoscalar and longitudinally polarized vector mesons at large $`s`$, $`t`$ and $`u`$ is investigated. The photon virtuality is not considered as a large scale; therefore the limit $`Q^20`$ is included in the investigation. This study of electroproduction is complementary to the case of DVEM where $`Q^2`$ is large and $`t`$ small. Based on the central assumption of the dominance of small parton virtualities and small intrinsic transverse momenta in the proton’s light-cone wave function we have shown that, like in Compton scattering , the electroproduction amplitudes factorize in hard parton-level subprocess amplitudes and soft proton matrix elements described by the same type of form factors as appear in Compton scattering. These form factors represent $`1/x`$-moments of SPDs. The soft mechanism bears resemblance to the dynamics controlling DVEM in many respects: The same parton-level subprocesses occur, the longitudinal cross section dominates (if $`|\mathrm{cos}\theta |`$ is small and $`Q^2`$ not too small) and the soft information on the proton is encoded in SPDs. Different is that, in the large momentum transfer region, a symmetric frame with zero skewedness can be chosen which entails the formation of $`1/x`$-moments of the SPDs, i.e. the appearance of new form factors. For asymptotically large momentum transfer the perturbative contribution will take the lead, the soft contribution, discussed here, then presents a power correction to it. We emphasize that the dimensional counting rule behaviour, i.e. $`s^7`$-scaling, approximately holds for photoproduction of mesons in the soft physics approach for a limited range of energy.
The new form factors, characteristic of large momentum transfer Compton scattering and electroproduction of mesons, can in principle be extracted from experiment by Rosenbluth-type separations . Their measurements would provide information on the large momentum transfer behavior of the proton SPDs and would allow to test models for them. Moreover, the experimental verification of their energy independence would constitute a severe test of the soft physics approach.
Based on a light-cone wave function overlap model for the form factors we have presented detailed predictions for electroproduction of pseudoscalar and longitudinally polarized vector mesons at moderately large photon virtuality. Although the soft physics approach also applies to large momentum transfer photoproduction of mesons it seems - as judged on the basis of the present data - that the contributions from the hadronic component of the photon dominate these reactions up to rather high energies. The kinematical region in which the soft physics approach is applicable to electroproduction, is accessible to experiments at the upgraded TJlab and at the proposed ELFE accelerator and EPIC collider. The measurement of large momentum transfer electroproduction of mesons is certainly difficult but seems feasible. We have not discussed electroproduction of flavored mesons and of transversally polarized vector mesons in this work. These processes involve flavor or helicity non-diagonal SPDs which are not directly related to those appearing in the processes we have investigated. As in DVEM , higher twist dynamics plays an important role in electroproduction of transversally polarized vector mesons since the leading twist, lowest order subprocess amplitudes are zero in this case.
## 7 Acknowledgment
We would like to thank Markus Diehl, Thorsten Feldmann and Rainer Jakob for many stimulating discussions. H.W. Huang thanks the Deutsche Forschungsgemeinschaft for support.
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# 1. Examples
## 1. Examples
As was probably first noticed by L. Ein and N. Shepherd-Barron (see \[ES\]), many examples of homaloidal polynomials arise from the theory of prehomogeneous vector spaces. Recall that a complex vector space $`V`$ is called prehomogeneous with respect to a linear rational representation of an algebraic group $`G`$ in $`V`$ if there exists a non-constant polynomial $`F`$ such that the complement of its set of zeroes is homogeneous with respect to $`G`$. The polynomial $`F`$ is necessarily homogeneous and an eigenvector for $`G`$ with some character $`\chi :G\text{GL}(1)`$. It generates the algebra of invariants for the group $`G_0=Ker(\chi )`$. The reduced part $`F_{red}`$ of $`F`$ (i.e. the product of irreducible factors of $`F`$) is determined uniquely up to a scalar multiple. A prehomogeneous space is called regular if the determinant of the Hesian matrix of $`F`$ is not identically zero. This definition does not depend on the choice of $`F`$. We shall call $`F`$ a relative invariant of $`V`$. Note that there is a complete classification of regular irreducible prehomogeneous spaces with respect to a reductive group $`G`$ (see \[KS\]).
###### Theorem (\[ES, EKP\])\]) 1
Let $`V`$ be a regular prehomogeneous vector space. Then its relative invariant is a homaloidal polynomial.
Here are some examples: Examples 1-4. 1. Any non-degenerate quadratic form $`Q`$ is obviously a homaloidal polynomial. The corresponding birational map is a projective automorphism. It is also a relative invariant for the orthogonal group $`O(Q)\times \text{GL}(1)`$ in its natural linear representation.
2. A reduced cubic polynomial $`F`$ on $`V`$ is a relative invariant for a regular prehomogeneous space with respect to a reductive group $`G`$ if and only if the pair $`(V,G)`$ is one of the following (up to a linear transformation):
2.1: $`G=\text{GL}(1)^3\text{GL}(3),V=^3`$, the action is natural, $`F=x_0x_1x_2`$.
2.2: $`G=\text{GL}(3)`$, $`V`$ is the space of quadratic forms on $`^3`$, the action is via the natural action on $`^3`$, $`F`$ is the discriminant function.
2.3: $`G=\text{GL}(3)\times \text{GL}(3)`$, $`V=\text{Mat}_3`$ is the space of complex $`3\times 3`$-matrices, the action is by $`(g,g^{})A=gAg^{}^1`$, the polynomial $`F`$ is the determinant.
2.4: $`G=\text{GL}(6),V=\mathrm{\Lambda }^2(^6)`$, the action is via the natural action on $`^6`$. The polynomial $`F`$ is the pffafian polynomial.
2.5: $`G=E_6\times \text{GL}(1)`$, $`V=^{27}=\text{Mat}_3\times \text{Mat}_3\times \text{Mat}_3`$ is its irreducible representation of minimal dimension. The polynomial $`F`$ is the Cartan cubic $`F(A,B,C)=|A|+|B|+|C|\text{Tr}(ABC)`$.
The last four examples correspond to the four Severi varieties: nonsingular nondegenerate subvarieties $`S`$ of $`^r`$ of dimension $`\frac{2r4}{3}`$ whose secant variety $`\text{Sec}(S)`$ is not equal to the whole space. The zero locus of the cubic $`F`$ in $`(V)`$ defines the secant variety. The singular locus of $`\text{Sec}(S)`$ is the Severi variety. According to a theorem from \[ES\], any homaloidal cubic polynomial $`F`$ such that the singular locus of $`F^1(0)`$ in $`(V)`$ is nonsingular coincides with one from examples 2.2-2.5.
3. Let us identify $`^{n^21}`$ with the space $`(\text{Mat}_n)`$. The map $`AA^1`$ is obbviously birational. It is given by the polar linear system of the polynomial $`Adet(A)`$. The polynomial is a relative invariant from Example 2.3 (extended to any dimension).
4. The polynomial $`F=x_0(x_0x_2+x_1^2)`$ is homaloidal. It is a relative invariant for a prehomogeneous space with respect to a non-reductive group.
## 2. Multiplicative Legendre transform
This section is almost entirely borrowed from \[EKP\]. Let $`F\text{Pol}_d(V)`$ be a homogeneous polynomial of degree $`d`$ on a complex vector space $`V`$ of dimension $`n+1`$. We denote by $`F^{}`$ or by $`dF`$ the derivative map $`VV^{},v(dF)_v`$. If no confusion arises we also use this notation for the associated rational map $`(V)(V^{})`$. If we choose a basis in $`V`$ and the corresponding dual basis in $`V^{}`$, we will be able to identify both spaces with $`^n`$, and the map $`F^{}`$ with the polar map defined in the introduction. Suppose $`F`$ is homaloidal, i.e. $`F^{}`$ defines a birational map $`(V)(V^{})`$. Then, obviously, $`d\mathrm{ln}F=F^{}/F`$ defines a birational map $`VV^{}`$.
###### Lemma 1
Let $`f`$ be a homogeneous function of degree $`k`$ on $`V`$ (defined on an open subset) such that $`det(\text{Hess}(\mathrm{ln}f))`$ is not identical zero. Then there exists a homogeneous function $`f_{}`$ on $`V^{}`$ of degree $`k`$ such that on some open subset of $`V`$
$$f_{}(d\mathrm{ln}f)=1/f.$$
$`(2.1)`$
Proof. Recall first the definition of the Legendre transform. Let $`Q`$ be a function on $`V`$ defined in an open neighborhood of a point $`v_0`$ such that $`det\text{Hess}(Q)(v_0)0`$. Let $`dQ(v_0)=p_0V^{}`$. Then the Legendre transform $`L(Q)`$ of $`Q`$ is the function $`L(Q)`$ on $`V^{}`$ defined in an neighborhood of $`p_0`$ such that
$$L(Q)(p)=p(v_p)Q(v_p),$$
$`(2.2)`$
where $`v_p`$ is the unique critical point of the function $`vp(v)Q(v)`$ in a neighborhood of $`v_0`$.
Since the critical point $`v_p`$ satisfies $`p=dQ(v_p)`$, we obtain from (2.2) the equality of functions on an neighborhood of $`v_p`$ in $`V`$
$$L(Q)(dQ(v))=dQ(v)(v)Q(v).$$
Now let us apply this to $`Q=\mathrm{ln}f`$. We have
$$L(\mathrm{ln}f)(d\mathrm{ln}f(v))=d\mathrm{ln}f(v)v\mathrm{ln}f(v).$$
Recall that a homogeneous function $`H`$ of degree $`k`$ satisfies the Euler formula:
$$kH(v)=dH(v).$$
Applying this to $`H=\mathrm{ln}f`$, we get
$$e^{L(\mathrm{ln}f)k}(d\mathrm{ln}f)=1/f.$$
It remains to define $`f_{}`$ by
$$\mathrm{ln}f_{}=L(\mathrm{ln}f)k.$$
$`(2.3)`$
It is immediately checked that it is homogeneous of degree $`k`$.
The function $`f_{}`$ is called the multiplicative Legendre transform of $`f`$.
###### Theorem 2 (\[EKP\])
Let $`F\text{Pol}_d(V)`$ such that $`det\text{Hess}(\mathrm{ln}F)`$ is not identical zero. Then $`F`$ is homaloidal if and only if its multiplicative transform $`F_{}`$ is a rational function. Moreover, in this case
$$d\mathrm{ln}F_{}=(d\mathrm{ln}F)^1.$$
$`(2.4)`$
Proof. Suppose $`F`$ is homaloidal. Then $`d\mathrm{ln}F`$ is a rational map of topological degree 1 in its set of definition. It follows from the definition of the Legendre transform that $`L(\mathrm{ln}F)`$ is one-valued on its set of definition. Differentiating (2.1) we obtain $`(d\mathrm{ln}F_{})(d\mathrm{log}F)=\text{id}.`$ This checks (2.4). Since $`d\mathrm{ln}F_{}=dF_{}/F`$ is a homogeneous rational function, the function $`F_{}`$ must be rational. Conversely, if $`F_{}`$ is rational, we get (2.4) locally, by differentiating (1). Since $`d\mathrm{ln}F_{}`$ is rational, we have (2.4) globally, and hence $`d\mathrm{ln}F`$ is invertible. This implies that $`dF`$ defines a birational map, and hence $`F`$ is homaloidal.
###### Corollary 1
Let $`F(x_0,\mathrm{},x_n)`$ be a homaloidal polynomial of degree $`k>2`$. Assume $`F_{}`$ is a reduced polynomial. Then
$$k|2(n+1).$$
Proof. By the previous theorem
$$dF_{}dF=F^{k1}(x)F_{}(x)(x_0,\mathrm{},x_n).$$
This implies that the image of the hypersurface $`F=0`$ under the birational map $`dF:^n^n`$ is contained in the set of base points of the polar linear system of $`F_{}`$. Since $`F_{}`$ is reduced the latter is a closed subset of codimension $`>1`$. Thus $`F=0`$ is contained in the set of critical points of $`dF`$ (considered as a map of vector spaces) and hence $`F`$ divides the Hessian determinant. The assertion follows from this.
A natural question posed in \[EKP\] is the following: For which homogenous polynomials $`F`$ its multiplicative Legendre transform $`F_{}`$ is a polynomial function?
A polynomial with this property will be called a homaloidal EKP-polynomial. It is easy to see that $`F_{}`$ has the same degree as $`F`$ and $`(F_{})_{}=F`$. It is conjectured that any homaloidal EKP-polynomial is a relative invariant of a regular prehomogeneous space. The converse is proved in \[EKP\]. In this case $`F_{}=F`$, up to a scaling.
A remarkable result of \[EKP\] is the following:
###### Theorem 3
A homaloidal EKP-polynomial of degree 3 coincides with one from Examples 2.
Example 5. Consider the polynomial $`F`$ from Example 4. We have
$$d\mathrm{ln}F=(\frac{2x_0x_2+x_1^2}{x_0(x_0x_2+x_1^2)},\frac{2x_1}{x_0x_2+x_1^2},\frac{x_0}{x_0x_2+x_1^2}).$$
Inverting this map we obtain
$$(d\mathrm{ln}F)^1=(\frac{8x_2}{4x_0x_2+x_1^2},\frac{4x_1}{4x_0x_2+x_1^2},\frac{4x_0x_2x_1^2}{(4x_0x_2+x_1^2)x_2})=d\mathrm{ln}\frac{(4x_0x_2+x_1^2)^2}{x_2}.$$
Thus the multiplicative Legendre transform of $`F`$ equals
$$F_{}=\frac{(4x_0x_2+x_1^2)^2}{x_2}.$$
It is a homogeneous rational but not polynomial function.
## 3. Plane polar Cremona transformations
Here we shall classify all homaloidal polynomials in three variables with no multiple factors.
Since the set of common zeroes of the polars $`_iF`$ is equal to the set of non-smooth points of the subscheme $`V(F)`$, this is equivalent to requiring that the polars $`_iF`$ have no common factors, i.e. the linear system $`𝒫_F`$ has no fixed part.
Let $`f:^2^2`$ be a rational map defined by homogeneous polynomials $`(P_0,P_1,P_2)`$ of degree $`d`$ without common factors. Let $`𝒥(f)k[x_0,x_1,x_2]`$ be the ideal generated by the polynomials $`P_0,P_1,P_2`$. The corresponding closed subscheme $`B_f=V(𝒥(f))`$ of $`^2`$ is the base locus subscheme of the linear system spanned by $`P_0,P_1,P_2`$. The quotient sheaf $`𝒪_^2/𝒥(f)`$ is artinian and we denote by $`\stackrel{~}{\mu }_x(f)`$ the length of its stalk at a point $`xV(𝒥(f))`$.
###### Lemma 2
$$\underset{x^2}{}\stackrel{~}{\mu }_x(f)=d^2d_t,$$
where $`d_t`$ is the degree of the map $`f`$.
Proof. See \[Fu\], 4.4.
Recall that for any singular point $`x`$ of $`V(F)`$ we have the conductor invariant $`\delta _x`$ defined as the length of the quotient module $`\overline{𝒪}_{C,x}/𝒪_{C,x}`$, where $`\overline{𝒪}_{C,x}`$ is the normalization of the local ring $`𝒪_{C,x}`$. Let $`r_x`$ denote the number of local branches of $`C`$ at $`x`$. We have the following
###### Lemma 3
Let $`\stackrel{~}{\mu }_x=\stackrel{~}{\mu }_x(f)`$, where $`f`$ is the map defined by the polar linear system $`𝒫_F`$. For any $`xC`$,
$$\stackrel{~}{\mu }_x2\delta _xr_x+1.$$
$`(3.1)`$
Proof. Without loss of generality we may assume that $`x=(1,0,0)`$. Let $`\stackrel{~}{P}(X,Y)`$ denote the dehomogenization of a homogeneous polynomial $`P`$ with respect to the variable $`x_0`$. Applying the Euler formula $`dF=x_0F_0+x_1F_1+x_2F_2`$, we obtain that
$$𝒥_x=(\stackrel{~}{F},\frac{\stackrel{~}{F}}{X},\frac{\stackrel{~}{F}}{Y})_x.$$
By Jung-Milnor’s formula (see \[Mi\], Theorem 10.5), the length $`\mu _x`$ of the module $`(k[X,Y]/(\frac{\stackrel{~}{F}}{X},\frac{\stackrel{~}{F}}{Y}))_x`$ is equal to $`2\delta _xr_x+1`$. It remains to observe that $`\stackrel{~}{\mu }_x\mu _x`$.
The next lemma is a well-known formula for the arithmetic genus of a plane curve.
###### Lemma 4
$$p_a(C)=(d1)(d2)/2=\underset{i=1}{\overset{h}{}}g_i+\underset{x}{}\delta _xh+1,$$
$`(3.2)`$
where $`h`$ is the number of irreducible components $`C_i`$ of $`C`$ and $`g_i`$ is the genus of the normalization of $`C_i`$.
The next formula is an easy consequence of the incidence relation count for pairs of lines, but just for fun we give a high-brow proof of this:
###### Corollary 2
Let $`\{L_1,\mathrm{},L_s\}`$ be a set of lines in $`^2`$. Let $`a_i`$ denote the number of points which belong to $`i2`$ distinct lines. Then
$$s(s1)=\underset{i=2}{\overset{s}{}}a_ii(i1).$$
$`(3.3)`$
Proof. We apply the previous formula to the curve $`L=L_1+\mathrm{}+L_s`$. Each singular point of $`L`$ lies on the intersection of $`i2`$ lines. It is isomorphic locally to the singular point of the affine curve given by an equation $`_{j=1}^i(\alpha _jX+\beta _jY)=0`$. It is easy to compute $`\delta _x`$. It is equal to $`i(i1)/2`$. Since $`r_x=i`$, by Lemma 4, we have
$$(s1)(s2)/2=\underset{i=2}{\overset{s}{}}a_ii(i1)/2s+1.$$
This is equivalent to the claimed formula.
###### Theorem 4
Let $`F`$ be a homaloidal polynomial in three variables without multiple factors. Then, after a linear change of variables, it coincides with one from Examples 1, 2.1, 4. In other words, $`C=V(F)`$ is one of the following curves:
(i) a nonsingular conic;
(ii) the union of three nonconcurrent lines;
(iii) the union of a conic and its tangent.
Proof. Since $`𝒫_F`$ is homaloidal, we can apply Lemma 2 to obtain
$$d^22d=\underset{xC}{}\stackrel{~}{\mu }_x.$$
$`(3.4)`$
By Lemma 3,
$$d^22d\underset{xC}{}(2\delta _xr_x+1).$$
By Lemma 4,
$$d^23d=2\underset{i=1}{\overset{h}{}}g_i+2\underset{xC}{}\delta _x2h.$$
$`(3.5)`$
Let $`C_1,\mathrm{},C_h`$ be irreducible components of $`C`$ and $`d_i=\mathrm{deg}C_i`$. Using (3.4) and (3.5), we obtain
$$\underset{i=1}{\overset{h}{}}(2d_i)=d+2h2\underset{i=1}{\overset{h}{}}g_i+\underset{xC}{}(r_x1)0.$$
$`(3.6)`$
The rest of the proof consists of analyzing this inequality. First observe that each point of intersection of two irreducible components gives a positive contribution to the sum $`_{i=1}^k(r_i1)`$. This immediately implies that $`d_i=1`$ for some $`i`$ unless $`C`$ is an irreducible conic. In the latter case it is obviously nonsingular (otherwise the polar linear system is a pencil). This is case (i) of the theorem. So we may assume that $`C_1,\mathrm{},C_s`$ are lines. It follows from (3.6) that
$$0\underset{i=s+1}{\overset{h}{}}(2d_i)2\underset{i=1}{\overset{h}{}}g_i+\underset{xC}{}(r_x1)s.$$
$`(3.7)`$
If $`s=1`$, then each point of intersection of $`C_1`$ with other component of $`C`$ contributes at least 1 to the sum $`_{i=1}^k(r_i1)`$. This shows that $`C=C_1+C_2`$, where $`L`$ intersects $`C_2`$ at one point and $`d_2=2`$. This is case (iii) of the theorem.
Assume that $`s2`$. Let $`x_1,\mathrm{},x_N`$ be the intersection points of the lines $`C_1,\mathrm{},C_s.`$ Let $`a_j`$ be the number of points among them which belong to $`j2`$ lines. Then $`_{j=2}^sa_j=N`$, and
$$\underset{xC}{}(r_x1)s\underset{i=1}{\overset{N}{}}(r_i1)s\underset{j=2}{\overset{s}{}}ja_jNs=\underset{j=2}{\overset{s}{}}(j1)a_js.$$
$`(3.8)`$
By (3.3),
$$s=\underset{j=2}{\overset{s}{}}\frac{j}{s1}a_j(j1).$$
Assume not all lines pass through one point, i.e. $`a_s=0`$. Then $`js1`$ for all $`j`$ with $`a_j0`$. In this case
$$s\underset{j=2}{\overset{s}{}}a_j(j1)$$
$`(3.9)`$
and the equality holds if and only if $`a_j=0`$ for all $`js1`$. If $`p_i`$ is a point lying on $`s1`$ lines, then the remaining line must intersect other lines at points different from $`p_i`$. This gives that $`a_20`$. So, if the equality holds, we have $`s=3`$ and $`a_2=N=3`$. If $`hs`$, then $`C_h`$ is of degree $`>1`$. Its points of intersection with three lines give positive contribution to the sum $`_{xx_1,\mathrm{},x_N}(r_x1)s`$. Thus (3.8) is a strict inequality contradicting (3.7). So $`C`$ is the union of three nonconcurrent lines which is case (ii) of the theorem.
It remains to consider the case when all lines pass through one point. In this case (3.7) implies that $`s<h`$. Then $`C_h`$ is of degree $`>1`$. Assume $`p_1C_h`$. Then $`r_1s+1`$ and
$$\underset{xC}{}(r_x1)s=(r_11s)+\underset{xp_1}{}(r_x1)0.$$
$`(3.9)`$
It follows from (3.7) that $`C_h`$ is a nonsingular conic. Since $`s2`$, one of the lines is not tangent to $`C_h`$ at $`p_1`$ and hence intersects $`C_h`$ at some point $`xp_1`$. Thus (3.9) is a strict inequality. This contradicts (3.7). If $`p_1C_h`$, then $`C_h`$ intersects each line so that we have $`_{xp_1}(r_x1)s`$ and
$$\underset{xC}{}(r_x1)s=(r_11s)+\underset{xp_1}{}(r_x1)s1>0.$$
Again a contradiction.
Let us notice the following combinatorial fact which follows from the proof of the above theorem in the case when $`C`$ is the union of lines:
###### Corollary 3
Let $`C`$ consist of $`s`$ lines $`l_1,\mathrm{},l_s`$. For each line $`l_i`$ let $`k_i`$ be the number of singular points of $`C`$ on $`l_i`$, and let $`t`$ be the total number of singular points. Assume that $`t>1`$. Then
$$\underset{i=1}{\overset{s}{}}(k_i1)t$$
and the equality takes place if and only if $`t=3,s=3`$.
Proof. Let $`d`$ be the degree of the map given by the polar linear system of the polynomial defining $`C`$. We resolve the indeterminacy points by blowing up the singular points of $`C`$. Let $`E_p`$ be the exceptional curve blow-up from the point $`p`$, $`h`$ be the class of a general line and $`m_p`$ be the multiplicity of a singular point $`p`$. Then
$$d=((s1)h\underset{p\text{Sing}(C)}{}(m_p1)E_p)^2=(s1)^2\underset{p\text{Sing}(C)}{}(m_p1)^2.$$
Let $`a_i=\mathrm{\#}\{p:m_p=i\}`$. Applying equality (3.3), we can rewrite it as follows:
$$d=s(s1)(s1)\underset{i=2}{\overset{s}{}}a_i(i1)i+\underset{i=2}{\overset{s}{}}a_i(i1)=$$
$$(s1)+\underset{i=2}{\overset{s}{}}a_i(i1)=s+1+\underset{i=2}{\overset{s}{}}ia_i\underset{i=2}{\overset{s}{}}a_i.$$
Now the standard incidence relation argument gives us
$$\underset{i=2}{\overset{s}{}}ia_i=\underset{p\text{Sing}(C)}{}m_p=\underset{i=1}{\overset{s}{}}k_i.$$
This allows us to rewrite the expression for $`d`$ in the form
$$d=1+\underset{i=1}{\overset{s}{}}(k_i1)t.$$
Now $`d1`$ unless all lines pass though one point and, by Theorem 4, $`d=1`$ if and only if $`s=3,t=3`$.
Remark As was explained to me by Hal Schenck, in the case of a real arrangement of lines the previous Corollary follows easily from the Euler formula applied to the cellular subdivision of $`^2`$ defined by the arrangement. One interprets the left-hand side as the number $`f_1`$ of edges, the right-hand side as the number $`f_0`$ of vertices and uses that $`f_0s`$ and $`f_2f_0+1`$ if the arrangement is not a pencil (see \[Gr\], pp.10 and 12).
Unfortunately, the argument used in the proof of Theorem 4 does not apply to non-reduced polynomials. However, the following conjecture seems to be reasonable:
###### Conjecture
Let $`F=A_1^{m_1}\mathrm{}A_s^{m_s}`$ be the factorization of $`F`$ into prime factors. Let $`G=A_1\mathrm{}A_s`$. Then the polar linear system $`𝒫_F`$ is homaloidal if and only if $`𝒫_G`$ is homaloidal.
## 4. Arrangements of hyperplanes in $`^3`$
Here we shall consider the special case when $`F=_{i=1}^nL_i`$ is the product of linear polynomials in four variables without multiple factors. Its set of zeroes is an arrangement of hyperplanes in $`^3`$.
Let $`𝒜=\{H_1,\mathrm{},H_N\}`$ be the set of planes $`\{L_i=0\}`$, $``$ be the set of lines which are contained in more than one plane $`H_i`$, and $`𝒫`$ be the set of points which are contained in more than two planes $`H_i`$. For any $`l`$, set
$$k_l=\mathrm{\#}\{i:lH_i\},a_l=\mathrm{\#}\{p𝒫:pl\}.$$
For any $`p𝒫`$ set
$$k_p=\mathrm{\#}\{i:pH_i\}.$$
We define $`d_𝒜`$ to be the degree of the polar linear system defined by $`F`$.
###### Lemma 5
$$d_𝒜=(N1)^3\underset{p𝒫}{}(k_p1)+\underset{l}{}(k_l1)(a_l1).$$
Proof. We can resolve the points of indeterminacy of $`𝒫_F`$ by first blowing up each points $`p𝒫`$ followed by blowing up the proper transforms of each line $`l`$. Let
$$D=\underset{p𝒫}{}(k_p1)E_p+\underset{l}{}(k_l1)E_l.$$
Here the notations are self-explanatary. We have (see \[Fu\])
$$d_𝒜=((N1)HD)^3,$$
where $`H`$ is the pre-image of a general plane in the blow-up. Using the standard formulae for the blow-up a smooth subvariety, we have
$$E_l^3=c_1(N_{\overline{l}})=[(4H2\underset{l,pl}{}E_p)\overline{l}2]=2a_l2.$$
Here $`\overline{l}`$ denotes the proper transform of the line $`l`$ under the blowing up the points from $`𝒫`$, and $`N_{\overline{l}}`$ is the normal bundle of $`\overline{l}`$. Next, we have
$$E_l^2E_p=1,E_p^3=1.$$
Collecting this together we get
$$D^3=\underset{l}{}(k_l1)^3(2a_l2)+\underset{p𝒫}{}(k_p1)^33\underset{l,pl}{}(k_l1)^2(k_p1),$$
$$HD^2=\underset{l}{}(k_l1)^2E_lH=\underset{l}{}(k_l1)^2,$$
$$H^2D=0.$$
This gives
$$d_𝒜=(N1)^33(N1)\underset{l}{}(k_l1)^2\underset{l}{}(k_l1)^3(2a_l2)$$
$$\underset{p𝒫}{}(k_p1)^3+3\underset{l,pl}{}(k_l1)^2(k_p1).$$
Observe now that
$$\underset{pl}{}(k_p1)=\underset{pl}{}k_pa_l=(a_lk_l+Nk_l)a_l=(a_l1)k_l+Na_l.$$
This allows us to rewrite the expression for $`d`$ as follows:
$$d_𝒜=(N1)^33(N1)\underset{l}{}(k_l1)^2\underset{l}{}(k_l1)^3(2a_l2)$$
$$\underset{p𝒫}{}(k_p1)^3+3\underset{l}{}(k_l1)^3(a_l1)+3(N1)\underset{l}{}(k_l1)^2=$$
$$(N1)^3\underset{p𝒫}{}(k_p1)+\underset{l}{}(k_l1)(a_l1).$$
This proves the lemma.
###### Lemma 6
Let
$$t_s=\mathrm{\#}\{p𝒫:k_p=s\},t_q(1)=\mathrm{\#}\{l:k_l=q\},$$
$$t_{sq}=\underset{l:k_l=q}{}\mathrm{\#}\{pl:k_p=s\}.$$
Then
$$\left(\genfrac{}{}{0pt}{}{N}{3}\right)=\underset{s}{}\left(\genfrac{}{}{0pt}{}{s}{3}\right)t_s\underset{s,q}{}\left(\genfrac{}{}{0pt}{}{q}{3}\right)(t_{sq}t_q(1)).$$
Proof. This is a three-dimensional analog of Corollary 2 to Lemma 4. It easily follows from the incidence relation count for triples of distinct planes and points and lines.
###### Corollary 4
$$d_𝒜=N1\underset{p𝒫}{}(k_p1)+\underset{l}{}(a_l1)(k_l1).$$
Proof. Combine the previous two lemmas.
###### Lemma 7
Let $`𝒜`$ be an arrangement of $`N`$ hyperplanes in $`^3`$ defined by a polynomial $`F`$. The following properties are equivalent:
(i) all planes pass through a point;
(ii) the partials of $`F`$ are linearly dependent
(ii) $`d_𝒜=0`$.
Proof. Obvious.
###### Lemma 8
Let $`𝒜`$ be an arrangement of $`N`$ planes. Let $`𝒜^{}`$ be a new arrangement obtained by adding one more plane to $`𝒜`$. Assume $`d_𝒜0`$. Then
$$d_𝒜^{}>d_𝒜.$$
Proof. Let
$$𝒫^{}=\{p𝒫:pH\},^{}=\{l:lH\},$$
$$^{\prime \prime }=\{l:pl\text{for any }p𝒫^{}\},$$
$$𝒩=\{lH(H_1\mathrm{}H_N)\}.$$
Note that each line $`l𝒩`$ is a double line and each line $`l^{\prime \prime }`$ contains one new singular point $`Hl`$ of multiplicity $`k_l+1`$. Applying the previous corollary, we obtain
$$d_𝒜^{}=N\underset{p𝒫𝒫^{}}{}(k_p1)\underset{p𝒫^{}}{}k_p\underset{l^{\prime \prime }}{}k_l+\underset{l^{}}{}k_l(a_l1)+$$
$$\underset{l^{}}{}(k_l1)a_l+\underset{l𝒩}{}(a_l^{}1),$$
where $`a_l^{}`$ denotes the number $`a_l`$ defined for the extended arrangement. Applying the corollary again, we get
$$d_𝒜^{}d_𝒜=1+(\underset{l(^{}^{\prime \prime }}{}(k_l1)\mathrm{\#}𝒫^{})+(\underset{l𝒩}{}(a_l^{}1)\mathrm{\#}^{\prime \prime })+\underset{l^{}}{}(a_l1).$$
$`(4.1)`$
For each $`p𝒫^{}`$ there exists a line $`l(^{}^{\prime \prime })`$ passing through $`p`$. Since $`k_l>1`$ for each line we see that $`_{l(^{}^{\prime \prime }}(k_l1)\mathrm{\#}𝒫^{}0`$. Now consider the arrangement of lines in the plane $`H`$ formed by the lines $`l𝒩`$. Its multiple points are the points of intersection of $`H`$ with lines in $`^{\prime \prime }`$. Applying Corollary 3 to Theorem 4, we see that $`_{l𝒩}(a_l^{}1)\mathrm{\#}^{\prime \prime }0`$ unless there is only one line in $`^{\prime \prime }`$ when this difference is equal to $`1`$. But in this case $`H`$ must contain at least one line from $``$ and hence there is an additional term $`_l^{}(a_l1)`$. If it is zero, then each line $`l^{}`$ contains only one singular point of the arrangement. This implies that all planes except maybe one contain $`l`$. In this case all planes pass through a point and $`d_𝒜=0`$. So the term is positive and we have proved the inequality $`d_𝒜^{}>d_𝒜`$.
###### Theorem 5
Let $`𝒜`$ be an arrangement of $`N`$ planes in $`^3`$ with $`d_𝒜=1`$. Then $`𝒜`$ is the union of four planes in general linear position.
Proof. By the previous lemma deleting any plane $`H`$ from the arrangement $`𝒜`$ defines an arrangement $`𝒜^{}`$ with $`d_𝒜=0`$. We may assume that $`H`$ does not pass through the common point of the planes from $`𝒜^{}`$. In the notation of the proof of the previous lemma, where the new arrangement is our $`𝒜`$ and the old one is $`𝒜\{H\}`$, we have $`\mathrm{\#}^{\prime \prime }=N1`$. Now the term $`(_{l𝒩}(a_l^{}1)\mathrm{\#}^{\prime \prime })`$ in (4.1) must be equal to zero since otherwise $`d_𝒜>1`$. By Lemma 6, $`N1=3`$. Thus $`N=4`$. Since $`d_𝒜0`$, the planes do not have a common point and hence the arrangement is as in the assertion of the theorem.
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# Effects of R-parity violation on CP asymmetries in Λ_𝑏→𝑝𝜋 decay
## Abstract
We have studied new CP violating effects in $`\mathrm{\Lambda }_bp\pi `$ decay mode, that can arise in Minimal Supersymmetric Standard Model with R-parity violation. We have estimated how much R-parity violation modifies the Standard Model predictions for CP asymmetries within the present bounds. We found that in the R-parity violating model, the rate asymmetry ($`a_{cp}`$) is suppressed (about 10 times) and the asymmetry parameter $`A(\alpha )`$ is enhanced (approximately $`𝒪(10^2)`$ times) with respect to the SM predictions.
One of the most important objects of the upcoming experiments at $`B`$ factories is to search for CP violation in as many $`B`$ decay modes as possible so as to establish the pattern of CP violation among various $`B`$ decays . This then may allow for an experimental test not only of the Standard Model (SM) Cabibbo-Kobayashi-Maskawa (CKM) paradigm of CP violation, but also many extensions of the SM that often contain new sources of CP violation. It is well known that CP violating $`B`$ decays might constitute an important hunting ground for new physics. This is particularly so since many CP violating asymmetries related to $`B`$ decays are predicted to be small in SM, are likely to be measured with high precision in the upcoming $`B`$ factories. Measurements larger than the SM predictions would definitely signal the presence of new physics. It is also interesting to study CP violation in bottom baryon system in order to find the physical channels which may have large CP asymmetries, even though the branching ratios for such processes are usually smaller than those for the corresponding bottom mesons. Recently some data on the bottom baryon $`\mathrm{\Lambda }_b`$ have appeared. For instance, OPAL has measured its lifetime and the production branching ratio for the inclusivs semileptonic decay $`\mathrm{\Lambda }_b\mathrm{\Lambda }l^{}\overline{\nu }X`$ . Furthermore, mesurements of the nonleptonic decay $`\mathrm{\Lambda }_b\mathrm{\Lambda }J/\psi `$ has also been reported . Certainly we expect more data in the future in the bottom baryon sector. In this paper we intend to study CP violation in the nonleptonic $`\mathrm{\Lambda }_bp\pi `$ decay in the Minimal Supersymmetric Standard Model (MSSM) with $`R`$-parity violation . The MSSM has been widely considered as a leading candidate for new physics beyond SM. In supersymmetric theories ‘$`R`$-parity’ is a discrete symmetry under which all standard model particles are even while their superpartner are odd. It is defined as $`R=(1)^{(3B+L+2S)}`$, where $`S`$ is the spin, $`B`$ is the baryon number and $`L`$ is the lepton number of the particle. An exact $`R`$-parity implies that superparticles could be produced in pairs and the lightest supersymmetric particle (LSP) is stable. However, $`B`$ and $`L`$ conservations are not ensured by gauge invariance and therefore it is worhwhile to investigate what happens to the CP asymmetries when $`R`$-parity is violated.
The most general Lorentz-invariant amplitude for the decay $`\mathrm{\Lambda }_bp\pi ^{}`$ can be written as
$$i\overline{u}_p(p_f)(a+b\gamma _5)u_{\mathrm{\Lambda }_b}(p_i)$$
(1)
The corresponding matrix element for $`\overline{\mathrm{\Lambda }}_b\overline{p}\pi ^+`$ is then
$$i\overline{v}_{\overline{p}}(p_f)(a^{}+b^{}\gamma _5)v_{\overline{\mathrm{\Lambda }}_b}(p_i)$$
(2)
It is convenient to express the transition amplitude in terms of S-wave (parity violating) and P-wave (parity conserving) amplitudes S and P as
$$S+P\sigma \widehat{𝐪}$$
(3)
where $`𝐪`$ is the proton momentum in the rest frame of $`\mathrm{\Lambda }_b`$ baryon and the amplitudes S and P are :
$$S=a\sqrt{\frac{\{(m_{\mathrm{\Lambda }_b}+m_p)^2m_\pi ^2\}}{16\pi m_{\mathrm{\Lambda }_b}^2}}P=b\sqrt{\frac{\{(m_{\mathrm{\Lambda }_b}m_p)^2m_\pi ^2\}}{16\pi m_{\mathrm{\Lambda }_b}^2}}$$
(4)
The experimental observables are the total decay rate $`\mathrm{\Gamma }`$ and the decay parameters $`\alpha `$, $`\beta `$ and $`\gamma `$ which govern the decay-angular distribution and the polarization of the final baryon. The decay rate is given as
$$\mathrm{\Gamma }=2|𝐪|\{|S|^2+|P|^2\}$$
(5)
and the dominant asymmetry parameter ($`\alpha `$) is given as
$$\alpha =\frac{2Re(S^{}P)}{\{|S|^2+|P|^2\}}$$
(6)
Similar observables for the antihyperon decays are $`\overline{\mathrm{\Gamma }}`$ and $`\overline{\alpha }`$ are given as
$$\overline{\mathrm{\Gamma }}=2|𝐪|\{|\overline{S}|^2+|\overline{P}|^2\},\overline{\alpha }=\frac{2Re(\overline{S}^{}\overline{P})}{\{|\overline{S}|^2+|\overline{P}|^2\}}$$
(7)
For $`\mathrm{\Lambda }_bp\pi ^{}`$ decay the CP violating rate asymmetry in partial decay rate ($`a_{cp}`$) and aymmetry parameter ($`A(\alpha )`$) are defined as follows,
$$a_{cp}=\frac{\mathrm{\Gamma }(\mathrm{\Lambda }_bp\pi ^{})\mathrm{\Gamma }(\overline{\mathrm{\Lambda }}_b\overline{p}\pi ^+)}{\mathrm{\Gamma }(\mathrm{\Lambda }_bp\pi ^{})+\mathrm{\Gamma }(\overline{\mathrm{\Lambda }}_b\overline{p}\pi ^+)},$$
(8)
$$A(\alpha )=\frac{\alpha +\overline{\alpha }}{\alpha \overline{\alpha }},$$
(9)
A nonzero value for $`a_{cp}`$ and $`A(\alpha )`$ will signal CP violation. The existence of such CP asymmetries require the interference of two decay amplitudes with different weak and strong phase differences. The weak phase difference arises from the superposition of various penguin contributions and the usual tree diagrams while the strong phases are induced by final state interactions (FSI). At the quark level, the strong phase diffences arise through the absorptive parts of penguin diagrams (hard final state interactions) and nonperturbatively (soft final state interactions) . In the absence of an argument that the parton-hadron duality should hold in exclusive processes, one can not exclude that the weak transition matrix elements receive phases originating from soft FSI. However the effects of soft FSI are extremely difficult to quantify. In the absence of a reliable theoretical calculation for soft FSI, we make the usual approximation of retaining the absorptive part from quark level calculation (hard FSI) for strong phase differences in our analysis.
We shall first consider the SM contributions to the transition amplitude. The effective Hamiltonian $`_{eff}`$ for the decay process $`\mathrm{\Lambda }_bp\pi ^{}`$ is given as
$$_{eff}=\frac{G_F}{\sqrt{2}}\{V_{ub}V_{ud}^{}[c_1(\mu )O_1^u(\mu )+c_2(\mu )O_2^u(\mu )]V_{tb}V_{td}^{}\underset{i=3}{\overset{10}{}}c_i(\mu )O_i(\mu )\}+\mathrm{h}.\mathrm{c}.,$$
(10)
where $`O_{1,2}`$ are the tree level current-current operators, $`O_{36}`$ are the QCD and $`O_{710}`$ are the electroweak penguin operators which are explicitly given in Ref. , $`c_i`$’s are the Wilson coefficients, which take care of the short-distance QCD corrections, are scheme and scale dependent. These unphysical dependences are cancelled by the corresponding scheme and scale dependences of the matrix elements of the operators. However, in the factorization approximation, the hadronic matrix elements are written in terms of form factors and decay constants, which are scheme and scale independent. So to achieve the cancellation, the various one loop corrections are absorbed into the effective Wilson coefficients $`c_i^{eff}`$, which are scheme and scale independent. The values of the effective Wilson coefficients for $`bd`$ transitions are explicitly evaluated in Ref. as :
$`c_1^{eff}=1.168c_2^{eff}=0.365c_3^{eff}=0.0224+i0.0038c_4^{eff}=(0.0454+i0.0115)`$ (11)
$`c_5^{eff}=0.0131+i0.0038c_6^{eff}=(0.0475+i0.0115)c_7^{eff}/\alpha =(0.0294+i0.0329)`$ (12)
$`c_8^{eff}/\alpha =0.055c_9^{eff}/\alpha =(1.426+i0.0329)c_{10}^{eff}/\alpha =0.48`$ (13)
These one loop corrections (to get $`c_i^{eff}`$’s ) result in imaginary parts for ($`c_i^{eff}`$’s) due to virtual quarks going on their mass shell.
The matrix elements of the operators can be calculated using the factorization approximation. In this approximation the hadronic matrix elements of the four quark operators $`(\overline{d}b)_{VA}(\overline{u}d)_{VA}`$ split into products of matrix elements one involving pion decay constant and the other dealt the baryonic form factors. The matrix elements of the $`(VA)(V+A)`$ i.e., ($`O_6`$ and $`O_8`$) operators can be obtained by Fierz reordering and using the Dirac equation as,
$$p\pi |O_6|\mathrm{\Lambda }_b=2\underset{q}{}\pi |\overline{d}(1+\gamma _5)q|0p|\overline{q}(1\gamma _5)b|\mathrm{\Lambda }_b$$
(14)
Using the Dirac equation the matrix element can be rewritten as
$$p\pi |O_6|\mathrm{\Lambda }_b=[R_1p|V_\mu |\mathrm{\Lambda }_bR_2p|A_\mu |\mathrm{\Lambda }_b]\pi |A_\mu |0,$$
(15)
with
$$R_1=\frac{2m_\pi ^2}{(m_bm_u)(m_d+m_u)},R_2=\frac{2m_\pi ^2}{(m_b+m_u)(m_d+m_u)},$$
(16)
where the quark masses are the current quark masses. The same relation works for $`O_8`$.
Thus under the factorization approximation the baryon decay amplitude is governed by a decay constant and baryonic transition form factors. The general expression for the baryon transition is given as
$`p(p_f)|V_\mu A_\mu |\mathrm{\Lambda }_b(p_i)`$ $`=`$ $`\overline{u}_p(p_f)\{f_1(q^2)\gamma _\mu +if_2(q^2)\sigma _{\mu \nu }q^\nu +f_3(q^2)q_\mu `$ (17)
$``$ $`[g_1(q^2)\gamma _\mu +ig_2(q^2)\sigma _{\mu \nu }q^\nu +g_3(q^2)q_\mu ]\gamma _5\}u_{\mathrm{\Lambda }_b}(p_i),`$ (18)
where $`q=p_ip_f`$. The values of the form factors at maximum momentum transfer are evaluated in nonrelativistic quark model and their $`q^2`$ dependence are determined using the pole dominance model with values as,
$$f_1(m_\pi ^2)=0.043m_if_3(m_\pi ^2)=0.009g_1(m_\pi ^2)=0.092m_ig_3(m_\pi ^2)=0.047,$$
(19)
where the particle masses are taken from .
Hence one obtains the amplitude for the decay mode $`\mathrm{\Lambda }_bp\pi ^{}`$ as (where the factor $`G_F/\sqrt{2}`$ is suppressed)
$`A(\mathrm{\Lambda }_bp\pi ^{})`$ $`=`$ $`if_\pi \overline{u}_p(p_f)[\{\lambda _u(a_1+a_4+a_{10}+(a_6+a_8)R_1)+\lambda _c(a_4+a_{10}+(a_6+a_8)R_1)\}`$ (20)
$`\times `$ $`\left(f_1(m_\pi ^2)(m_im_f)+f_3(m_\pi ^2)m_\pi ^2\right)`$ (21)
$`+`$ $`\{\lambda _u(a_1+a_4+a_{10}+(a_6+a_8)R_2)+\lambda _c(a_4+a_{10}+(a_6+a_8)R_2)\}`$ (22)
$`\times `$ $`(g_1(m_\pi ^2)(m_i+m_f)g_3(m_\pi ^2)m_\pi ^2)\gamma _5]u_{\mathrm{\Lambda }_b}(p_i),`$ (23)
where $`m_i`$ and $`m_f`$ are the masses of the initial and final baryons respectively. The coefficients $`a_1,a_2\mathrm{}a_{10}`$ are combinations of the effective Wilson coefficients given as
$$a_{2i1}=c_{2i1}^{eff}+\frac{1}{(N_c)}c_{2i}^{eff}a_{2i}=c_{2i}^{eff}+\frac{1}{(N_c)}c_{2i1}^{eff}i=1,2\mathrm{}5,$$
(24)
where $`N_c`$ is the number of colors, taken to be $`N_c=3`$ for naive factorization. Thus one obtains the S and P-wave amplitudes using eqns. (1), (4) and (17), in units of $`(10^9)`$ as
$`S=\lambda _u(34.6030.7115i)\lambda _c(2.782+0.7115i)`$ (25)
$`P=\lambda _u(74.0561.521i)\lambda _c(5.95+1.521i),`$ (26)
with $`\lambda _i=V_{ib}V_{id}`$. Now we shall proceed to evaluate the R-parity violating ($`\overline{)}R_p`$) amplitude. In the MSSM the most general R-parity violating superpotential is given as
$$W_{\overline{)}R_p}=\lambda _{ijk}L_iL_jE_k^c+\lambda _{ijk}^{}L_iQ_jD_k^c+\lambda _{ijk}^{\prime \prime }U_i^cD_j^cD_k^c,$$
(27)
where $`i,j,k`$ are the generation indices and we assume that possible bilinear terms $`\mu _iL_iH_2`$ can be rotated away. $`L_i`$ and $`Q_i`$ are the $`SU(2)`$-doublets for lepton and quark superfields and $`E_i^c`$, $`U_i^c`$ and $`D_i^c`$ are the singlet superfields. $`\lambda _{ijk}`$ and $`\lambda _{ijk}^{\prime \prime }`$ are antisymmetric under the interchange of the first two and last two indices. The first two terms violate lepton numbers where as the last term violates baryon number. For our purpose we will consider only the lepton number violation contributions. As the $`\lambda `$ type couplings do not contribute to the nonleptonic decays we obtain from eqn. (27) the following effective Hamiltonian due to the exchange of sleptons as
$$_{\overline{)}R_p}^{eff}=\underset{n,p,q=1}{\overset{3}{}}\frac{\lambda _{npi}^{}\lambda _{nql}^{}}{M_{\stackrel{~}{l}_n}^2}V_{kq}V_{jp}^{}(\overline{d}_iP_Lu_j)(\overline{u}_kP_Rd_l)$$
(28)
with $`P_{L,R}=(1\gamma _5)/2`$. From the above effective Hamiltonian we calculate the amplitude $`𝒜_{\overline{)}R_p}(\mathrm{\Lambda }_bp\pi )`$ using the factorization approximation. The matrix elements of the $`(SP)(S+P)`$ operators are obtained using the Dirac equation of motion. Assuming $`V_{CKM}`$ is given by only down-type quark sector we obtain the dominant transition amplitude to be
$`𝒜_{\overline{)}R_p}^{eff}`$ $`=`$ $`{\displaystyle \underset{n,=2,3}{}}{\displaystyle \frac{\lambda _{npi}^{}\lambda _{nql}^{}}{M_{\stackrel{~}{l}_n}^2}}V_{11}V_{11}^{}\times if_\pi \overline{u}(p_f)\left[R_1\right(f_1(m_\pi ^2)(m_im_f)+f_3(m_\pi ^2)m_\pi ^2)`$ (29)
$`+`$ $`R_2(g_1(m_\pi ^2)(m_i+m_f)g_3(m_\pi ^2)m_\pi ^2)\gamma _5]u_{\mathrm{\Lambda }_b}(p_i).`$ (30)
Now considering the slepton mass to be 100 GeV, the present bounds on $`\lambda _{ijk}^{}`$ are
$$\lambda _{211}^{}<0.09\lambda _{213}^{}<0.09\lambda _{311}^{}<0.16\lambda _{313}^{}<0.16$$
(31)
we obtain the S and P-wave $`R_{\overline{)}R_p}`$ amplitudes to be
$$S_{\overline{)}R_p}<1.626\times 10^9P_{\overline{)}R_p}<3.474\times 10^9$$
(32)
After obtaining the transition amplitude in SM and $`R_{\overline{)}R_p}`$ model we now proceed to estimate the CP asymmetries. The parity violating (S wave) and parity conserving (P wave) amplitudes can be explicitly written as
$`S=\lambda _uS_u+\lambda _cS_c+S_{\overline{)}R_p}`$ (33)
$`P=\lambda _uP_u+\lambda _cP_c+P_{\overline{)}R_p}`$ (34)
where $`\lambda _i=V_{ib}V_{id}`$, are the product of CKM matrix elements which contain the weak phases. The strong phases which arise from the perturbative penguin diagrams at one loop level, are contained in $`S_{u/c}`$ and $`P_{u/c}`$ i.e., $`S_u=|S_u|e^{i\delta _u}`$ etc. The corresponding quantities for the antihyperon decays are given as
$`\overline{S}=\left(\lambda _u^{}S_u+\lambda _c^{}S_c+S_{\overline{)}R_p}\right)`$ (35)
$`\overline{P}=\lambda _u^{}P_u+\lambda _c^{}P_c+P_{\overline{)}R_p}`$ (36)
Thus the CP violating rate asymmetry is given as,
$$a_{cp}=\frac{2[Im(\lambda _u\lambda _c^{})(Im(S_uS_c^{})+Im(P_uP_c^{}))+Im(\lambda _u)[S_{\overline{)}R_p}Im(S_u)+P_{\overline{)}R_p}Im(P_u)]]}{A}$$
(37)
where
$`A`$ $`=`$ $`[|\lambda _u|^2(|S_u|^2+|P_u|^2)+|\lambda _c|^2(|S_c|^2+|P_c|^2)+(|S_{\overline{)}R_p}|^2+|P_{\overline{)}R_p}|^2)`$ (38)
$`+`$ $`2Re(\lambda _u\lambda _c^{})(Re(S_uS_c^{})+Re(P_uP_c^{}))+2{\displaystyle \underset{i=u,c}{}}Re(\lambda _i)(S_{\overline{)}R_p}Re(S_i)+P_{\overline{)}R_p}Re(P_i))]`$ (39)
Using the Wolfenstein parametrization for CKM matrix elements with values $`A=0.815`$, $`\lambda =0.2205`$, $`\rho =0.175`$ and $`\eta =0.37`$, we obtain the branching ratio and CP violating observables in RPV model using eqns. (5), (9) and (27) as
$`Br(\mathrm{\Lambda }_bp\pi )<1.6\times 10^4`$ (40)
$`a_{cp}0.3\%`$ (41)
$`A(\alpha )8.8\times 10^3`$ (42)
The corresponding quantities in the SM ($`S_{\overline{)}R_p}=0`$ and $`P_{\overline{)}R_p}=0`$) are given as
$`Br(\mathrm{\Lambda }_bp\pi )=0.9\times 10^6`$ (43)
$`a_{cp}=8.3\%`$ (44)
$`A(\alpha )=2.3\times 10^5`$ (45)
It can be seen from eqns. (42) and (45) that the effects of R-parity and lepton number violating couplings significantly modify the SM results of the branching ratio and CP asymmetry parameters for the decay mode $`\mathrm{\Lambda }_bp\pi `$. The branching ratio and the asymmetry parameter ($`A(\alpha )`$) in RPV model are approximately $`𝒪(10^2)`$ times larger than the SM contributions whereas the rate asymmetry $`a_{cp}`$ is nearly 10 times smaller than the SM result.
To summarize, in this work we have studied the effects of R-parity violating couplings on the direct CP asymmetry parameters in $`\mathrm{\Lambda }_bp\pi `$ decay mode. Assuming factorization, we have used the nonrelativistic quark model to evaluate the form factors at maximum momentum transfer ($`q_m^2`$) and the extrapolation of the form factors from $`q_m^2`$ to the required $`q^2`$ value is done by using the pole dominance. Although there are significant uncertainties in our estimates as we have used the factorization approximation to evaluate the matrix elements of the four-quark current operators and taken all the R-parity violating couplings to be real, it is probably safe to say that the asymmetry parameter $`A(\alpha )`$ in $`\mathrm{\Lambda }_bp\pi `$ decay is significantly larger than the corresponding asymmetry in the Standard Model.
The author would like to thank Professor M. P. Khanna and Dr. A. K. Giri for many useful discussions and also to CSIR, Govt. of India, for financial support.
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# Identifying Old Tidal Dwarf Irregulars
## 1 Introduction
Models of interacting galaxies have shown that tidal forces in the interaction can produce long tails of stars and gas that have been pulled out of the interacting galaxies (Toomre & Toomre 1972). Zwicky (1956) pointed out the possibility that self-gravitating objects might develop in these tidal tails that could then evolve into small galaxies. Since then, concentrations of stars and gas that are probable “tidal dwarfs” in the making have been observed at the tips of tidal tails in several interacting systems (for example, the Antennae system; Mirabel et al. 1992). Numerical modelling confirms that bound, galaxy-sized clumps can form in tidal tails (Barnes & Hernquist 1992; Elmegreen, Kaufman, & Thomasson 1993).
These tidal dwarfs, once the tidal tail itself has dispersed and the parent galaxies have moved off, could bear a striking resemblance to small, independent, Im-type galaxies (Schweizer 1974). The tidal dwarfs are small, gas-rich, morphologically disorganized, and already have on-going star formation (see also Mirabel, Lutz, & Maza 1991). Furthermore, the properties measured for tidal dwarfs are well within the range of properties seen for normal, relatively isolated irregular galaxies (Mirabel, Dottori, & Lutz 1992; Duc & Mirabel 1994; Hibbard et al. 1994; Hunter 1997).
Gravitational interactions are an on-going process in the Universe that began when galaxies themselves first formed. Therefore, this mechanism for forming irregular galaxies has been taking place for the age of the Universe. Hunsberger, Charlton, & Zaritsky (1996), for example, estimate that as many as one-half of the current dwarf galaxy population of compact groups may have been formed from the interactions of giant spiral galaxies. The formation of dwarf irregulars in compact groups is accelerated because of the increased crowdedness and potential for interactions there. However, interactions can and do occur outside of compact groups of galaxies as well.
One must then consider that any given dwarf irregular galaxy, including field galaxies, could have been formed in one of two ways: traditional formation from collapse of a primordial cloud of gas early in the age of the Universe, and tidal dwarf formation from an interaction of larger galaxies at any time during the history of the Universe. Because the time scale since the formation of a tidal dwarf can be large, a tidal dwarf could appear to be relatively isolated if the formation took place many Gyrs ago.
Because the formation mechanism of traditional dwarfs and of tidal dwarfs are different, some key characteristics of these two groups of galaxies could also be different, as outlined by Barnes & Hernquist (1992) and Elmegreen et al. (1993). In this paper we examine a sample of irregular galaxies with these observational differences in mind and ask whether any seemingly normal Im galaxy might be a candidate for an old tidal dwarf. The observational characteristics of Im galaxies and distinctions with tidal dwarfs are too imprecise at this time to do more than point out candidate tidal dwarf systems, but it is a way to begin thinking about the issue that not all irregulars may have had the same initial conditions.
Dwarf galaxies, galaxies that are intrinsically not luminous, come in a variety of types, including spirals, irregulars, ellipticals, and spheroidals, and Gerola, Carnevali, & Salpeter (1983) suggested galaxy interactions as a means for forming dwarf ellipticals. However, the most common type of dwarf is a gas-rich Im galaxy, and we will primarily concentrate on irregulars in this discussion.
## 2 Selection Criteria for Old Tidal Dwarfs
Since tidal dwarfs form from material drawn out of larger galaxies, their properties could differ significantly from those of traditional dwarf galaxies. Barnes & Hernquist (1992) suggest that tidal dwarfs formed from parent galaxies with dark matter distributed in extended massive haloes (Barnes 1992) will have $`<`$5% of the dwarf’s mass in dark matter. Thus, tidal dwarfs should have low mass-to-light ratios. Tidal dwarfs may also have unusual metal abundances since they were made from material first processed in a much larger galaxy. Elmegreen et al. (1993) have also suggested that the age distribution of their stellar populations may also be peculiar since there will be old stars from the spiral along with stars that have formed in a burst of star formation when the dwarf first formed and stars that have formed at a much slower pace since then.
Some of these properties, like the amount of dark matter, will not depend on how long ago the tidal dwarf formed but others will. Tidal dwarfs forming out of modern day spirals are observed to be metal-rich because they have formed out of material already processed in a giant spiral (Duc & Mirabel 1999, Hunsberger et al. 2000); most dwarf irregulars today, on the other hand, are metal poor compared to all but the extreme outer parts of spiral disks. However, a tidal dwarf that formed from a spiral many Gyrs ago when spirals were still metal poor ($``$6 Gyrs ago for material drawn from the central regions of a galaxy like M33, $``$9 Gyrs ago for a galaxy like the Milky Way \[Mollá, Ferrini, & Díaz 1997\]). could have a metallicity today that is consistent with that of traditional, metal-poor dwarfs.
Another problem comes in disentangling evolutionary effects from initial conditions. For example, a galaxy which underwent a strong burst of star formation several Gyrs ago might be a tidal dwarf that formed then or it could be a tiny galaxy with a peculiar star formation history. We do know that many irregulars evolve with star formation rates that vary in amplitude by factors of a few, as would be expected in small galaxies (Ferraro et al. 1989; Tosi et al. 1991; Greggio et al. 1993; Marconi et al. 1995; Gallart et al. 1996a,b; Aparicio et al. 1997a,b; Dohm-Palmer et al. 1998, Gallagher et al. 1998, Gallart et al. 1999). However, others show evidence of higher amplitude variations either currently or in the past, some of which may be statistically significant (Israel 1988, Tolstoy 1996, Dohm-Palmer et al. 1997, Greggio et al. 1998, Tolstoy et al. 1998). Why statistically large variations occur in some seemingly isolated irregulars and not in others is not clear although arguments have been made that tiny galaxies should evolve via periodic starbursts (see, for example, Gerola, Seiden, & Schulman 1980). Clearly, we do not understand how irregular galaxies evolve well enough yet to be able to say whether a particular evolutionary history can only be consistent with a tidal dwarf formation scenario.
With these problems in mind, we examine properties of tidal dwarfs that are potentially different from those of traditional dwarfs. We are asking the question: Do any normal Im galaxies that we know of have properties that are consistent with a tidal dwarf origin?
Tidal dwarfs observed currently at the ends of tidal tails have M<sub>B</sub> of $``$14 to $``$19 (Mirabel et al. 1992; Duc & Mirabel 1994; Hibbard et al. 1994), and a luminosity function of tidal dwarfs in compact groups extends to as bright as $``$18 in M<sub>R</sub> (Hunsberger et al. 1996, converted to H$`{}_{0}{}^{}=65`$ km s<sup>-1</sup> Mpc<sup>-1</sup>). An older tidal dwarf, however, might have faded from the glory days of formation. Therefore, we expect to find tidal dwarfs with M$`{}_{B}{}^{}>18`$. In Table 1 we keep a running list of dwarfs that stand out in the properties discussed below as possible tidal dwarf candidates.
### 2.1 Lack of Dark Matter
Barnes & Hernquist (1992) found that two interacting galaxies with extended massive dark haloes produced tidal dwarfs with $`<`$5% of its mass in dark matter. (How this property might be effected by a different distribution of dark matter in the parent galaxies is not explored.) This lack of dark matter in tidal dwarfs is potentially the most useful distinguishing feature since that property will not change with time and it does not depend on when the dwarf formed. Most disk galaxies have rotation curves that require the presence of dark matter, and some studies of irregulars have argued that irregulars are just as dominated by, or even more dominated by, dark matter than spirals. Here we look for signs of unusual dark matter properties among irregulars through rotation curves and the traditional Tully-Fisher relationship.
Rotation curves of irregular galaxies are a mixed bag. There are those that look like normal disk rotation curves: they rise as a solid body, peak, and level off or even begin to fall (for example, DDO 154: Carignan & Purton 1998). In other irregulars the rotation curve rises but never peaks, presumably because the rotation curves have not been observed far enough out.
However, there are also irregulars that have been found to have no measureable rotation with upper limits on V<sub>rot,max</sub>sin $`i`$ of $``$7.5 km s<sup>-1</sup>. These include DDO 69 ($`=`$Leo A; Young & Lo 1996), DDO 99, DDO 120, DDO 143 (Swaters 1999), DDO 155 ($`=`$GR 8; Lo, Sargent, & Young 1993; but see also Carignan, Beaulieu, & Freeman who interpret the velocity field as having some rotation in the inner 250 pc); DDO 187 (Swaters 1999), DDO 210, DDO 216 ($`=`$Pegasus Dwarf), LGS3 (Lo, Sargent, & Young 1993); Sag DIG (Young & Lo 1997), and NGC 4163 (Swaters 1999). This is in contrast to low surface brightness spirals that, in spite of their very low surface brightness levels, nevertheless, rotate at high speeds compared to irregulars (de Blok, McGaugh, & van der Hulst 1996). Galaxies with no measureable organized rotation may be good candidates for no dark matter.
The maximum rotation speeds of irregulars and a sample of spirals taken from Broeils (1992) are shown in Figure 3. Upper limits to the rotation speed are used to put the galaxies with no measureable rotation on this plot. One can see that at an M<sub>B</sub> of about $``$15 and fainter the irregular galaxies deviate strongly from the relationship between M<sub>B</sub> and V<sub>rot,max</sub> determined for spirals. The deviation from the relationship is in the sense that many irregulars have rotation speeds that are too small for their luminosity. This is in the sense that one expects for galaxies that have too little dark matter.
In a study of dark matter in late-type dwarf galaxies, Swaters (1999) has shown that there may not be a lot of dark matter in the disks of irregulars, but dark matter in haloes is still required to explain the rotation curves. However, in some galaxies the evidence for any dark matter at all is not strong (for example, NGC 1569: Stil 1999). Swaters (1999) found 5 galaxies in his study, including 3 Ims (DDO 50, DDO 125, DDO 143), for which a maximum disk fit to the rotation curve leaves no room for dark matter. He suggests that one of these galaxies, DDO 125, is a good candidate for a tidal dwarf since it lies near giant HI streamers associated with the larger irregular galaxy NGC 4449 (Hunter et al. 1998).
In Figure 3 we consider an alternative means of looking at the mass, and hence dark matter, in galaxies: the traditional Tully-Fisher relationship (Tully & Fisher 1977). The Tully-Fisher relationship is shown in Figure 3a where we plot log W$`{}_{}{}^{ci}{}_{20}{}^{}`$ against M$`{}_{}{}^{i}{}_{B}{}^{}`$. W$`{}_{}{}^{ci}{}_{20}{}^{}`$ is the full width at 20% intensity of the integrated HI profile, corrected for instrumental broadening, random motions, and the inclination of the galaxy (Broeils 1992). The inclinations were determined using minor-to-major axis ratios from Swaters (1999) and de Vaucouleurs et al. (1991,$`=`$RC3) and assuming an intrinsic ratio of 0.3 (Hodge & Hitchcock 1966, van den Bergh 1988). The M$`{}_{}{}^{i}{}_{B}{}^{}`$ have been corrected for internal absorption, also dependent on the galaxy inclination, according to Broeils (1992) with a reduction of a factor of 4 in the absorption compared to spirals to better match the observed lower dust contents of irregulars. W$`{}_{}{}^{ci}{}_{20}{}^{}`$ should measure approximately twice the maximum rotation speed and so should be related to the total mass in the galaxy. This plot has the advantage over the plot shown in Figure 3 that integrated HI profiles are available for far more galaxies than are velocity fields. We see that there is much more scatter among the irregular galaxies than the spiral sample, but the scatter for the brighter irregulars is distributed around the relationship defined by the spirals with more falling below the relationship than above it. However, at the low end of the relationship—low HI widths and low luminosities, several irregulars are too bright for their HI widths. This deviation above the relationship is consistent with the possibility that those galaxies have less dark matter than other galaxies of that luminosity although statistics on low luminosity galaxies are poor.
Carignan & Beaulieu (1989), Swaters (1999), and McGaugh et al. (2000) have observed that the Tully-Fisher relationship begins to break down for low luminosity galaxies although the galaxies in their samples deviate in the sense of being underluminous for their rotation speed. Milgrom & Braun (1988) suggest that the relationship is maintained if M$`{}_{}{}^{i}{}_{B}{}^{}`$ is replaced with the total luminous mass, including gas. We have examined this possibility in Figure 3b. We have estimated the mass in stars from M$`{}_{}{}^{i}{}_{B}{}^{}`$ and a stellar mass-to-light ratio that is appropriate for a galaxy forming stars at a constant rate for 10 Gyr with a normal stellar initial mass function (Larson & Tinsley 1978). This is clearly a rough approximation and will not apply equally well to all of the galaxies in the sample. However, this is most likely to break down for the dwarf galaxies for which errors of factors of even 10 in stellar mass will not make very much difference to the mass in gas plus stars because the masses are usually dominated by the gas. The mass in gas is HI plus He (1.34$`\times `$M<sub>HI</sub>) and does not include molecular gas since this quantity is not known for most irregulars. The resulting relationship looks similar to that in Figure 3a, and, if one compares equal logarithmic intervals, the scatter is no less. This is in contrast to what McGaugh et al. (2000) found: that galaxies with W$`{}_{20}{}^{ci}<180`$ km s<sup>-1</sup> in their sample, that extends to a mass of $`10^7`$ M$`_{\mathrm{}}`$, are brought into agreement with spirals once the gas content is taken into account.
In Figure 3b we see that only a few galaxies deviate more than most irregulars. Several, such as DDO 155, NGC 4163, and the Sm galaxy DDO 135, fall below the relationship for the rest of the galaxies. This is in the opposite sense of what we are looking for. Of those galaxies above the relationship, only IC 1613, DDO 210, DDO 216, and SagDIG, because they are upper limits (all but IC 1613 have no measureable rotation), have the potential to fall further outside the relationship than the bulk of the galaxies. (We ignore the spiral DDO 80, possibly interacting, discussed by Broeils .)
### 2.2 Metallicity
At the time of formation the tidal dwarf will take on the metallicity of the material drawn from the parent spiral, and, since spirals have evolved more rapidly than irregulars, this could make the tidal dwarf more metal rich (Schweizer 1978). In fact, many tidal dwarfs in formation observed today are too metal rich for their luminosity (Duc & Mirabel 1997, Hunsberger et al. 2000). For normal irregulars, Richer & McCall (1995) found that the scatter in the metallicity-luminosity relationship increases for galaxies with M$`{}_{B}{}^{}>15`$. This could imply that low luminosity galaxies have a more diverse evolutionary background. We have examined the position of irregular galaxies on a plot of oxygen abundance versus M<sub>B</sub>, adapted from Hunter & Hoffman (1999). In Figure 3 we show the oxygen abundance plotted against M<sub>B</sub> for a sample of spiral and irregular galaxies along with Richer and McCall’s relationship for low-luminosity galaxies. The scatter even among the spirals is substantial. Although most of the irregulars with no measureable rotation deviate substantially from the relationship, the scatter among all of the galaxies is too high to be able to say that they deviate more than most. Clearly the lower luminosity end of this relationship needs to be explored further.
We have also included a few Im galaxies in the Virgo cluster for comparison. Potentially the higher density of clusters like Virgo will result in a higher population of tidal dwarfs as has been found in compact clusters (Hunsberger et al. 1996). Vílchez (1995) had suggested that irregulars in Virgo are in fact more metal rich than field irregulars, but he assumed the high metallicity branch of the double-valued relationship between emission-line ratio and abundance. On the other hand, Lee, McCall, & Richer (1998) found metallicities for Virgo irregulars that are consistent with the relationship. In our plot, where we assumed the lower branch of the metallicity, the Virgo Cluster galaxies do not stand out from the general scatter.
However, there are also complications to using metalliciy as a tidal dwarf indicator: 1) If the material that forms the tidal dwarf comes primarily from the outer part of the spiral, it could be just as metal poor as an irregular. 2) If a tidal dwarf formed many Gyr ago, the starting metallicity would be lower than if it had formed today. 3) The metallicity will change as the dwarf evolves and how it changes is convolved with how it evolves. 4) There are substantial observational uncertainties in determining the oxygen abundance.
### 2.3 Structure
Irregular galaxies are generally believed to be disk systems, although thicker than spirals Hodge & Hitchcock 1966, van den Bergh 1988). However, there is some debate about even this basic property of irregulars (Sung et al. 1998). However, we see no reason why a dwarf formed in a tidal tail would have to be a disk. Furthermore, the irregulars that have no measureable disk rotation could be ones that are not disk-shaped. Intriguingly, Patterson & Thuan (1996) examined the surface photometry and scale lengths of a sample of dwarf irregular galaxies and found that they divided into two groups. One group has scale lengths like those of dwarf ellipticals and twice those of BCDs, and the other is comparable to BCDs and half that of dEs. Could these two groups also be related to the two origins? At this point, we cannot tell which group would be the tidal dwarfs. In addition, studies of the irregulars WLM and NGC 3109 have shown that those irregulars have a halo of old stars in addition to their disk (Minniti & Zijlstra 1996; Minniti, Zijlstra, & Alonso 1999). They point out that by contrast all of the old stars in the LMC are in its disk. However, whether this is normal or abnormal for irregulars is not yet clear.
### 2.4 Stellar Population
Another possible oddity of a tidal dwarf is its stellar population. Elmegreen et al. (1993) argue that the tidal dwarf should consist of a small fraction ($``$40%) of old stars from the parent spiral, a strong starburst at formation, and a normal distribution of mixed ages of stars formed since the galaxy’s formation. Unfortunately, not many color magnitude diagrams of irregulars go deep enough for analysis to pull out limits on star formation histories more than a few Gyrs into the past. Of those that can put some limits up to 10 Gyrs ago, 6 galaxies appear to have normal star formation histories (DDO 216: Aparicio, Gallart, & Bertelli 1997a; LGS3: Aparicio et al. 1997b; IC 1613: Cole et al. 1999; NGC 3109: Minniti et al. 1999; NGC 6822: Gallart et al. 1996a; WLM: Minniti & Zijlstra 1997). However, one galaxy may fit the pattern expected of a tidal dwarf: DDO 69. Tolstoy et al. (1998) put a limit of $`<`$10% of the total in a very old stellar population, with the majority of the star formation taking place within the past 1.5 Gyr. In the scenario of Elmegreen et al., DDO 69 would have formed about 2 Gyrs ago.
## 3 Discussion
We have examined properties of a sample of irregular galaxies from the perspective of features that might distinguish galaxies formed in tidal interactions at some time shorter than a Hubble time from those formed from collapse of a primordial gas cloud a Hubble time ago. We have considered the lack of dark matter predicted by models as manifested in rotation speeds and the Tully-Fisher relationship, the fact that tidal dwarfs may have formed from enriched material, structure, and peculiar stellar populations. However, using these features to identify old tidal dwarfs is currently imprecise. Abundances and star formation histories are entangled in other evolutionary and observational effects, and not enough is known about the amount and location of dark matter and the true structure of irregulars.
Nevertheless, we have identified candidates for old tidal dwarfs, and they are listed in Table 1. We have also listed the distance to the nearest large galaxy. A little over one-quarter of the galaxies in this list are in the Local Group. Eighty-five percent of the galaxies are within 0.5 Mpc of a large galaxy; and one lies near supergiant gas streamers wrapped around a nearby Im galaxy.
Because of the difficulties in identifying old tidal dwarfs, these galaxies can only be considered candidates at this point. In addition this is not an exhaustive list, and we have not included representative samples of other groups of dwarfs including dwarf ellipticals and dwarf spheroidals. The peculiar galaxy IZw18, for example, has the peculiar stellar populations and complex kinematics that might make it a candidate.
Clearly, it is important to understand the formation and evolutionary processes of the most common galaxy in the universe: irregular and other dwarf galaxies. The fact that irregulars could potentially be formed in more than one way complicates our ability to interpret the properties of the galaxies that we see today. How can we improve our understanding of irregulars so that differences due to different origins might be more apparent? We need to better understand the kinematics and structures of irregular and dwarf galaxies. This includes the gas and stellar kinematics and velocity dispersions from which we can infer the distribution and amount of dark matter and the stellar structure of the galaxy. We also need more very deep studies of stellar populations of irregulars, particularly probing the extremes of galaxy characteristics. Only once there is a statistically significant sample of star formation histories can we begin to see trends. Finally, we need numerical simulations that can show whether interactions are feasible, perhaps between two unequal mass partners, that can produce a tidal dwarf and still leave the parent spiral intact. This is particularly important for the Local Group system in which we identify 6 candidate old tidal dwarfs, but the obvious parents are relatively normal looking spirals.
This discussion originated while EWR was a summer student at Lowell Observatory under the Research Experiences for Undergraduates program operated by Northern Arizona University and funded by the National Science Foundation under grant number 9423921. Support for this work was provided by the Lowell Research Fund and the Friends of Lowell Observatory.
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# Fulfillment of expectations of precise measurements of the Casimir force
## I Introduction
The Casimir force (see for a review) between closely spaced macroscopic bodies is an effect of quantum electrodynamics (QED) and for this reason it could be predicted very accurately. The force acting between nonideal bodies can be found using the Lifshitz theory , where it depends on optical properties of used materials. Knowledge of these properties is the weakest element in the theory restricting the accuracy that can be achieved. Experiments measuring the Casimir force are of great importance because they are sensitive to the presence of new fundamental forces predicted in many modern theories (see and references therein). To distinguish a new force from the background, we should be able to calculate the Casimir force with a precision better than the experimental one. In the series of recent experiments this force has been measured with the torsion pendulum (TP) in the range of distances $`0.66\mu m`$ and with the atomic force microscope (AFM) in the range $`0.10.9\mu m`$. The corresponding precisions were 5% and 1%, respectively.
For two ideal plates the famous Casimir formula for the force per unit area is
$$F_c^{pl}(a)=\frac{\pi ^2\mathrm{}c}{240a^4},$$
(1)
where $`a`$ is the distance between plates. In the experiments the force is measured between metallized disc and sphere because for two plates it is difficult to keep them parallel. In this case (1) has to be modified with the proximity force theorem (PFT) . This theorem allows to evaluate the force by adding the contributions of various distances as if they were independent and for plate and sphere it is reduced to
$$F(a)=2\pi R\underset{a}{\overset{R+a}{}}F^{pl}\left(x\right)𝑑x,$$
(2)
where $`R`$ is the radius of curvature of the spherical surface. The PFT approximation is good for $`Ra`$ that holds true in all the experiments. If we use the Casimir expression (1) for the force in (2), then the force between plate and sphere will be
$$F_c^0(a)=\frac{\pi ^3\mathrm{}c}{360}\frac{R}{a^3}.$$
(3)
Eq. (3) was deduced for ideally conducting bodies at zero temperature and three kinds of corrections have been considered to take into account their real properties. The correction due to finite metal conductivity was found on the base of the free electron model, where the optical properties of a metal were described by the only parameter $`\omega _p`$ which is the plasma frequency. The force including corrections up to the second order is
$$F_c^p\left(a\right)=F_c^0\left(a\right)\left[14\frac{c}{a\omega _p}+\frac{72}{5}\left(\frac{c}{a\omega _p}\right)^2\right].$$
(4)
For typical plasma frequency $`\omega _p10^{16}s^1`$ and separations $`a1\mu m`$ the correction will be more than 10%. Correction due to finite temperature has been found for ideal conductors and the resulting force is given by
$$F_c^T\left(a\right)=F_c^0\left(a\right)\left(1+\frac{720}{\pi ^2}f\left(\xi \right)\right),$$
(5)
where $`\xi =k_BTa/\mathrm{}ca(\mu m)/7.61`$ for $`T=300^{}K`$. The function $`f(\xi )`$ is expressed via an infinite sum but it can be represented approximately as $`f(\xi )=(\xi ^3/2\pi )\zeta (3)(\xi ^4\pi ^2/45)`$ for $`\xi <1/2`$. The temperature correction is negligible for the AFM experiments since $`\xi `$ is small in the important separation range $`0.10.3\mu m`$ and is only a minor correction in condition of the TP experiment , where the important separation range was $`0.63\mu m`$. The general form of the correction due to surface distortions has been found in . If the bodies are covered by distortions with characteristic amplitudes $`A_1`$ and $`A_2`$, then the force up to the second order in the relative amplitudes of the distortions has the form
$$F_c^d\left(a\right)=F_c^0\left(a\right)[1+3(f_1\frac{A_1}{a}f_2\frac{A_2}{a})+$$
$$6(f_1^2\frac{A_1^2}{a^2}2f_1f_2\frac{A_1A_2}{a^2}+f_2^2\frac{A_2^2}{a^2})],$$
(6)
where the functions $`f_{1,2}(x,y)`$ describe distribution of the distortions on the surfaces and $`\mathrm{}`$ denotes averaging over the surface area. Corrections due to surface roughness are very important for the experiment .
At first the experimental data were treated using these corrections to Eq. (3), but it was realized soon that at least the conductivity correction has to be considered on more reliable basis. In more realistic approach the Lifshitz theory was used to evaluate the force between bodies . Similar but technically a little bit different method was developed in . In these approaches the force depends on the dielectric function of the bodies at imaginary frequencies $`\epsilon \left(i\omega \right)`$. It has to be expressed with the dispersion relation via the imaginary part of the function $`\epsilon \left(\omega \right)`$ on the real axis which can be directly measured. However, in any of the experiments the information on $`\epsilon \left(\omega \right)`$ was not collected and the handbook data were used instead. Such data are good only to make an estimate for the Casimir force with the accuracy much worse than the experimental one. The reason is that the dielectric function depends in substantial degree on the sample preparation procedure as will be discussed below. Nevertheless, it is possible to find a reliable upper limit on the Casimir force using only well defined parameters of perfect crystalline materials. In this paper we will discuss in detail this limit and its comparison with the existing experimental data.
The paper is organized as follows. In Sec. II we give a general expression for the Casimir force between sphere and plate made of nonideal materials at nonzero temperature. Then, to treat the experimental data, the expression for the force is generalized for the case of layered bodies. The choice of dielectric functions and parameters for the used materials is described in Sec. III. In Sec. IV we define the boundary values of the optical parameters and find the upper limit on the force in conditions of each independent experiment. Possible reasons for discrepancy between theory and experiment are discussed in Sec. V. Our conclusions are given in the last Section.
## II Theory
Let us discuss first a reliable way to evaluate the Casimir force in the experimental configurations. The force per unit area between parallel plates arising as a result of electromagnetic fluctuations is generalized by the Lifshitz theory , where the real material is taken into account by its dielectric function at imaginary frequencies $`\epsilon \left(i\zeta \right)`$:
$$F^{pl}(a)=\frac{kT}{\pi c^3}\underset{n=0}{\overset{\mathrm{}}{}}^{}\zeta _n^3\underset{1}{\overset{\mathrm{}}{}}𝑑pp^2\left\{\left[G_1^2e^{2p\zeta _na/c}1\right]^1+\left[G_2^2e^{2p\zeta _na/c}1\right]^1\right\}.$$
(7)
Here prime over the sum sign means that $`n=0`$ term is taken with the coefficient $`1/2`$ and
$$G_1=\frac{p+s}{ps},G_2=\frac{\epsilon \left(i\zeta _n\right)p+s}{\epsilon \left(i\zeta _n\right)ps},$$
$$s=\sqrt{\epsilon \left(i\zeta _n\right)1+p^2},\zeta _n=\frac{2\pi nkT}{\mathrm{}}.$$
(8)
It is supposed that both bodies were made of identical materials. The function $`\epsilon \left(i\zeta _n\right)`$ cannot be measured directly but can be expressed via imaginary part of the dielectric function $`\epsilon ^{\prime \prime }\left(\omega \right)`$ on the real axis with the dispersion relation
$$\epsilon \left(i\zeta \right)1=\frac{2}{\pi }\underset{0}{\overset{\mathrm{}}{}}𝑑\omega \frac{\omega \epsilon ^{\prime \prime }\left(\omega \right)}{\omega ^2+\zeta ^2}.$$
(9)
Information on $`\epsilon ^{\prime \prime }\left(\omega \right)`$ can be extracted from the data on reflectivity and absorptivity of electromagnetic waves with the frequency $`\omega `$ for a given material.
Applying PFT to Eq. (7) one can find the force between sphere and plate. The integration in (2) can be done analytically and we find
$$F(a)=\frac{kTR}{c^2}\underset{n=0}{\overset{\mathrm{}}{}}^{}\zeta _n^2\underset{1}{\overset{\mathrm{}}{}}𝑑pp\mathrm{ln}\left[\left(G_1^2e^{2p\zeta _na/c}1\right)\left(G_2^2e^{2p\zeta _na/c}1\right)\right].$$
(10)
Special care needs to treat the first $`n=0`$ term. The formal reason is that $`\zeta _n^2`$ becomes zero but the integral over $`p`$ diverges. The physical reason is that this term corresponds to the static limit when for metallic bodies $`\epsilon \mathrm{}`$. This means that any parameter characterizing the dielectric function of a metal cannot appear in the $`n=0`$ term in contrast with a dielectric for which it will depend on the static permittivity of the material. In the $`\epsilon \mathrm{}`$ limit the functions $`G_{1,2}`$ become $`G_1=G_2=1`$. The formal problem is overcome by introducing the integration over a new variable $`x=2p\zeta _na/c`$ and after that one can take $`\zeta _n=0`$ for the $`n=0`$ term. The resulting contribution of the first term in the force corresponds to the classical limit $`F_{cl}\left(a\right)`$ for metals
$$F_{cl}\left(a\right)=\frac{kTR}{4a^2}\zeta \left(3\right),$$
(11)
where $`\zeta \left(n\right)`$ is the zeta-function. Note that in this limit the force does not depend on the metal parameters as it should be for a static field.
The bare Casimir force (3) is reproduced from Eq. (10) in the limit $`\epsilon \mathrm{}`$ and $`T0`$. The finite conductivity correction also can be derived from (10). To this end one considers the limit of small temperature when the sum in (10) can be replaced by the integral and supposes that the dielectric function of the metal covering the bodies is described by the free electron plasma model. In this model $`\epsilon \left(i\zeta \right)`$ is
$$\epsilon \left(i\zeta \right)=1+\frac{\omega _p^2}{\zeta ^2},$$
(12)
where $`\omega _p`$ is the free electron plasma frequency. Typical value of the frequency $`\omega _p10^{16}s^1`$ is larger than fluctuation frequencies $`\zeta c/a`$ giving the main contribution in (10). Then one can expand the functions $`G_{1,2}`$ in (10) in powers of the parameter $`\zeta /\omega _p`$ and performing necessary integrations one finds exactly the result (4) for the conductivity corrections <sup>1</sup><sup>1</sup>1The correction is actually connected with finite density of free electrons (finite $`\omega _p`$) since the metal conductivity is still infinite for the plasma model. Nevertheless, we will not change the fixed terminology.. In this way the corrections up to the fourth order were found in recent paper . The temperature correction (5) is also reproduced from (10) in the limit of ideal metals $`\epsilon \mathrm{}`$. In this case the linear in temperature correction does not survive since the $`n=0`$ term (11) is exactly canceled by the linear in $`T`$ contribution from the rest terms in the sum. As a result the leading correction behaves only as $`\xi ^3`$.
The expression (10) differs from those used in and in two respects. First, in the cited papers the integration connected with the PFT was not done analytically that complicated numerical analysis. Second, the zero temperature limit has been taken. This limit was also considered in , though the PFT integral was evaluated explicitly. It seems a reasonable approximation at small separations because the temperature correction in (5) is proportional to $`\xi ^3`$ and, therefore, is small. However, one should remember that this correction was derived in the limit of ideal conductors $`\epsilon \mathrm{}`$. For a real conductor it will be proportional to $`\xi `$ as expected for difference between sum and integral and will be important (for details see ). We have computed the force according to (10) and with the integral instead of the sum at the smallest separation $`a=100nm`$ tested in the experiments. For the plasma model (12) with $`\omega _p=210^{16}s^1`$ we have found that the difference between the sum and integral is $`2.5pN`$ for $`T=300^{}K`$. It becomes $`4pN`$ for the Drude dielectric function (see Eq. (18) below) with the damping frequency $`\omega _\tau =510^{13}s^1`$. These values exceed the conservative estimate for the experimental errors $`2pN`$ .
In the AFM experiments an additional $`Au_{0.6}Pd_{0.4}`$ layer of $`20nm`$ or $`8nm`$ thick was on the top of $`Al`$ metallization of the bodies to prevent aluminum oxidation. It has to be included into consideration. This layer is transparent for the electromagnetic waves with high frequencies $`c/a`$ since adsorption, proportional to $`\epsilon ^{\prime \prime }\left(\omega \right)`$, is small. For this reason the layer was ignored in . However, the force depends on $`\epsilon (i\zeta )`$ for which the low frequencies dominate in the dispersion relation (9) because of large $`\epsilon ^{\prime \prime }\left(\omega \right)`$ and that is why we cannot neglect the $`Au/Pd`$ layer. To take it into account, one has to generalize expression for the force (7) to the case of layered bodies. Suppose that the top layer has the thickness $`h`$ and its dielectric function is $`\epsilon _1`$. The bottom layer is thick enough to be considered as infinite and let its dielectric function be $`\epsilon _2`$. The method described in for deriving Eq. (7) can be easily generalized for layered plates. We have to add only the matching conditions for the Green functions on the layers interface. After some algebra the result will look exactly as (7) but with more complex $`G_{1,2}`$:
$$G_1=\frac{\left(s_1+s_2\right)\left(p+s_1\right)e^{\zeta _ns_1h/c}+\left(s_1s_2\right)\left(ps_1\right)e^{\zeta _ns_1h/c}}{\left(s_1+s_2\right)\left(ps_1\right)e^{\zeta _ns_1h/c}+\left(s_1s_2\right)\left(p+s_1\right)e^{\zeta _ns_1h/c}},$$
$$G_2=\frac{\left(\epsilon _2s_1+\epsilon _1s_2\right)\left(\epsilon _1p+s_1\right)e^{\zeta _ns_1h/c}+\left(\epsilon _2s_1\epsilon _1s_2\right)\left(\epsilon _1ps_1\right)e^{\zeta _ns_1h/c}}{\left(\epsilon _2s_1+\epsilon _1s_2\right)\left(\epsilon _1ps_1\right)e^{\zeta _ns_1h/c}+\left(\epsilon _2s_1\epsilon _1s_2\right)\left(\epsilon _1p+s_1\right)e^{\zeta _ns_1h/c}},$$
(13)
where $`s_{1,2}`$ are defined similar to $`s`$ in (8). The force between plate and sphere is given by (10) with the above $`G_{1,2}`$.
To see qualitatively the effect of an additional layer, we found the finite conductivity correction up to the second order in this case. Than for the force one has
$$F_c^p(a,h)=F_c^0\left(a\right)\left[14K(h)\frac{c}{a\omega _{1p}}+\frac{72}{5}\left(K(h)\frac{c}{a\omega _{1p}}\right)^2\right],$$
(14)
where the function $`K(h)`$ depends on the plasma frequencies of the layers $`\omega _{1p}`$ , $`\omega _{2p}`$ and the thickness of the top layer $`h`$
$$K(h)=\frac{\omega _{1p}+\omega _{2p}\mathrm{tanh}\left(h\omega _{1p}/c\right)}{\omega _{2p}+\omega _{1p}\mathrm{tanh}\left(h\omega _{1p}/c\right)}.$$
(15)
When $`h=0`$ the force will depend only on $`\omega _{2p}`$ and in the case $`h\mathrm{}`$ on $`\omega _{1p}`$ as it should be. The effect of the top layer disappears if the plasma frequencies coincide. The top layer will be negligible if $`h\omega _{1p}/c1`$. For typical plasma frequencies $`10^{16}s^1`$ it is definitely not the case even for $`h=8nm`$. The opposite conclusion made in was based on the too small value of $`\omega _p`$ for gold as will be discussed below (see also ). Eq. (14) is not very good approximation and was discussed only for qualitative understanding of the effect. For actual calculations we will use the exact equations (10), (13).
Importance of a thin metallic layer on the body surfaces has been stressed first in . The general expression for the Casimir force between layered bodies has been presented in but was not used their for actual calculations. Significant role of the $`Au/Pd`$ layer in the AFM experiments was indicated in our preprint , where it was demonstrated that the effect far exceeds the experimental errors. This conclusion was supported in , where the expressions (13) for $`G_{1,2}`$ were confirmed using a different method to deduce them. However, the authors were uncertain on applicability of (13) for thin films with $`h<25nm`$ because the spatial dispersion of the dielectric function can be important for such films. We discuss this effect in Sec. V where we argue that the spatial dispersion can be neglected because of very short mean free path for the electrons in thin films.
## III The dielectric function
Now we are able to evaluate the Casimir force in real geometry of the experiments if there is information on the dielectric functions of used materials: $`Au`$, $`Al`$, and $`Au_{0.6}Pd_{0.4}`$ alloy. Strictly speaking, one has to measure these functions in wide range of wavelengths on the same samples which are used for the force measurement. It was not done in all of the experiments and to draw any conclusion from them we have to make some assumptions on the dielectric functions. At low frequencies $`Au`$ and $`Al`$ are well described by the Drude dielectric function :
$$\epsilon =\epsilon ^{}+i\epsilon ^{\prime \prime },$$
$$\epsilon ^{}\left(\omega \right)=1\frac{\omega _p^2}{\omega ^2+\omega _\tau ^2},\epsilon ^{\prime \prime }\left(\omega \right)=\frac{\omega _p^2\omega _\tau }{\omega \left(\omega ^2+\omega _\tau ^2\right)},$$
(16)
where $`\omega _p`$ is the free electron plasma frequency and $`\omega _\tau `$ is the Drude damping frequency. A simple test for validity of the Drude model is behavior of the material resistivity which is defined as
$$\rho \left(\omega \right)=Im\left(\frac{1}{\epsilon _0\left(1\epsilon \left(\omega \right)\right)\omega }\right)=\frac{\omega _\tau }{\epsilon _0\omega _p^2},$$
(17)
where $`\epsilon _0`$ is the free space permittivity. The resistivity is frequency independent within the Drude approximation. For high-purity single-crystal samples of $`Au`$ and $`Al`$ (entries 2 in Table I) the frequency behavior of the resistivity in the infrared range of wavelengths $`3\mu m<\lambda <32\mu m`$ is shown in Fig. 1. The data on the dielectric functions were taken from , where the results from many original works are collected. The data for $`\epsilon ^{}\left(\omega \right)`$ and $`\epsilon ^{\prime \prime }\left(\omega \right)`$ can be fitted with (16) to find the parameters $`\omega _p`$ and $`\omega _\tau `$. The points and fitting curves for $`\epsilon ^{\prime \prime }\left(\omega \right)`$ are shown in the same Fig. 1. Palladium definitely cannot be described by (16) since its resistance significantly changes in the infrared range. However, it is known experimentally that amorphous metallic alloys can be described by the Drude approximation . The physical explanation for this is associated with large Drude damping of the compounds like $`Au_{0.6}Pd_{0.4}`$.
Of course, at higher frequencies when interband transitions are reached the Drude approximation fails. Nevertheless, it is very helpful since low frequencies dominate in the dispersion relation. Extrapolating (16) to all frequencies one finds
$$\epsilon \left(i\zeta \right)=1+\frac{\omega _p^2}{\zeta \left(\zeta +\omega _\tau \right)}.$$
(18)
Let us estimate the relative error inserted in (18) due to extrapolation. If $`\omega _0`$ is the frequency of the first resonance for a given metal, then the contribution in $`\epsilon \left(i\zeta \right)`$ of the frequency range $`\omega _0<\omega <\mathrm{}`$, where the Drude model does not valid, will be $`\left(\omega _p/\zeta \right)^2`$ $`\left(\omega _\tau /\omega _0\right)`$ for $`\zeta \omega _0`$. This contribution one can take as an estimate for the absolute error and, therefore, for the relative error one has $`\omega _\tau /\omega _0`$. For typical values $`\omega _\tau 10^{14}s^1`$ and $`\omega _0>10^{15}s^1`$ the error can be as large as 10% but error in the force is smaller. If we will use (18) for the force computation and change $`\omega _p`$ by 5% ( 10% correction to $`\epsilon \left(i\zeta \right)`$ at all frequencies) then the force is changed less than 2%. Moreover, since the interband transitions give a correction to (18) which is frequency dependent, it reduces the correction to the force further. Of course, we can take the interband transitions into consideration exactly using the handbook data in visible-ultraviolet range which are not very sensitive to the purity and defect density as it happens in the infrared range. However, we are intended to establish the upper limit on the force using $`\omega _p`$ in (18) which is definitely larger than any real value and for this reason we can neglect the interband transitions.
Therefore, in all cases of interest we can use Eq. (18) to describe the dielectric function of a material on the imaginary axis. The question is how we should extract the parameters $`\omega _p`$ and $`\omega _\tau `$ from the data. We proceeded as follows. The data for the complex refraction index $`n+i\kappa =\sqrt{\epsilon }`$ have been taken from . First, the validity of the Drude approximation was checked by calculating the frequency dependence of the material resistivity according to (17). In the investigated cases the resistivity is more or less constant in the wavelength range $`\lambda >2\mu m`$ ($`\omega <9.410^{14}s^1`$). This range gives the most important contribution to the dispersion relation (9) and, therefore, it is the range where we have to extract the Drude parameters. Using $`\omega _p`$ found from the high frequency region can happen to be wrong. For example, sometime the plasma frequency is estimated using the transition point in the reflectivity dependence on frequency. It works, not very good though, for $`Al`$ but gives considerably smaller value for $`Au`$ than that found from fitting $`\epsilon \left(\omega \right)`$ in the infrared range. Probably such an estimate was taken in for $`Au`$ where very small value $`\omega _p=3.610^{15}s^1`$ was used.
The optimal fitting procedure is described in . The damping frequency is evaluated first from the ratio $`\left(1\epsilon ^{}\right)/\epsilon ^{\prime \prime }`$ which depends linearly on frequency in the Drude model
$$\frac{1\epsilon ^{}}{\epsilon ^{\prime \prime }}=\frac{\omega }{\omega _\tau }.$$
(19)
After that $`\omega _p^2`$ can be extracted from $`1\epsilon ^{}`$ by linear fit. The results together with the statistical errors are collected in Table I for those data in which include the optical behavior of $`Au`$ and $`Al`$ in the infrared region. This table clearly demonstrates that the Drude parameters depend significantly on the sample which is used to measure the optical data. These samples contained different densities of the defects (such as impurity atoms, vacancies, dislocations, etc.) that influence their optical properties. In this sense there are no universal material parameters. Reproducible parameters one can get only for high-purity single-crystals. In this connection all the attempts to use the handbook data for the Casimir force calculation can be considered only as estimates and cannot claim on high precision.
Actually in any of the experiment we do not know the Drude parameters even with 10% accuracy. That is because the optical properties of evaporated or spattered films which cover the bodies can be quite different from those of bulk materials and depend on technological details of film preparation. It is known, for example, that the film density is typically 0.7 from that of the bulk material if it was not annealed. For the resistivity of spattered and evaporated $`Au`$ the value $`\rho _0=8.2\mu \mathrm{\Omega }cm`$ has been reported in contrast with the bulk resistivity $`2.25\mu \mathrm{\Omega }cm`$. If a metal is evaporated or spattered on a substrate, it has a large number of defects. Relatively thick metallic films ($`>100nm`$) are usually exist in polycrystalline form. Defects will reduce the concentration of free electrons $`n`$ which defines the plasma frequency of the material. They also will increase the damping frequency $`\omega _\tau `$ and resistivity since the mean free path of electrons will shorten. To minimize these undesirable in practical applications effects the films are usually annealed at high temperature. In the experiments it was not reported were the bodies annealed or not but one can say definitely that it was not done in the AFM experiments because the polysterene ball cannot exist at the annealing temperature. Even more defects present in thin films ($`2030nm`$) which are usually amorphous. This explains why thin films have very large resistivity in comparison with the bulk material. Entries 1 and 4 for $`Al`$ in Table I correspond to the data for thick film samples. They support our expectations that the plasma frequency for films should be smaller and the resistivity larger than those for the bulk material.
## IV The upper limit
Though we cannot use the handbook data to evaluate the force, one can constrain it for a given experiment. This statement is based on the observation that the force (10) increases every time when $`\omega _p`$ increases or $`\omega _\tau `$ decreases. It has simple physical meaning: the force becomes larger when the metal reflectivity increases. For us it is important that any technological procedures will reduce $`\omega _p`$ and increase $`\omega _\tau `$ for a given material. A perfect single-crystal will have the largest plasma frequency and the smallest $`\omega _\tau `$ and these parameters are well defined. One can use them to get the upper limit on the Casimir force. The plasma frequency $`\omega _p`$ is defined by the concentration of free electrons in the metal $`n`$ and their effective mass $`m_e^{}`$
$$\omega _p=\sqrt{\frac{e^2n}{m_e^{}\epsilon _0}},$$
(20)
where $`e`$ is the electron charge. For good metals, which we are concerned, $`m_e^{}`$ is close but larger than the electron mass. It will be helpful for what follows to use Eq. (17) and instead of the damping frequency $`\omega _\tau `$ take the static resistivity $`\rho \left(0\right)=\rho _0`$ as a parameter. The later can be directly measured for any material.
### A Torsion pendulum experiment
In the TP experiment the quartz lens and plate were covered first with $`Cu`$ of thickness $`0.5\mu m`$ and then with $`Au`$ of the same thickness. The $`Au`$ layer is thick enough to be considered as infinite and $`Cu`$ will not play any role. We will find the upper limit on the electron concentration if suppose that every $`Au`$ atom produce a free electron with the mass $`m_e^{}=m_e`$. Then for the $`Au`$ plasma frequency one finds $`\omega _p^{Au}=1.3710^{16}s^1`$. The resistivity for crystalline gold is $`\rho _0^{Au}=2.25\mu \mathrm{\Omega }cm`$. One can compare these parameters with that given in Table I to make sure that they correspond to the limit values. Substituting these parameters in (18) and calculating the force according to (10) one finds the upper limit on the Casimir force $`F^{max}\left(a\right)`$ in the TP experiment.
To compare the upper limit on the force with the measured force $`F^{exp}`$, it is more convenient to consider the residual force defined as<sup>*</sup><sup>*</sup>*Note that Lamoreaux used different definition of the residual force $`F^{exp}\left(a_i\right)F_c^0\left(a_i\right)`$.
$$\mathrm{\Delta }F(a_i)=F^{exp}\left(a_i\right)F^{max}\left(a_i\right),$$
(21)
where $`a_i`$ are the separations for which the force has been measured. Theory and experiment will be in agreement if $`\mathrm{\Delta }F`$ will not be positive within the experimental errors. The original experimental data were presented for the lens curvature radius $`R=11.3cm`$ and the residual force in this case is shown in Fig. 2a. It clearly indicates the presence of some unexplained force at the smallest separations. However, later the author recognized that he was working with aspheric lens which had the curvature radius $`R=12.5\pm 0.3cm`$ in the place where the force was measured. The correction was published in erratum . The points for $`\mathrm{\Delta }F(a_i)`$ with the corrected $`R`$ are presented in Fig. 2b. This time the prediction obviously does not contradict to the experiment but dealing with the upper limit we cannot conclude that there is an agreement, either.
The question about surface distortions in TP experiment has been raised in . Surfaces of the bodies have not been examined in but roughness of the order of $`3040nm`$ is quite typical for a metallic film on a polished substrate and correlates with the substrate roughness. Quartz optics is used for near UV light and its surface has to be polished with a precision at least $`\lambda /10`$, where $`\lambda 300nm`$ is the UV wavelength. It supports the value above which is routinely observed with atomic force or tunnel microscope. According to (6) the short-scale stochastic distortions give only a few percent correction even for the smallest separation $`a=0.6\mu m`$. Large-scale deviations seem potentially more dangerous since the correction can be the first order in $`A/a`$ especially if we take into account that the lens was not spherical. The radius $`r_{int}`$ of the interaction area one can estimate as $`r_{int}\sqrt{Ra}1000\mu m`$. Therefore, only small area on the lens takes part in the interaction. In this place the lens can be represented as part of the parabolic surface
$$z=\frac{r^2}{2\left(R+\mathrm{\Delta }R\right)}\frac{r^2}{2R}\frac{1}{2}\frac{r^2}{R^2}\mathrm{\Delta }R,$$
(22)
where $`\mathrm{\Delta }R`$ is the error in the curvature radius. Here the second term describes the error in the plate-lens separation and it can be taken as the distortion amplitude. This amplitude is maximal for $`r=r_{int}`$ and for the relative amplitude one has an estimate
$$\frac{A}{a}\frac{1}{2}\frac{\mathrm{\Delta }R}{R}1.210^2.$$
(23)
This value is rather small and according to (6) the correction to the force will be less but comparable with the experimental errors. Moreover, negligible role of the large-scale distortions is actually an experimental fact. The region of the plate and sphere used for the force measurement in was varied by tilting the lens with the adjustment screws and there was no evidence for any variation of the force depending on the region used for the measurement.
The first attempt to evaluate the Casimir force was undertaken by the author of TP experiment who takes into account the first finite conductivity correction but used very small $`\omega _p`$ for $`Au`$. Lamoreaux was the first who recognized the necessity of more rigorous approach to the force evaluation and importance of thin films on the metallic surfaces . His numerical results were not quite good due to the delicate problem with choice of $`\omega _p`$ which we discussed above. The matter has been settled in with the result which coincide with ours. However, our statement is that the calculated force represents the upper limit but not the force itself. The reason is that evaporated $`Au`$ film will have smaller plasma frequency than $`1.3710^{16}s^1`$ due to large number of defects in the film. To know the exact value of the force, one has to measure the dielectric function of the bodies but not to take it from a handbook.
### B AFM experiment
Let us discuss now the upper limit on the Casimir force for the AFM experiment . The plasma frequency for $`Al`$ can be restricted using (20) if one supposes that every atom produces 3 free electrons. It gives $`\omega _p^{Al}=2.4010^{16}s^1`$ that coincide with the largest value in Table I. The resistivity of perfect crystals is $`\rho _0^{Al}=2.65\mu \mathrm{\Omega }cm`$. Since we successfully predicted the plasma frequencies for the best samples of $`Au`$ and $`Al`$, the same way one can use to estimate $`\omega _p`$ for $`Au/Pd`$. If each $`Au`$ atom gives one and $`Pd`$ atom gives not more than two free electrons, then $`\omega _p^{Au/Pd}=1.6910^{16}s^1`$. This alloy is used in microelectronics and resistivity of the bulk material is known to be $`\rho _0^{Au/Pd}30\mu \mathrm{\Omega }cm`$ in accordance with the statement that alloys have large Drude damping. These data allow to find the upper limit on the force using (10) with the functions $`G_{1,2}`$ defined in (13).
Before comparing the upper limit with the measured force we have to discuss a few additional aspects concerning the experiment. Real surface of the bodies is always distorted and the distortions are especially important to treat the data in . The distortion statistics were analysed with the atomic force microscope . The force has to be averaged over the distorted surfaces and we use for this the procedure developed in . The major distortions are the large separate crystals situated irregularly on the surfaces with a typical lateral size of $`200nm`$. The height of the highest and intermediate distortions is about $`h_1=40nm`$ and $`h_2=20nm`$, respectively. The homogeneous stochastic background of the averaged height $`h_0/2=5nm`$ fills the surface between the major distortions. The character of roughness on the plate and on the ball is quite similar. The part of the surface occupied by distortions with the height $`h_1`$, $`h_2`$, and $`h_0/2`$ was measured as $`v_1=0.11`$, $`v_2=0.25`$, and $`v_0=0.64`$, respectively. These values are the probabilities for the corresponding distortion to appear. The body surface is defined in such a way that averaging over distortions gives zero result. Then the averaged force is the sum of local forces for all possible kinds of distortions which face each other taken with the corresponding probabilities
$$F^{dist}\left(a\right)=\underset{i,j=0}{\overset{2}{}}v_iv_jF(a_{ij}),$$
(24)
where $`a_{ij}`$ are the local separations defined in . For us it will be important that the minimal local separation is $`a_{11}=a54.8nm`$. This procedure seems quite reliable but, of course, large distortions give the feeling of uncertainty. The upper limit on the Casimir force has to be averaged with the corresponding roughness parameters according to (24).
The raw force $`F_m`$ measured in the experiments consists of a few components
$$F_m=F_c(a_1+a_0)+F_e\left(a_1+a_0\right)+C\left(a_1+a_0\right).$$
(25)
Here $`a_1`$ is the separation from the voltage applied to the piezo corrected to the cantilever deflection, $`a_0`$ is the parameter chosen in such a way that $`a=a_1+a_0`$ is the absolute separation between bodies, the first term in (25) is the Casimir force, the second term is the electrostatic force corresponding to the measured contact potential $`29mV`$, the third term represents the linearly increasing coupling of the scattered light into the photodiodes (see for details). The parameters $`a_0`$ and $`C`$ were determined at large separations, where $`F_c`$ is represented by the bare force (3). Then the Casimir force can be extracted from the raw data with the help of (25). Similar way to find the Casimir force was used in . Of course, the separation $`a`$ has to be defined as the distance between averaged surfaces of $`Au/Pd`$ layers. However, the role of these layers have been underestimated in . In the Casimir force was found from (25) but $`a`$ was interpreted as the absolute separation between $`Al`$ surfaces . Effectively the $`Au/Pd`$ layers were changed by $`Al`$ which has larger $`\omega _p`$ and, therefore, the force calculated theoretically was overestimated.
It was indicated that the thickness of $`Au/Pd`$ layer is less than $`20nm`$, that is why for calculations we use the conservative value $`h=15nm`$ . This change makes the force only larger. The experimental points from (triangles) and theoretical upper limit on the force including the roughness correction (solid line) are shown in Fig. 3 in the small separations range $`a<250nm`$. If the top layer is changed by $`Al`$, it enlarges the force on $`15pN`$ at the smallest separation $`a=120nm`$. It is clear that the top layer definitely cannot be ignored in the force evaluation. Variation of $`\omega _p^{Al}`$ on 10% gives only $`1pN`$ change in the force because of screening effect of the top layer. The same variation in $`\omega _p^{Au/Pd}`$ changes the force on $`2pN`$. The resistivity variation of the $`Au/Pd`$ layer on 30% gives $`1pN`$ effect. At larger separations all the effects become smaller.
We can see from Fig. 3 that the upper limit is smaller than the force measured at small separations and the difference is significant. This conclusion contradicts to that in , where good agreement between theory and experiment has been reported (dotted line) based on the detailed theoretical analysis. We have already stressed the importance of the $`Au/Pd`$ layer but it is not the only reason of deviations. It comes also from poor behavior of the finite conductivity correction used in at small separations $`a`$. The correction was based on a simple interpolating formula for the force between two plates
$$F_c^{plates}(a)=F_c^{0plates}(a)\left(1+\frac{11}{3}\frac{c}{a\omega _p}\right)^{\frac{16}{11}}$$
(26)
which is applicable in wider range of separations than the expansion up to the second order (4). It was used to calculate the conductivity correction between sphere and plate. Applying the proximity force theorem to (26) one gets
$$F(a)=3F_c^0(a)\underset{1}{\overset{\mathrm{}}{}}\frac{dx}{x^4}\left(1+\frac{11}{3x}\frac{c}{a\omega _p}\right)^{\frac{16}{11}},$$
(27)
where the upper limit was moved to infinity since $`aR`$. This integral can be expressed via the Gauss hypergeometric function but in it was expanded in the series up to the fourth order
$$F(a)=F_c^0(a)\left[14\frac{c}{a\omega _p}+\frac{72}{5}\left(\frac{c}{a\omega _p}\right)^2\frac{152}{3}\left(\frac{c}{a\omega _p}\right)^3+\frac{532}{3}\left(\frac{c}{a\omega _p}\right)^4\right]$$
(28)
and the equations (24) and (28) were used to get the theoretical prediction for the corrected Casimir force (dotted line in Fig. 3). We found that the interpolating curve which we got by numerical intergation of (27) is very close to the exact force evaluated according to (10), (8) with the parameters $`\omega _p=1.8810^{16}s^1`$, $`\omega _\tau =0`$ from and $`T=300^{}K`$ in all range of distances. The small difference between the curves is the temperature effect which disappears when $`T0`$. However, the expansion (28) works very bad at $`a<100nm`$. Although the separation in the experiment exceeds $`100nm`$, the local distance in (24) can be as small as $`65nm`$, where (28) is absolutely unacceptable. The same is true if one uses the expansion up to the fourth order found in directly form (10). The dashed line in Fig. 3 shows the force calculated according to (24), (27) with the parameters above. The divergence of the solid and dashed curves is the effect of the top layer and different parameters used for $`Al`$. The larger $`\omega _p^{Al}`$ which we are using for the upper limit is partly compensated by the top layer and that is why the dashed curve lies not too far from the solid one.
### C Improved AFM experiment
Very important progress has been achieved in where controlled metal evaporation and smaller thickness of $`Au/Pd`$ layer ($`h=8nm`$ instead of $`20nm`$) allowed to reduce the surface roughness to the level when the correction to the force becomes practically unimportant. Also the contact potential has been considerably reduced and the parameter $`a_0`$ defining the absolute separation of the surfaces has been found in independent electrostatic measurements. Interaction between metallized ball and corrugated plate has been probed in . It is not a subject of our consideration here. The data for flat plate and sphere in this work are actually in very good agreement with that given in and we will not discuss them specially. The roughness parameters have been reduced to the following : $`h_1=14nm,v_1=0.05;h_2=7nm,v_2=0.11;h_0/2=2nm,v_0=0.84`$. Unfortunately, the unjustified assumption that the force is insensitive to the presence of $`Au/Pd`$ layer has been inserted in the procedure of the Casimir force extraction from the raw data and the following relation has been used instead of (25)
$$F_m=F_c(a+2h)+F_e\left(a\right)+Ca.$$
(29)
For this reason we cannot directly use the data in to compare with the theoretical prediction. It becomes obvious if we plot in the same figure the measured force in the experiments and (see Fig. 4). One would expect that for thicker $`Au/Pd`$ layer the force has to be smaller, but the actual relation is opposite and the difference is large. Fortunately, it is easy to restore the right data. In $`a_0`$ was found from an independent electrostatic measurement and the constant $`C`$ was determined at large separations when the shift on $`2h=16nm`$ in $`F_c`$ argument was practically unimportant. The measured Casimir force $`F_{cm}`$ was expressed as
$$F_{cm}(a+2h)=F_m(a)F_e(a)Ca,$$
(30)
but the points in Fig. 4 taken from represent the force as a function of true separation. Therefore, the force presented in the figure was calculated as
$$F_{cm}(a)=F_m(a2h)F_e(a2h)C(a2h).$$
(31)
The $`Au/Pd`$ layer certainly cannot be ignored and the right expression for the measured Casimir force must be
$$F_{cm}(a)=F_m(a)F_e(a)Ca.$$
(32)
It is obvious that to restore the right data one has to shift the open squares in Fig. 4 on $`2h=16nm`$ to larger separations. After this shift a good agreement between two different experiments is reached. To be absolutely sure that the right transformation was done we have tried to reproduce the measured force directly from the raw data presented in and , where the procedure has been described in details. The data were available only for one scan and for this reason our calculations had restricted precision, but it was enough to make a conclusion on reproducibility. To check the procedure, we successfully reproduced the force from the raw data in . The force found from the raw data in according to (32) agrees much better with the shifted points than with the ones presented in . These detailed explanations are given not only to answer the criticism of our preprint but also because of great importance of the conclusion. It is stated in that the points have to be shifted to smaller separations. Even ignoring the arguments above it is obvious from Fig. 4 that such a shift would give drastic disagreement between two experiments made by the same method.
The upper limit on the Casimir force in conditions of the experiment one can find exactly as was explained above. The only difference is the other set of roughness parameters, but in this case the roughness correction is on the level of experimental errors. The experimental points from shifted on $`16nm`$ as was discussed above and the corresponding upper limit (solid line) are presented in Fig. 5. Again we can see that the upper limit is smaller than the measured force and deviation increases at smaller separations. Moreover, even if we replace the $`8nm`$ thick top layer by $`Al`$ with the maximal plasma frequency $`\omega _p=2.4010^{16}s^1`$, this disagreement will not disappear as shows the dashed line. The residual force defined according (21) in the experiments (triangles) and (open squares) is shown in Fig. 6. It clearly demonstrates the presence of some unexplained attractive force which is decreasing rapidly when the distance between bodies increased. The points from two different experiments are in reasonable agreement with each other that means that the residual force is reproducible. The residual force becomes larger if we deviate the parameters from their limit values but the agreement between two experiments is not broken.
## V Discussion
Let us discuss now possible reasons for disagreement between experiment and theoretical expectations. As was mentioned above the main problem is the values of the material parameters which can significantly deviate from their handbook values for evaporated or spattered metallic films. The idea of this paper was to find the upper limit on the force instead of the force itself. It allowed to use only well defined parameters of perfect single-crystal materials. We took the largest values for the plasma frequencies and the smallest ones for the resistivities. Any possible deviation from these values will make the force only smaller and disagreement between theory and experiment will be larger.
Some doubts were raised about the possibility to describe the thin top layer by a dielectric function which depends only on frequency. It was stated that the spatial dispersion can be important for thin films because the distance traveled by electron during one period of the field can be larger than the film thickness
$$\frac{v_F}{\omega }>h,$$
(33)
where $`v_F`$ is the velocity of the electron on the Fermi surface. This dimensional effect is really exist (see, for example, ) but it is difficult for observation at room temperature. The reason is that for thin metallic films the mean free path for electrons is very short ($`<100\AA `$) because of large concentration of the defects. Typically the resistivity of very thing films is on the level of $`100\mu \mathrm{\Omega }cm`$. Then for $`\omega _p10^{16}s^1`$ from (17) one finds $`\omega _\tau 10^{15}s^1`$. The mean free path is estimated as
$$l=\frac{v_F}{\omega _\tau }10\AA $$
(34)
that is smaller than the used $`Au/Pd`$ film thickness and, therefore, the spatial dispersion can be neglected. That is why the standard dependence for the dielectric function $`\epsilon \left(\omega \right)`$ is widely used in optics of metals up to the film thickness in a few nanometers when quantum effects become involved. For the same reason one can neglect the anomalous skin effect for evaporated (spattered) films even for thick ones. Extremely high resistivity $`2000\mu \mathrm{\Omega }cm`$ for $`60nm`$ thick $`Au/Pd`$ film has been reported in . However, the authors themselves stress that the resistance of the film is, to all appearance, dominated by grain boundaries but optical properties of the film are quite usual. This example shows once more that the details of the spattering technology cannot be ignored.
In uncertainty was expressed about applicability of the proximity force theorem. At the moment there is no any work where the force between sphere and plate was calculated from ”the first principles”. There are some heuristic approaches allowing to calculate nonadditive Casimir force which agree well with the result found by using PFT (see additional discussion and references in ). The PFT states that the main contribution to the force can be found by adding the contributions of various distances as if they were independent and it is applicable to nonadditive forces. An example is the electrostatic force which is nonadditive because the surface charge density is nonuniform for curved surfaces. One can check that the PFT gives in this case the correct result with the accuracy $`a/R`$. For the discussed experiments the correction is very small ($`a/R0.001`$). Even if this term appears with a large coefficient, say $`10`$, the correction will be only on the level of the experimental errors.
In the AFM experiments the electrostatic attraction between bodies because of contact potential was carefully taken into account. Of course, the aluminum surfaces were partly oxidized and electrons could be trapped in the oxide. These charges can be potentially dangerous if the $`Au/Pd`$ film is not continuous. In this case the trapped charges and their images in underlying aluminum will be the source of the dipole field. Then an additional force can arise as a result of dipole-dipole interaction. However, it is difficult to make a reliable estimate for this effect because we do not know the concentration of trapped charges and the size of islands in $`Au/Pd`$ layer or even do discontinuities exist at all for the used layer thickness (it depends on details of the covering procedure). In this connection to make the experiment absolutely clear, it is preferable to use $`Au`$ instead of $`Al`$ metallization because its non-reactive surface has strong advantage over $`Al`$. It excludes also additional uncertainties connected with $`Au/Pd`$ layer. One can use as well silver or copper but they are not as inert as gold. It is difficult to measure the dielectric function at the wavelengths larger than $`30\mu m`$ but this range gives an important contribution to the dispersion relation. That is why the material behavior in this range has to be predictable. One can say definitely that the materials of platinum group cannot be used since they are not described by the Drude dielectric function at low frequencies.
One can speculate that the observed discrepancy is explained by a new Yukawa force mediated by a light scalar boson. Then interaction of two atoms is described by the Yukawa potential
$$V_Y\left(r\right)=\alpha N_1N_2\frac{\mathrm{}c}{r}\mathrm{exp}(r/\lambda ),$$
(35)
where $`\alpha `$ is a dimensionless interaction constant, $`\lambda `$ is the Compton wavelength of a particle responsible for the interaction, and $`N_{1,2}`$ is the number of nucleons in atoms of the interacting bodies. An additional advantage of $`Au`$ metallization is higher density of the bodies coating. In this case the Yukawa force will be enlarged roughly by the factor $`\left(\rho _{Au}/\rho _{Al}\right)^250`$, where $`\rho _{Au,Al}`$ are the material densities. If the observed discrepancy has relation with the Yukawa interaction, the AFM experiment with $`Au`$ metallization of the bodies will definitely reveal this new force even without detailed knowledge of optical properties of the metallization.
## VI Conclusion
We have analysed the results of recent precise measurements of the Casimir force using the Lifshitz theory to evaluate the force. Layered structure of the bodies coating was taken into account in the frame of the Lifshitz approach. It was stressed that the force cannot be predicted with necessary accuracy if there is no detailed information on the dielectric function of the bodies coating. Fortunately, all the used materials ($`Au`$, $`Al`$ and $`Au_{0.6}Pd_{0.4}`$) are well described in terms of Drude parameters $`\omega _p`$ and $`\omega _\tau `$ in the infrared range which dominates in the dispersion relation for the dielectric function $`\epsilon \left(i\zeta \right)`$. It was noted that one can find the upper limit on the Casimir force that realized for perfect single-crystal materials for which electrical and optical properties are well defined. The surface roughness and linear in temperature corrections were taken into consideration. It was shown that the upper limit on the Casimir force does not contradict to the result of the torsion pendulum experiment . The main conclusion of the paper is that the upper limit is smaller than the observed force in the AFM experiments and the difference far exceeds experimental errors and theoretical uncertainties for small separations between bodies. The simplest modification of the experiment is proposed allowing to reveal origin of the discrepancy.
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# Perturbative renormalization of the first two moments of non-singlet quark distributions with overlap fermions
## 1 Introduction
In recent years remarkable progress has been made towards formulations of lattice fermions which have no doublers and possess an exact chiral symmetry without giving up desirable features like flavor symmetry, locality, unitarity or gauge-invariance. The Ginsparg-Wilson relation
$$\gamma _5D+D\gamma _5=a\frac{1}{\rho }D\gamma _5D$$
(1)
has been recognized as the fundamental condition in this context, and a Dirac operator $`D`$ which satisfies this relation can indeed describe chiral fermions on the lattice. One of the possible solutions of the Ginsparg-Wilson relation has been found by Neuberger starting from the overlap formalism . In the massless case the overlap-Dirac operator is
$$D_N=\frac{1}{a}\rho \left[1+\frac{X}{\sqrt{X^{}X}}\right],$$
(2)
where
$$X=D_W\frac{1}{a}\rho $$
(3)
in terms of the usual Wilson-Dirac operator
$$D_W=\frac{1}{2}\left[\gamma _\mu (_\mu ^{}+_\mu )ar_\mu ^{}_\mu \right],$$
(4)
$$_\mu \psi (x)=\frac{1}{a}\left[U(x,\mu )\psi (x+a\widehat{\mu })\psi (x)\right].$$
(5)
In the range $`0<\rho <2r`$ the right spectrum of massless fermions is obtained . For a quark of bare mass $`m_0`$ the overlap-Dirac operator becomes
$$\left(1\frac{1}{2\rho }am_0\right)D_N+m_0.$$
(6)
Since additive mass renormalization is forbidden, one avoids altogether a source of systematic errors that is always present with Wilson fermions .
Lüscher has shown that a fermion obeying the Ginsparg-Wilson relation possesses an exact chiral symmetry of the general form
$$\delta \psi =ϵ\gamma _5\left(1\frac{c}{\rho }aD\right)\psi ,\delta \overline{\psi }=ϵ\overline{\psi }\left(1\frac{1c}{\rho }aD\right)\gamma _5.$$
(7)
A fully legitimate form of global chiral symmetry is thus preserved for finite lattice spacing. The expected value of the global anomaly (which comes from the non-invariance of the fermionic integration measure under the transformations above), the Atiyah-Singer index theorem on the lattice and the chiral Ward identities are then fully attained before taking the continuum limit. Thus, the central point in this framework is not to insist on the lattice form of the canonical chiral transformations (in fact actions satisfying Eq. (1) are not chirally invariant in the canonical sense), but rather be satisfied with a modified form of the chiral symmetry. The Nielsen-Ninomyia theorem is then circumvented altogether, and it is possible to define chiral fermions without doublers or breaking of flavor symmetry or other unpleasant drawbacks. The chiral transformations (7) depend also on the interaction, but this does not forbid the non-perturbative construction of chiral gauge theories on the lattice with exact gauge invariance . For recent reviews on the interesting properties of fermions satisfying the Ginsparg-Wilson relation, see Refs. .
We will be interested in the following in performing chiral-invariant computations by using the overlap-Dirac operator (2). This has been proven to be local <sup>1</sup><sup>1</sup>1The locality is not meant here to be strict locality, but in the larger sense that the strength of the interaction decays exponentially with the distance. , and although simulations with the Neuberger operator look computationally very demanding, progress is under way . There has been activity also on the analytic side, and some 1-loop calculations with overlap fermions have been already carried out . The most recent and advanced calculations have featured the relation between the $`\mathrm{\Lambda }`$ parameter in the lattice scheme defined by the overlap operator and in the $`\overline{\mathrm{MS}}`$ scheme , and the renormalization factors of the quark bilinears $`\overline{\psi }\mathrm{\Gamma }\psi `$ .
In this paper we present the calculation in lattice QCD of the renormalization factors of a few operators which measure the lowest two moments of non-singlet quark distributions. The generic operators which measure their $`n`$-th moments are $`\overline{\psi }\gamma _\mu D_{\mu _1}\mathrm{}D_{\mu _n}\psi `$ and $`\overline{\psi }\gamma _\mu \gamma _5D_{\mu _1}\mathrm{}D_{\mu _n}\psi `$ for unpolarized and polarized Structure Functions respectively. Contrary to what happens with Wilson fermions, in the present case, thanks to the exact chiral symmetry that we maintain, the renormalization constants for every pair of corresponding unpolarized and polarized operators which differ by a $`\gamma _5`$ matrix are the same.
This paper is organized as follows: in Sect. 2 we introduce the various operators that we have studied, in Sect. 3 their renormalization on the lattice is discussed, and in Sect. 4 some details about the perturbative calculations are given. Finally, in Sect. 5 we present our results. In the Appendices one can find the Feynman rules that we have used as well as some analytic results.
## 2 Moments of Structure Functions
The hadronic physics in Deep Inelastic Scattering is contained in the matrix elements of the T-product of the hadronic currents
$$id^4xe^{iqx}p|T(J_\mu (x)J_\nu (0))|p,$$
(8)
$`q`$ being the momentum transfer and $`p`$ the target momentum. An operator product expansion on the light cone of the kind
$$J(x)J(0)\underset{n,i,l}{}C_l^{n,i}(x^2)x^{\mu _1}\mathrm{}x^{\mu _n}O_{\mu _1\mathrm{}\mu _n}^{(n,i)}(0)$$
(9)
describes the physics in the Bjorken limit, in which scaling is reached and the Structure Functions depend only on the Bjorken variable $`x_B=q^2/(2pq)`$. The moments of the Structure Functions, which measure the moments of quark distributions, are directly connected to the forward matrix elements of the local operators appearing in the light-cone expansion. Since the divergence of the Wilson coefficients is governed by the twist (dimension minus spin) of the corresponding operators, infinite towers of operators of increasing dimension and spin but with the same twist (dimension minus spin) appear at a given order in $`1/q^2`$. The dominant contribution is given by twist two, which in the flavor non-singlet case means the symmetric traceless operators
$`O_{\mu \mu _1\mathrm{}\mu _n}`$ $`=`$ $`\overline{\psi }\gamma _{\{\mu }D_{\mu _1}\mathrm{}D_{\mu _n\}}{\displaystyle \frac{\lambda ^a}{2}}\psi `$ (10)
$`O_{\mu \mu _1\mathrm{}\mu _n}^{(5)}`$ $`=`$ $`\overline{\psi }\gamma _{\{\mu }\gamma _5D_{\mu _1}\mathrm{}D_{\mu _n\}}{\displaystyle \frac{\lambda ^a}{2}}\psi ,`$ (11)
where $`\lambda ^a`$ are flavor matrices (which in the following will be omitted and implicitly understood). These operators measure the moment of unpolarized and polarized Structure Functions respectively, that is the moments $`x_B^n`$ of quark momentum distributions and the moments $`(\mathrm{\Delta }x_B)^n`$ of quark helicity distributions, at leading twist. Higher-twist operators give the sub-dominant contributions, and in particular twist-4 operators (which include also 4-fermion operators) measure the $`1/q^2`$ corrections to these moments .
Since we are considering flavor non-singlet operators, there can be no mixing with operators like $`\mathrm{Tr}_\rho F_{\mu _1\rho }D_{\mu _2}\mathrm{}F_{\rho \mu _n}`$ and $`\mathrm{Tr}_\rho \stackrel{~}{F}_{\mu _1\rho }D_{\mu _2}\mathrm{}F_{\rho \mu _n}`$ which measure the gluon distributions. Mixing with these operators is prohibited even when one considers unquenched calculations. Singlet quark operators instead do mix with gluon operators. For more detailed discussions of Deep Inelastic Scattering on the lattice, see Refs. .
The operators in the light-cone expansion (9) of which we compute the renormalization factors are the following:
$`O_{v_2,d}`$ $`=`$ $`\overline{\psi }\gamma _{\{1}D_{4\}}\psi `$ (12)
$`O_{v_2,e}`$ $`=`$ $`\overline{\psi }\gamma _4D_4\psi {\displaystyle \frac{1}{3}}{\displaystyle \underset{i=1}{\overset{3}{}}}\overline{\psi }\gamma _iD_i\psi ,`$ (13)
which measure the first moment of quark momentum distributions,
$`O_{a_2,d}`$ $`=`$ $`\overline{\psi }\gamma _{\{1}\gamma _5D_{4\}}\psi `$ (14)
$`O_{a_2,e}`$ $`=`$ $`\overline{\psi }\gamma _4\gamma _5D_4\psi {\displaystyle \frac{1}{3}}{\displaystyle \underset{i=1}{\overset{3}{}}}\overline{\psi }\gamma _i\gamma _5D_i\psi ,`$ (15)
which measure the first moment of quark helicity distributions,
$`O_{v_3,d}`$ $`=`$ $`\overline{\psi }\gamma _{\{4}D_1D_{2\}}\psi `$ (16)
$`O_{v_3,e}`$ $`=`$ $`\overline{\psi }\gamma _{\{4}D_1D_{1\}}\psi {\displaystyle \frac{1}{2}}{\displaystyle \underset{i=2}{\overset{3}{}}}\overline{\psi }\gamma _{\{4}D_iD_{i\}}\psi ,`$ (17)
which measure the second moment of quark momentum distributions, and
$`O_{a_3,d}`$ $`=`$ $`\overline{\psi }\gamma _{\{4}\gamma _5D_1D_{2\}}\psi `$ (18)
$`O_{a_3,e}`$ $`=`$ $`\overline{\psi }\gamma _{\{4}\gamma _5D_1D_{1\}}\psi {\displaystyle \frac{1}{2}}{\displaystyle \underset{i=2}{\overset{3}{}}}\overline{\psi }\gamma _{\{4}\gamma _5D_iD_{i\}}\psi ,`$ (19)
which measure the second moment of quark helicity distributions. Symmetrization in all Lorentz indices is understood. We use $`\stackrel{}{D}=\stackrel{}{D}\stackrel{}{D}`$, with the following lattice discretizations:
$`\stackrel{}{D_\mu }\psi (x)`$ $`=`$ $`{\displaystyle \frac{1}{2a}}\left[U(x,\mu )\psi (x+a\widehat{\mu })U^{}(xa\widehat{\mu },\mu )\psi (xa\widehat{\mu })\right]`$ (20)
$`\overline{\psi }(x)\stackrel{}{D_\mu }`$ $`=`$ $`{\displaystyle \frac{1}{2a}}\left[\overline{\psi }(x+a\widehat{\mu })U^{}(x,\mu )\overline{\psi }(xa\widehat{\mu })U(xa\widehat{\mu },\mu )\right].`$ (21)
We have chosen the Lorentz indices of each operator appearing in the continuum expansion in two different ways, so that they fall in two different representations of the hypercubic group (the symmetry group of the lattice, remnant of the Lorentz symmetry), and on the lattice they will then renormalize in a different way. The representations where the indices are all different from each other are least likely to mix with other operators, however one needs more components of the hadron momentum different from zero when simulating the corresponding matrix elements, and this can lead to more lattice artifacts. Sometimes one of the representations can be easier to handle in practice, depending on the trade-off between the amount of mixing and the number of non-zero momentum components. For example, in the case of the second moment the operator (16) is multiplicatively renormalized, but it can be more advantageous to use the other representation (operator (17)) when it is important that fewer components of the momenta are different from zero, although in this case one has then to deal with a mixing.
The operator (17) (together with the corresponding polarized (19)) is the only one in the list above that is not multiplicatively renormalized on the lattice. What happens is that that the two operators
$`O_A`$ $`=`$ $`\overline{\psi }\gamma _4D_1D_1\psi {\displaystyle \frac{1}{2}}{\displaystyle \underset{i=2}{\overset{3}{}}}\overline{\psi }\gamma _4D_iD_i\psi ,`$ (22)
$`O_B`$ $`=`$ $`\overline{\psi }\gamma _1D_4D_1\psi +\overline{\psi }\gamma _1D_1D_4\psi {\displaystyle \frac{1}{2}}{\displaystyle \underset{i=2}{\overset{3}{}}}\overline{\psi }\gamma _iD_4D_i\psi {\displaystyle \frac{1}{2}}{\displaystyle \underset{i=2}{\overset{3}{}}}\overline{\psi }\gamma _iD_iD_4\psi ,`$ (23)
which are not symmetrized in their indices, do not receive the same 1-loop corrections on the lattice , and so the operator (17),
$$O_{v_3,e}=\frac{1}{3}\left(O_A+O_B\right),$$
(24)
does not go into itself under lattice renormalization, because the symmetric combination is lost. In the polarized case, an analogous situation occurs for
$`O_{A^5}`$ $`=`$ $`\overline{\psi }\gamma _4\gamma _5D_1D_1\psi {\displaystyle \frac{1}{2}}{\displaystyle \underset{i=2}{\overset{3}{}}}\overline{\psi }\gamma _4\gamma _5D_iD_i\psi ,`$ (25)
$`O_{B^5}`$ $`=`$ $`\overline{\psi }\gamma _1\gamma _5D_4D_1\psi +\overline{\psi }\gamma _1\gamma _5D_1D_4\psi `$
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=2}{\overset{3}{}}}\overline{\psi }\gamma _i\gamma _5D_4D_i\psi {\displaystyle \frac{1}{2}}{\displaystyle \underset{i=2}{\overset{3}{}}}\overline{\psi }\gamma _i\gamma _5D_iD_4\psi .`$
A detailed discussion of these effects can be found in Sect. 5.
We mention that for the second moment there exists also a third independent lattice representation, in which all indices are equal. However, it mixes with a lower dimensional operator with a power-divergent coefficient ($`1/a^2`$), and so we do not consider it here.
## 3 Renormalization
To relate the numbers obtained from Monte Carlo simulations to physical continuum quantities, a lattice renormalization of the relevant matrix elements is necessary. The connection is given by
$$O_i^{\mathrm{cont}}=\underset{j}{}\left(\delta _{ij}\frac{g_0^2}{16\pi ^2}\left(R_{ij}^{\mathrm{lat}}R_{ij}^{\mathrm{cont}}\right)\right)O_j^{\mathrm{lat}},$$
(27)
where
$$O_i^{\mathrm{cont},\mathrm{lat}}=\underset{j}{}\left(\delta _{ij}+\frac{g_0^2}{16\pi ^2}R_{ij}^{\mathrm{cont},\mathrm{lat}}\right)O_j^{\mathrm{tree}}$$
(28)
are the continuum and lattice 1–loop expressions respectively, and the tree-level matrix element is the same in both cases. The differences $`\mathrm{\Delta }R_{ij}=R_{ij}^{\mathrm{lat}}R_{ij}^{\mathrm{cont}}`$ enter then in the renormalization factors
$$Z_{ij}(a\mu ,g_0)=\delta _{ij}\frac{g_0^2}{16\pi ^2}\mathrm{\Delta }R_{ij}(a\mu )$$
(29)
which connect the lattice to the continuum. These renormalization factors are independent of the state, depend only on the scale $`a\mu `$ and are gauge-invariant. Using them, a matrix element obtained from Monte Carlo simulations can be renormalized to a continuum scheme.
Although there are in principle also non-perturbative methods with which one can determine the renormalization factors, perturbation theory is still important. In fact, it can happen that for some operators a window for the non-perturbative signal is difficult to obtain, and then perturbation theory remains the only possibility of computing the relevant renormalization factors. For the Neuberger operator, extracting the non-perturbative renormalization factors in addition to simulating the bare matrix elements could turn out to be computationally very expensive. In general, perturbative lattice renormalization is important as a hint and a guide for non–perturbative renormalization studies, and even more when mixings are present, which are generally more intricate than in the continuum case, and more transparent when looked at in perturbation theory, especially if some amounts of mixings are small. Perturbative renormalization can in any case be very useful in checking and understanding results obtained with non-perturbative methods.
In the continuum we renormalize the operators in the $`\overline{\mathrm{MS}}`$ scheme. As perturbative renormalization condition on the lattice we impose that the 1-loop amputated matrix elements at a certain reference scale $`\mu `$ are equal to the corresponding bare tree–level quantities. For lattice matrix elements of multiplicatively renormalized operators computed between one-quark states this condition means
$`p|O^{\mathrm{lat}}(\mu )|p|_{p^2=\mu ^2}`$ $`=`$ $`Z_O(a\mu ,g_0)Z_\psi ^1(a\mu ,g_0)p|O^{(0)}(a)|p|_{p^2=\mu ^2}^{1\mathrm{loop}}`$ (30)
$`=`$ $`p|O^{(0)}(a)|p|_{p^2=\mu ^2}^{\mathrm{tree}},`$
where $`Z_\psi `$ is the wave-function renormalization, computed from the quark self-energy.
The 1-loop lattice matrix elements of the operators we consider here, which are at most logarithmically divergent, will have the form
$`p|O^{(0)}(a)|p|^{1\mathrm{loop}}`$ $`=`$ $`p|O^{(0)}(a)|p|^{\mathrm{tree}}\times `$
$`\left(1+{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left(\gamma _O\mathrm{log}a^2p^2+V_O+T_O+2S\right)\right),`$
where $`V_O`$ is the finite contribution of the vertex and sails diagrams (a, b and c in Fig. 1), $`T_O`$ refers to the tadpole arising from the operator (d in Fig. 1), S is the finite contribution (proportional to $`ip/`$) of the quark self-energy of one leg, including the leg tadpole (e and g, or f and h, in Fig. 1), and $`C_F=\frac{N_c^21}{2N_c}`$. The $`Z_O`$ factor for an operator $`O`$ will then be given by
$$Z_O(a\mu ,g_0)=1\frac{g_0^2}{16\pi ^2}C_F\left(\gamma _O\mathrm{log}a^2\mu ^2+B_O\right),$$
(32)
with
$$B_O=V_O+T_O+S.$$
(33)
We will call “proper” contributions the ones that exclude the self-energy diagrams. They correspond to the diagrams a-d in Fig. 1.
Although the theory defined by the overlap-Dirac operator (2) has an exact chiral symmetry and no lattice artifacts of order $`a`$ are present in the action, matrix elements of operators still possess corrections of order $`a`$ and therefore they need to be improved. The improved operator corresponding to $`\overline{\psi }\gamma _{\{\mu }D_{\mu _1}\mathrm{}D_{\mu _n\}}\psi `$ is
$$\overline{\psi }\left(1\frac{1}{2}aD_N\right)\gamma _{\{\mu }D_{\mu _1}\mathrm{}D_{\mu _n\}}\left(1\frac{1}{2}aD_N\right)\psi ,$$
(34)
and the big bonus with overlap fermions is that the renormalization factors for the improved operator and for the unimproved operator without the rotations $`(1\frac{1}{2}aD_N)`$ are the same. What happens is that in 1-loop amplitudes a factor $`D_N`$ can combine with a quark propagator, but since it has an $`a`$ in front and (contrary to the Wilson case) there is no $`1/a`$ factor in the propagator, as additive mass renormalization is forbidden by chiral symmetry, the contribution of $`D_N`$ to the renormalization factors is zero . Thus, we can simulate the improved operator (34) and renormalize it as if it were the unimproved operator, which saves a lot in the perturbative calculations.
Chiral symmetry, in addition to avoiding the extrapolations to zero quark mass in the simulations, gives also useful constraints on the mixing patterns and on the values of the renormalization factors. In general, mixings that in the Wilson case were allowed by the breaking of chirality are now forbidden, although this has less relevance here than in other problems like in weak interactions . An important consequence here is that, similarly to what happens in the case of the bilinears where we have $`Z_S=Z_P`$ and $`Z_V=Z_A`$ , chiral symmetry forces pairs of corresponding unpolarized and polarized operators, like $`O_{v_2,d}`$ and $`O_{a_2,d}`$, to have the same renormalization constants <sup>2</sup><sup>2</sup>2This is strictly true for improved operators like in Eq. (34), which transform like chiral multiplets, but since they have the same renormalization constants as the unimproved operators, these symmetry relations are preserved..
## 4 The computations
The interaction vertices and the propagator of the overlap-Dirac operator are much more complicated than the ones in the Wilson formulation, and this causes the computations to be rather cumbersome (see Appendices). Computer programs need then to be introduced, and we have performed the calculations using an ensemble of routines written in the symbolic manipulation language FORM. These routines are an extension of the ones used to do calculations with the Wilson action in various occasions .
The outputs of the FORM codes, which correspond to the results of the analytic calculations, are fed to Fortran programs which perform the numerical integrations. To treat the divergent integrals we use the Kawai method , in which they are split in integrals independent of external momenta plus continuum integrals. Divergent lattice integrals which contain both gluon and overlap-quark propagators are reduced to divergent integrals containing only gluon propagators, which are computed in a exact way as in , plus a finite rest.
Many checks have been performed on our results. We did the computations in a general covariant gauge <sup>3</sup><sup>3</sup>3The gluon propagator that we use is
$$G_{\mu \nu }(k)=\frac{1}{4_\rho \mathrm{sin}^2\frac{k_\rho }{2}}\left(\delta _{\mu \nu }(1\alpha )\frac{4\mathrm{sin}\frac{k_\mu }{2}\mathrm{sin}\frac{k_\nu }{2}}{4_\lambda \mathrm{sin}^2\frac{k_\lambda }{2}}\right).$$
(35) , and the cancellations of the gauge-dependent parts between the continuum scheme and the lattice results is a strong check on the good behavior of the FORM codes, as well as of the integration routines.
We have used the codes to compute also the renormalization constants for all quark bilinears, and we agree with the results given in Feynman-gauge by Alexandrou et al. , as well as with their self-energy. In the case of the leg tadpole and of the vertex graph of the scalar current $`\overline{\psi }\psi `$ we have also made the computations entirely by hand (although only in the Feynman gauge), and successfully checked them against the FORM output expressions. These analytic results are reported in Appendix B. We tried to do a calculation by hand for a first moment operator and for the rest of the leg self-energy, but they involved a huge amount of manipulations. We rely also on the vast amount of results which have been produced in different occasions using these FORM codes limited to the Wilson formulation , and these results have by now a certain number of cross-checks. With overlap fermions we need to introduce a new propagator and new vertices, but the operators are unaltered, and the same Wilson expressions for their expansion in $`a`$ can be used here. Many of the routines (for example the gamma algebra reduction) are also exactly the same as in the Wilson case.
We checked in all cases that the polarized operators have the same $`Z`$s as the corresponding unpolarized operators. This is a rather strong check, as the analytic results are widely different for the two cases, and it is only after performing the numerical integrations and seeing that the difference of the two results is much smaller than the precision of our integrals that we able to say that the renormalization factors are indeed equal. Finally, since our codes are able to do calculations both in Dimensional Regularization and using a mass regularization, we also checked that using two different intermediate regularizations while treating the divergences with the Kawai method we get the same results.
We obtain the $`B_O`$ constants with 5 digits after the decimal points, by doing the numerical integrations on a $`40^4`$ regular grid and using the method proposed by Lüscher and Weisz in to accelerate convergence. We have checked that these digits do not vary when a $`60^4`$ integration grid is used. In any case, to get more significant digits, grids of $`80^4`$ or $`100^4`$ points would still not be enough, and they would require an enormous computational effort, as a typical FORM output for the second moment operators already contains thousands of terms. Comparing our results for the self-energy and the bilinears with the numbers in Ref. , we can see that a few times a difference of one unit on the fifth digit can be noticed, and this agrees with our estimate of errors. In that paper the results are given in terms of the constants $`b=\frac{B}{16\pi ^2}`$, and this explains the presence of two more significant digits after the decimal point.
## 5 Results
We give in this Section the renormalization factors of the operators considered in this work. That these factors are identical for pairs of corresponding unpolarized and polarized operators has been explicitly verified for all operators. We give then the numerical results only for the unpolarized operators, for $`r=1`$.
We first consider the 1-loop contributions of the proper diagrams (a-d in Fig. 1), which for the operators that do not mix are
$`O_{v_2,d}^{\mathrm{proper}}={\displaystyle \frac{g_0^2}{16\pi ^2}}C_F`$ $`[({\displaystyle \frac{5}{3}}+(1\alpha ))\mathrm{log}a^2p^2`$
$`+V_{v_2,d}^{\alpha =1}(1\alpha )\mathrm{\hspace{0.17em}6.850272}+T_{v_2,d}]O_{v_2,d}^{\mathrm{tree}},`$
$`O_{v_2,e}^{\mathrm{proper}}={\displaystyle \frac{g_0^2}{16\pi ^2}}C_F`$ $`[({\displaystyle \frac{5}{3}}+(1\alpha ))\mathrm{log}a^2p^2`$
$`+V_{v_2,e}^{\alpha =1}(1\alpha )\mathrm{\hspace{0.17em}6.850272}+T_{v_2,e}]O_{v_2,e}^{\mathrm{tree}},`$
$`O_{v_3,d}^{\mathrm{proper}}={\displaystyle \frac{g_0^2}{16\pi ^2}}C_F`$ $`[({\displaystyle \frac{19}{6}}+(1\alpha ))\mathrm{log}a^2p^2`$
$`+V_{v_3,d}^{\alpha =1}(1\alpha )\mathrm{\hspace{0.17em}7.369693}+T_{v_3,d}]O_{v_3,d}^{\mathrm{tree}}.`$
The contributions of the sails and vertices, $`V_O`$, have been for convenience separated in the Feynman gauge results ($`V_{v_2,d}^{\alpha =1}`$, $`V_{v_2,e}^{\alpha =1}`$ and $`V_{v_3,d}^{\alpha =1}`$), tabulated in Table 1 for various values of the parameter $`\rho `$, and the remaining parts proportional to $`(1\alpha )`$, which are instead independent of $`\rho `$. The analytic expressions of the latter are very complicated functions of $`\rho `$ containing thousands of terms, and the numerical cancellation of this dependence is a rather non-trivial check of our computations. These numbers are also independent of the lattice representation, and furthermore they have the same value as with the Wilson action , so they seem to be to a certain extent independent of the particular fermion action chosen.
The contributions $`T_O`$ of the operator tadpoles (d in Fig. 1) are shown in Table 2, which contains also their results for the other operator, $`O_{v_3,e}`$. However, since this operator is not multiplicatively renormalized on the lattice, we consider in its place the two operators $`O_A`$ and $`O_B`$ introduced in Sect. 2. They mix with each other and we can write their renormalization as <sup>4</sup><sup>4</sup>4In Refs. another choice for the two operators that mix was made, and $`O_{v_3,e}`$ and an operator of mixed symmetry were considered instead of $`O_A`$ and $`O_B`$.
$`\widehat{O}_A`$ $`=`$ $`Z_{AA}O_A+Z_{AB}O_B`$ (37)
$`\widehat{O}_B`$ $`=`$ $`Z_{BA}O_A+Z_{BB}O_B.`$
In terms of the bare operators $`O_A`$ and $`O_B`$, the renormalized operator $`\widehat{O}_{v_3,e}`$ appearing in the DIS light-cone expansion will be
$$\widehat{O}_{v_3,e}=\frac{1}{3}\left(\widehat{O}_A+\widehat{O}_B\right)=\frac{1}{3}(Z_{AA}+Z_{BA})O_A+\frac{1}{3}(Z_{AB}+Z_{BB})O_B,$$
(38)
and since on the lattice $`Z_{AA}+Z_{BA}Z_{AB}+Z_{BB}`$, the result is that $`O_{v_3,e}`$ is not multiplicatively renormalized. In fact, the 1-loop contributions coming from the proper diagrams are in this case
$`O_A^{\mathrm{proper}}={\displaystyle \frac{g_0^2}{16\pi ^2}}C_F`$ $`[({\displaystyle \frac{7}{6}}+(1\alpha ))\mathrm{log}a^2p^2`$
$`+V_{AA}^{\alpha =1}(1\alpha )\mathrm{\hspace{0.17em}8.685568}+T_{AA}]O_A^{\mathrm{tree}}`$
$`+{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F`$ $`\left[\mathrm{log}a^2p^2+V_{AB}^{\alpha =1}+{\displaystyle \frac{1}{3}}(1\alpha )+T_{AB}\right]O_B^{\mathrm{tree}}`$
$`O_B^{\mathrm{proper}}={\displaystyle \frac{g_0^2}{16\pi ^2}}C_F`$ $`\left[2\mathrm{log}a^2p^2+V_{BA}^{\alpha =1}+{\displaystyle \frac{2}{3}}(1\alpha )+T_{BA}\right]O_A^{\mathrm{tree}}`$ (39)
$`+{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F`$ $`[({\displaystyle \frac{13}{6}}+(1\alpha ))\mathrm{log}a^2p^2`$
$`+V_{BB}^{\alpha =1}(1\alpha )\mathrm{\hspace{0.17em}7.703026}+T_{BB}]O_B^{\mathrm{tree}},`$
where the finite contributions of vertex and sails in the Feynman gauge are shown in Table 3.
We stress again that this mixing is a pure lattice effect, deriving from the breaking of the Lorentz group to the hypercubic group. In the continuum we do have $`Z_{AA}+Z_{BA}=Z_{AB}+Z_{BB}`$, and thus the operator $`O_{v_3,e}`$ is multiplicatively renormalized, and has the same $`Z`$ as the operator with different indices $`O_{v_3,d}`$, as they both belong to the same representation of the Lorentz group. We remind also that in the polarized case one encounters exactly the same situation with the same numerical results, i.e. $`Z_{AA}=Z_{A^5A^5}`$, $`Z_{AB}=Z_{A^5B^5}`$ etc., as we have explicitly verified.
To complete the computation of the renormalization factors, we have now to add to the proper diagrams the 1-loop amplitudes of the self-energy and tadpole of one leg which are proportional to $`ip/`$,
$$\mathrm{\Sigma }_1=\frac{g_0^2}{16\pi ^2}C_F\left[\alpha \mathrm{log}a^2p^2+S^{\alpha =1}+(1\alpha )\mathrm{\hspace{0.17em}4.792010}\right],$$
(40)
where the Feynman-gauge finite results $`S^{\alpha =1}`$ are given in Table 4. Putting everything together, we get the expressions of the renormalized operators on the lattice for overlap fermions, which for $`\rho =1`$ are:
$`\widehat{O}_{v_2,d}`$ $`=`$ $`\left[1{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left({\displaystyle \frac{8}{3}}\mathrm{log}a^2\mu ^253.25571+(1\alpha )\right)\right]O_{v_2,d}^{\mathrm{tree}}`$
$`\widehat{O}_{v_2,e}`$ $`=`$ $`\left[1{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left({\displaystyle \frac{8}{3}}\mathrm{log}a^2\mu ^252.82171+(1\alpha )\right)\right]O_{v_2,e}^{\mathrm{tree}}`$
$`\widehat{O}_{v_3,d}`$ $`=`$ $`\left[1{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left({\displaystyle \frac{25}{6}}\mathrm{log}a^2\mu ^268.99341+{\displaystyle \frac{3}{2}}(1\alpha )\right)\right]O_{v_3,d}^{\mathrm{tree}}`$
$`\widehat{O}_{v_3,e}`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left[1{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left({\displaystyle \frac{25}{6}}\mathrm{log}a^2\mu ^273.29777+{\displaystyle \frac{3}{2}}(1\alpha )\right)\right]O_A^{\mathrm{tree}}`$
$`+{\displaystyle \frac{1}{3}}\left[1{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left({\displaystyle \frac{25}{6}}\mathrm{log}a^2\mu ^267.83586+{\displaystyle \frac{3}{2}}(1\alpha )\right)\right]O_B^{\mathrm{tree}}.`$
We give here for comparison the numerical values obtained with the usual Wilson action, without any improvement. We have computed them again and checked them with the results in the literature . In this case however chiral invariance is broken and the renormalization factors for the unpolarized operators
$`\widehat{O}_{v_2,d}^{\mathrm{Wilson}}`$ $`=`$ $`\left[1{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left({\displaystyle \frac{8}{3}}\mathrm{log}a^2\mu ^23.16486+(1\alpha )\right)\right]O_{v_2,d}^{\mathrm{tree}}`$
$`\widehat{O}_{v_2,e}^{\mathrm{Wilson}}`$ $`=`$ $`\left[1{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left({\displaystyle \frac{8}{3}}\mathrm{log}a^2\mu ^21.88259+(1\alpha )\right)\right]O_{v_2,e}^{\mathrm{tree}}`$
$`\widehat{O}_{v_3,d}^{\mathrm{Wilson}}`$ $`=`$ $`\left[1{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left({\displaystyle \frac{25}{6}}\mathrm{log}a^2\mu ^219.00763+{\displaystyle \frac{3}{2}}(1\alpha )\right)\right]O_{v_3,d}^{\mathrm{tree}}`$
$`\widehat{O}_{v_3,e}^{\mathrm{Wilson}}`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left[1{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left({\displaystyle \frac{25}{6}}\mathrm{log}a^2\mu ^221.78271+{\displaystyle \frac{3}{2}}(1\alpha )\right)\right]O_A^{\mathrm{tree}}`$
$`+{\displaystyle \frac{1}{3}}\left[1{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left({\displaystyle \frac{25}{6}}\mathrm{log}a^2\mu ^218.46640+{\displaystyle \frac{3}{2}}(1\alpha )\right)\right]O_B^{\mathrm{tree}}`$
differ from the ones of the polarized operators
$`\widehat{O}_{a_2,d}^{\mathrm{Wilson}}`$ $`=`$ $`\left[1{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left({\displaystyle \frac{8}{3}}\mathrm{log}a^2\mu ^24.09933+(1\alpha )\right)\right]O_{a_2,d}^{\mathrm{tree}}`$
$`\widehat{O}_{a_2,e}^{\mathrm{Wilson}}`$ $`=`$ $`\left[1{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left({\displaystyle \frac{8}{3}}\mathrm{log}a^2\mu ^24.27705+(1\alpha )\right)\right]O_{a_2,e}^{\mathrm{tree}}`$
$`\widehat{O}_{a_3,d}^{\mathrm{Wilson}}`$ $`=`$ $`\left[1{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left({\displaystyle \frac{25}{6}}\mathrm{log}a^2\mu ^219.56159+{\displaystyle \frac{3}{2}}(1\alpha )\right)\right]O_{a_3,d}^{\mathrm{tree}}`$
$`\widehat{O}_{a_3,e}^{\mathrm{Wilson}}`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left[1{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left({\displaystyle \frac{25}{6}}\mathrm{log}a^2\mu ^222.39940+{\displaystyle \frac{3}{2}}(1\alpha )\right)\right]O_{A^5}^{\mathrm{tree}}`$
$`+{\displaystyle \frac{1}{3}}\left[1{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left({\displaystyle \frac{25}{6}}\mathrm{log}a^2\mu ^219.25837+{\displaystyle \frac{3}{2}}(1\alpha )\right)\right]O_{B^5}^{\mathrm{tree}}.`$
What one can notice is that the values of the renormalization factors for overlap fermions are in general much larger than for Wilson fermions. This seems to be true for most values of $`\rho `$, and is to trace largely to the quark self-energy and to the operators tadpoles. It can be observed that when in the Wilson formulation one improves the theory by adding the clover term to the action (with $`c_{sw}=1`$) and by canceling $`O(a)`$ effects on the operators at tree level, the renormalization factors become also somewhat larger, as in the examples below :
$`\widehat{O}_{v_2,d}^{\mathrm{Wilson},\mathrm{imp}}`$ $`=`$ $`\left[1{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left({\displaystyle \frac{8}{3}}\mathrm{log}a^2\mu ^215.816+(1\alpha )\right)\right]O_{v_2,d}^{\mathrm{tree}}`$
$`\widehat{O}_{v_3,d}^{\mathrm{Wilson},\mathrm{imp}}`$ $`=`$ $`\left[1{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left({\displaystyle \frac{25}{6}}\mathrm{log}a^2\mu ^229.815+{\displaystyle \frac{3}{2}}(1\alpha )\right)\right]O_{v_3,d}^{\mathrm{tree}}`$
$`\widehat{O}_{v_3,e}^{\mathrm{Wilson},\mathrm{imp}}`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left[1{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left({\displaystyle \frac{25}{6}}\mathrm{log}a^2\mu ^239.192+{\displaystyle \frac{3}{2}}(1\alpha )\right)\right]O_A^{\mathrm{tree}}`$
$`+{\displaystyle \frac{1}{3}}\left[1{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left({\displaystyle \frac{25}{6}}\mathrm{log}a^2\mu ^222.141+{\displaystyle \frac{3}{2}}(1\alpha )\right)\right]O_B^{\mathrm{tree}}.`$
One could then speculate whether the large renormalization factors that we have obtained for overlap fermions are related to the fact that the Neuberger action is improved and the $`Z`$s correspond to the ones of improved operators. A calculation with Wilson fermions with which to compare in this sense our overlap results would be one in which the operators are improved beyond tree level, but although we know in some cases the contributions of the operator counterterms that are needed for the full improvement, unfortunately not all their coefficients are yet determined .
We give finally the 1-loop results for the various matrix elements in the $`\overline{\mathrm{MS}}`$ scheme, so that one can complete the connection with the continuum as in Eq. (27). In the continuum there is only one renormalization constant for all possible operators (unpolarized and polarized) measuring the first moment, and only one for all possible operators measuring the second moment:
$`\widehat{O}_{v_2}^{\overline{\mathrm{MS}}}`$ $`=`$ $`\left[1{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left({\displaystyle \frac{8}{3}}\mathrm{log}{\displaystyle \frac{p^2}{\mu ^2}}{\displaystyle \frac{40}{9}}+(1\alpha )\right)\right]O_{v_2}^{\mathrm{tree}}`$ (45)
$`\widehat{O}_{v_3}^{\overline{\mathrm{MS}}}`$ $`=`$ $`\left[1{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left({\displaystyle \frac{25}{6}}\mathrm{log}{\displaystyle \frac{p^2}{\mu ^2}}{\displaystyle \frac{67}{9}}+{\displaystyle \frac{3}{2}}(1\alpha )\right)\right]O_{v_3}^{\mathrm{tree}}.`$
The connection of overlap lattice fermions with the continuum $`\overline{\mathrm{MS}}`$ is then given by the gauge-invariant factors
$`\widehat{O}_{v_2,d}^{\overline{\mathrm{MS}}}`$ $`=`$ $`\left[1{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left({\displaystyle \frac{8}{3}}\mathrm{log}a^2\mu ^248.81127\right)\right]O_{v_2,d}^{\mathrm{lat}}`$
$`\widehat{O}_{v_2,e}^{\overline{\mathrm{MS}}}`$ $`=`$ $`\left[1{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left({\displaystyle \frac{8}{3}}\mathrm{log}a^2\mu ^248.37727\right)\right]O_{v_2,e}^{\mathrm{lat}}`$
$`\widehat{O}_{v_3,d}^{\overline{\mathrm{MS}}}`$ $`=`$ $`\left[1{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left({\displaystyle \frac{25}{6}}\mathrm{log}a^2\mu ^261.54897\right)\right]O_{v_3,d}^{\mathrm{lat}}`$
$`\widehat{O}_{v_3,e}^{\overline{\mathrm{MS}}}`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left[1{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left({\displaystyle \frac{25}{6}}\mathrm{log}a^2\mu ^265.85333\right)\right]O_A^{\mathrm{lat}}`$
$`+{\displaystyle \frac{1}{3}}\left[1{\displaystyle \frac{g_0^2}{16\pi ^2}}C_F\left({\displaystyle \frac{25}{6}}\mathrm{log}a^2\mu ^260.39142\right)\right]O_B^{\mathrm{lat}}.`$
It looks increasingly difficult to go to higher moments, as the number of terms that are present in the FORM outputs and that need to be numerically integrated becomes very large. This has at the moment limited our calculations to second moment operators, but we hope to be able to compute the renormalization of third-moment operators in the near future.
Finally, we remark that all the renormalization constants presented here can be considered as computed in the unquenched case, because we limit ourselves to 1-loop computations with non-singlet quark operators, where internal quark loops never have the chance to come to play.
## Acknowledgment
I have enjoyed stimulating discussions with Robert Edwards and Leonardo Giusti. I would also like to thank Mark Alford for critically reading the manuscript. Both the FORM and Fortran computations have been done at MIT on a few Pentium III PCs running on Linux. This work has been supported in part by the U.S. Department of Energy (DOE) under cooperative research agreement DE-FC02-94ER40818.
## Appendix A Feynman rules for the Neuberger operator
We give here the Feynman rules which are needed to perform 1-loop calculations using the overlap-Dirac operator. An explicit derivations of these rules is given in .
The quark propagator is
$$S(k)=\frac{i_\mu \gamma _\mu \mathrm{sin}ak_\mu }{2\rho \left(\omega (k)+b(k)\right)}+\frac{a}{2\rho },$$
(47)
where
$`\omega (k)`$ $`=`$ $`{\displaystyle \frac{1}{a}}\sqrt{{\displaystyle \underset{\mu }{}}\mathrm{sin}^2ak_\mu +\left(2r{\displaystyle \underset{\mu }{}}\mathrm{sin}^2{\displaystyle \frac{ak_\mu }{2}}\rho \right)^2}`$
$`b(k)`$ $`=`$ $`{\displaystyle \frac{1}{a}}\left(2r{\displaystyle \underset{\mu }{}}\mathrm{sin}^2{\displaystyle \frac{ak_\mu }{2}}\rho \right).`$ (48)
The vertices needed for 1-loop calculations can be entirely given in terms of the vertices of the Wilson action
$`W_{1\mu }(p_1,p_2)`$ $`=`$ $`g_0\left(i\gamma _\mu \mathrm{cos}{\displaystyle \frac{a(p_1+p_2)_\mu }{2}}+r\mathrm{sin}{\displaystyle \frac{a(p_1+p_2)_\mu }{2}}\right)`$ (49)
$`W_{2\mu }(p_1,p_2)`$ $`=`$ $`{\displaystyle \frac{1}{2}}ag_0^2\left(i\gamma _\mu \mathrm{sin}{\displaystyle \frac{a(p_1+p_2)_\mu }{2}}+r\mathrm{cos}{\displaystyle \frac{a(p_1+p_2)_\mu }{2}}\right)`$ (50)
(where $`p_1`$ and $`p_2`$ are the quark momenta flowing in and out of the vertices) and of the quantity
$$X_0(p)=\frac{1}{a}\left(i\underset{\mu }{}\gamma _\mu \mathrm{sin}ap_\mu +2r\underset{\mu }{}\mathrm{sin}^2\frac{ap_\mu }{2}\rho \right).$$
(51)
The quark-quark-gluon vertex in the overlap theory has the expression
$`V_{1\mu }(p_1,p_2)`$ $`=`$ $`\rho {\displaystyle \frac{1}{\omega (p_1)+\omega (p_2)}}\times `$ (53)
$`\left[W_{1\mu }(p_1,p_2){\displaystyle \frac{1}{\omega (p_1)\omega (p_2)}}X_0(p_2)W_{1\mu }^{}(p_1,p_2)X_0(p_1)\right],`$
and the quark-quark-gluon-gluon vertex is
$`V_{2\mu \nu }(p_1,p_2)=\delta _{\mu \nu }\rho {\displaystyle \frac{1}{\omega (p_1)+\omega (p_2)}}\times `$
$`\left[W_{2\mu }(p_1,p_2){\displaystyle \frac{1}{\omega (p_1)\omega (p_2)}}X_0(p_2)W_{2\mu }^{}(p_1,p_2)X_0(p_1)\right]`$
$`+{\displaystyle \frac{1}{2}}\rho {\displaystyle \frac{1}{\omega (p_1)+\omega (p_2)}}{\displaystyle \frac{1}{\omega (p_1)+\omega (k)}}{\displaystyle \frac{1}{\omega (k)+\omega (p_2)}}\times `$
$`[X_0(p_2)W_{1\mu }^{}(p_2,k)W_{1\nu }(k,p_1)+W_{1\mu }(p_2,k)X_0^{}(k)W_{1\nu }(k,p_1)`$
$`+W_{1\mu }(p_2,k)W_{1\nu }^{}(k,p_1)X_0(p_1)`$
$`{\displaystyle \frac{\omega (p_1)+\omega (k)+\omega (p_2)}{\omega (p_1)\omega (k)\omega (p_2)}}X_0(p_2)W_{1\mu }^{}(p_2,k)X_0(k)W_{1\nu }^{}(k,p_1)X_0(p_1)].`$
## Appendix B Some analytic results
We give here the analytic results for the leg tadpole and for the vertex of the scalar current, in the Feynman gauge for $`r=1`$. In order to be able to write them in a compact form it is convenient to introduce the abbreviations
$$M_\lambda =\mathrm{cos}\frac{k_\lambda }{2},N_\lambda =\mathrm{sin}\frac{k_\lambda }{2},s_\lambda =\mathrm{sin}k_\lambda ,s^2=\underset{\lambda }{}s_\lambda ^2,$$
(55)
$$b=b(k)=2\underset{\lambda }{}\mathrm{sin}^2\frac{k_\lambda }{2}\rho ,$$
(56)
$$D=2\rho \left(\omega (k)+b(k)\right),$$
(57)
$$A=\frac{\omega ^2(k)}{\rho ^2}=1\frac{4}{\rho }\underset{\lambda }{}\mathrm{sin}^2\frac{k_\lambda }{2}+\frac{1}{\rho ^2}\left(\underset{\lambda }{}\mathrm{sin}^2k_\lambda +\left(2\underset{\lambda }{}\mathrm{sin}^2\frac{k_\lambda }{2}\right)^2\right).$$
(58)
The result for the 1-loop leg tadpole is then given by
$`{\displaystyle \frac{1}{2}}g_0^2{\displaystyle \frac{d^4k}{2\pi ^4}G(k)\left(1\frac{4}{\rho }\right)}`$ (59)
$`+g_0^2{\displaystyle }{\displaystyle \frac{d^4k}{2\pi ^4}}G(k){\displaystyle \frac{1}{\rho ^2(1+\sqrt{A})^2}}\times `$
$`{\displaystyle \underset{\lambda }{}}[M_\lambda ^2+N_\lambda ^2+(1+{\displaystyle \frac{1}{\sqrt{A}}})(s_\mu ^2+b(M_\mu ^2N_\mu ^2))`$
$`+{\displaystyle \frac{2+\sqrt{A}}{\rho \sqrt{A}}}(b(M_\lambda ^2N_\lambda ^2)+s^2)]`$
$`+g_0^2{\displaystyle }{\displaystyle \frac{d^4k}{2\pi ^4}}G^2(k){\displaystyle \frac{1}{\rho ^2(1+\sqrt{A})^2}}\times `$
$`{\displaystyle \underset{\lambda }{}}\left[2s_\mu ^2N_\lambda ^2+\left(1+{\displaystyle \frac{1}{\sqrt{A}}}\right)\left(2s_\mu ^2(b+2M_\mu ^2M_\lambda ^2)\right)\right].`$
The first term comes from the part of the $`V_2`$ vertex (A) containing $`W_2`$ and $`W_2^{}`$, and its value for $`\rho =1`$ is
$$g_0^2\frac{Z_0}{2}\left(1\frac{4}{\rho }\right)|_{\rho =1}=\frac{g_0^2}{16\pi ^2}\mathrm{\hspace{0.17em}36.69915},$$
(60)
while the 1-loop result for the whole leg tadpole is smaller:
$$\frac{g_0^2}{16\pi ^2}\mathrm{\hspace{0.17em}23.35975}.$$
(61)
Adding now the value $`\frac{g_0^2}{16\pi ^2}\mathrm{\hspace{0.17em}14.27088}`$ for the diagram e in Fig. 1, which is much harder to compute by hand and would be a very lengthy analytic expression anyway, gives the result $`\frac{g_0^2}{16\pi ^2}\mathrm{\hspace{0.17em}37.63063}`$ for the complete self-energy in the Feynman gauge for $`\rho =1`$.
The results for the 1-loop vertex diagram of the scalar operator (the sails are not present in this case) can be written in the form
$$g_0^2\frac{d^4k}{2\pi ^4}G(pk)\frac{1}{(1+\sqrt{A})^2}\left[\left(\frac{s^2}{D^2}+\frac{1}{4\rho ^2}\right)X+\frac{1}{\rho D}Y\right],$$
(62)
for small $`p`$, where
$`X`$ $`=`$ $`{\displaystyle \underset{\lambda }{}}[(M_\lambda ^2N_\lambda ^2)+{\displaystyle \frac{1}{\rho \sqrt{A}}}2b(M_\lambda ^2+N_\lambda ^2)`$
$`+{\displaystyle \frac{1}{\rho ^2A}}((s^2b^2)(M_\lambda ^2N_\lambda ^2)+2bs^2)]`$
$`Y`$ $`=`$ $`{\displaystyle \underset{\lambda }{}}[s^2+{\displaystyle \frac{1}{\rho \sqrt{A}}}2s^2(M_\lambda ^2+N_\lambda ^2)`$ (63)
$`+{\displaystyle \frac{1}{\rho ^2A}}s^2(2b(M_\lambda ^2N_\lambda ^2)+s^2b^2)].`$
We have checked that all these results obtained by hand correspond with the outputs of the FORM programs.
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# Three point functions and the effective lagrangian for the chial primary fields in 𝐷=4 supergravity on 𝐴𝑑𝑆₂×𝑆²
## I Introduction
Among all known examples of the $`AdS`$/CFT correspondence , , ,, the least understood is the $`AdS_2`$/CFT<sub>1</sub> case. The $`D=1`$ conformal field theory (CFT), or conformal quantum mechanics (CQM), has not been formulated and therefore no quantitative comparison between the two sides of the duality has been made. However, there have been conjectures,, that the dual $`CFT`$ is given by the $`n`$-particle, $`𝒩=4`$ superconformal Calogero model, which has yet to be constructed for arbitrary $`n`$. By going to the second-quantized formulation, this becomes a $`1+1`$ dimensional field theory. (See also ref., , , .)
As a first step toward investigating this conjecture we consider the supergravity side of the duality. The low energy limit of type II(A or B) string theory on $`T^6`$ is the $`𝒩=8`$ supergravity theory of Cremmer and Julia . Off-shell this theory has an $`SO(8)`$ symmetry while on-shell the symmetry is enhanced to an $`E_{7(7)}`$ duality. We consider the fluctuations of the fields around a classical configuration called Bertotti-Robinson solution, which gives the $`AdS_2\times S^2`$ spacetime. This solution is the near-horizon limit of four D-branes intersecting over a string. The equations for the fluctuations of the fields were considered up to linear order in ref., to obtain the physical spectrum of the theory. <sup>*</sup><sup>*</sup>* See also ref., for the spectrum of the minimal $`D=4,𝒩=2`$ supergravity on $`AdS_2\times S^2`$ and the $`D=10`$, IIA supergravity on ‘quasi’ $`AdS_2\times S^8`$, respectively.
In this paper, we compute the cubic interactions among the bosonic chiral primary fields. Such computations were done for $`D=10`$ IIB supergravity on $`AdS_5\times S^5`$ and $`D=6`$ supergravity on $`AdS_3\times S^3`$ in ref. and , respectively. The computation involves the standard elements of previous works, namely nonlinear field redefinitions, which then lead to three point functions of chiral operators. The results have some similarities in form with the examples in the higher dimensional $`AdS`$ spaces. It would be interesting to see whether one can make computations in a supersymmteric Calogero model, to be compared with these results.
## II Supergravity in $`D=4`$
In this section, we review following , the equations of motion of $`D=4`$, $`𝒩=8`$ supergravity. We mainly follow the four dimensional formalism and notation of ref. . We are interested in the equation of motion for the metric $`g_{\widehat{\mu }\widehat{\nu }}`$, vectors $`B_{\widehat{\mu }}^{AB}`$, and scalars $`W_{ABCD}`$, where the capital roman letters are $`SU(8)`$ indices. The scalars are packaged in a $`E_{7(7)}`$ matrix $`𝒱`$ parametrizing the coset $`\frac{E_{7(7)}}{SU(8)}`$,
$$_{\widehat{\mu }}𝒱𝒱^1=\left(\begin{array}{cc}Q_{\widehat{\mu }[A}^{[C}\delta _{B]}^{D]}& P_{\widehat{\mu }ABCD}\\ \overline{P}_{\widehat{\mu }}^{ABCD}& \overline{Q}_{\widehat{\mu }[C}^{[A}\delta _{D]}^{B]}.\end{array}\right)$$
(1)
$`P`$’s parametrize the coset manifold and can be expressed in terms of the scalar fields $`W`$, and $`Q`$’s are the $`SU(8)`$ gauge fields, but for our purposes it is enough to note that the $`SU(8)`$ gauge symmetry can be fixed to the so called symmetric gauge where $`𝒱=\mathrm{exp}(X)`$,
$$X=\left(\begin{array}{cc}0& W_{ABCD}\\ \overline{W}^{ABCD}& 0\end{array}\right),$$
(2)
and $`W_{ABCD}`$ is complex, completely antisymmetric in $`A,B,C,D`$, and satisfies the constraint
$`\overline{W}^{ABCD}={\displaystyle \frac{1}{24}}ϵ^{ABCDEFGH}W_{EFGH}.`$ (3)
We will use the following notation for indices: $`\widehat{\mu },\widehat{\nu }=0,1,2,3`$ are $`D=4`$ coordinates, $`\lambda ,\mu ,\nu \mathrm{}=0,1`$ are $`AdS_2`$ coordinates and $`\alpha ,\beta ,\gamma \mathrm{}=2,3`$ are $`S^2`$ coordinates.
The bosonic part of the $`D=4`$ supergravity action is:
$$=\left(\frac{1}{4}eR(\omega ,e)+\frac{1}{8}eF_{\widehat{\mu }\widehat{\nu }}^{MN}(B)\stackrel{~}{H}_{MN}^{\widehat{\mu }\widehat{\nu }}(B,𝒱)\frac{1}{24}eP_{\widehat{\mu }ABCD}\overline{P}^{\widehat{\mu }ABCD}\right)$$
(4)
where
$$F_{\widehat{\mu }\widehat{\nu }}^{AB}=2_{[\widehat{\mu }}B_{\widehat{\nu }]}^{AB}$$
(5)
$`\stackrel{~}{H}_{MN}^{\widehat{\mu }\widehat{\nu }}`$ are defined in the Appendix 1. For our purpose, we only need the leading expansion of the $`\stackrel{~}{H}`$ and $`P`$ in $`W^{ABCD}`$:
$`G_{\widehat{\mu }\widehat{\nu }}^{MN}\stackrel{~}{H}_{\widehat{\mu }\widehat{\nu }MN}^{(B)}`$ $`=`$ $`G_{\widehat{\mu }\widehat{\nu }}^{MN}(1+W+\overline{W}+W^2+\overline{W}^2`$ (8)
$`{\displaystyle \frac{1}{3}}W\overline{W}W{\displaystyle \frac{1}{3}}\overline{W}W\overline{W}+W^3+\overline{W}^3)_{MNPQ}G^{\widehat{\mu }\widehat{\nu }PQ}`$
$`+iG_{\widehat{\mu }\widehat{\nu }}^{MN}(W\overline{W}+W^2\overline{W}^2+O(W^3))_{MNPQ}\stackrel{~}{G}^{\widehat{\mu }\widehat{\nu }PQ}+O(W^4),`$
$`P_{\widehat{\mu }ABCD}`$ $`=`$ $`_{\widehat{\mu }}W_{ABCD}+O(W^3).`$ (9)
## III equations of motion for the bosonic chiral primary fields
We consider the classical configuration called Bertotti-Robinson solution,
$`ds^2`$ $`=`$ $`{\displaystyle \frac{1}{z^2}}(dx_0^2+dz^2)+d\mathrm{\Omega }_2^2,`$ (10)
$`R_{\mu \lambda \nu \sigma }`$ $`=`$ $`(g_{\mu \nu }g_{\lambda \sigma }g_{\mu \sigma }g_{\lambda \nu })`$ (11)
$`R_{\alpha \gamma \beta \delta }`$ $`=`$ $`(g_{\alpha \beta }g_{\gamma \delta }g_{\alpha \delta }g_{\gamma \beta })`$ (12)
$`\overline{F}_{\alpha \beta }^{12}`$ $`=`$ $`ϵ_{\alpha \beta },`$ (13)
$`\overline{F}_{\alpha \beta }^{AB}`$ $`=`$ $`0(A1,2),`$ (14)
$`W_{ABCD}`$ $`=`$ $`0`$ (15)
which breaks the $`SU(8)`$ internal symmetry into $`SU(6)\times SU(2)`$. The bulk fields of interest are the fluctuations about this background,
$`g_{\widehat{\mu }\widehat{\nu }}`$ $`=`$ $`\overline{g}_{\widehat{\mu }\widehat{\nu }}+h_{\widehat{\mu }\widehat{\nu }},`$ (16)
$`F_{\widehat{\mu }\widehat{\nu }}^{AB}`$ $`=`$ $`\overline{F}^{AB_{\widehat{\mu }\widehat{\nu }}}+2_{[\widehat{\mu }}b_{\widehat{\nu }]}^{AB}`$ (17)
and $`W_{ABCD}`$. One can consider the classical equations for these fluctuations to linear order and organize the spectrum into the multiplets of $`SU(6)\times SU(2)`$,. In this paper we expand the equations of motion up to the second order in these fluctuations to obtain the interaction terms.
To find the chiral primary fields on $`AdS_2`$ we expand the fields in spherical harmonics on $`S^2`$. The expansions are quite simple in this case as all harmonic functions on the 2-sphere can be expressed in terms of just the scalar spherical harmonics $`Y_{lm}`$. $`l`$ is the quantum number labeling the Casimir of the representation,
$$_\alpha ^\alpha Y_{lm}=l(l+1)Y_{lm}.$$
(18)
The expansions of the bosonic fluctuations are then given by (denoting the $`l,m`$ indices collectively by $`I`$)
$`h_{\mu \nu }`$ $`=`$ $`{\displaystyle H_{\mu \nu }^IY_I}`$ (20)
$`h_{\mu \alpha }`$ $`=`$ $`{\displaystyle (B_{1\mu }^I_\alpha Y_I+B_{2\mu }^Ie_{\alpha \beta }^\beta Y_I)}`$ (21)
$`h_{\alpha \beta }`$ $`=`$ $`{\displaystyle (\varphi _1^I_\alpha _\beta Y_I+\varphi _2^Ie_{(\alpha }^\gamma _{\beta )}_\gamma Y_I+\varphi _3^Ig_{\alpha \beta }Y_I)}`$ (22)
$`b_\mu ^{AB}`$ $`=`$ $`{\displaystyle b_\mu ^{(I)AB}Y_I}`$ (23)
$`b_\alpha ^{AB}`$ $`=`$ $`{\displaystyle (b_1^{(I)AB}_\alpha Y_I+b_2^{(I)AB}e_{\alpha \beta }^\beta Y_I)}`$ (24)
$`W_{ABCD}`$ $`=`$ $`{\displaystyle W_{ABCD}^IY_I}.`$ (25)
Before substituting the expansions into the equations of motion we can first simplify the expansions by fixing gauge symmetries by imposingSince we project out only the physical modes, actually it is enough to impose $`\varphi _1=\varphi _2=0`$. From now on, we write $`\varphi _3\varphi ,B_{2\mu }B_\mu ,b_2b`$.
$$\varphi _1^I=\varphi _2^I=B_{1\mu }^I=b_1^{(I)AB}=0.$$
(26)
At the linearized level, the bosonic fields can be decomposed into the eigenstates of the $`AdS_2`$ Laplacian,,
$`\varphi ^I`$ $`=`$ $`2lT^I+2(l1)\stackrel{~}{T}^I`$ (27)
$`b^{(I)[12]}`$ $`=`$ $`T^I\stackrel{~}{T}^I`$ (28)
$`a^{(I)[12]}`$ $``$ $`ϵ^{\mu \nu }_\mu b_\nu ^{(I)[12]}=(l+1)S^I+l\stackrel{~}{S}^I`$ (29)
$`B^I`$ $``$ $`ϵ^{\mu \nu }_\mu B_\nu ^I=2S^I+2\stackrel{~}{S}^I`$ (30)
$`a^{(I)[MN]}`$ $``$ $`ϵ^{\mu \nu }_\mu b_\nu ^{(I)[MN]}=(l1)U^{(I)[MN]}+(l+2)\stackrel{~}{U}^{(I)[MN]}(M,N1,2)`$ (31)
$`(W\overline{W})^{(I)12MN}`$ $`=`$ $`i_x^2(U^{(I)[MN]}\stackrel{~}{U}^{(I)[MN]})`$ (32)
$`b^{(I)[MN]}`$ $`=`$ $`V^{(I)[MN]}+\stackrel{~}{V}^{(I)[MN]}(M,N1,2)`$ (33)
$`(W+\overline{W})^{(I)12MN}`$ $`=`$ $`lV^{(I)[MN]}+(l+1)\stackrel{~}{V}^{(I)[MN]}`$ (34)
where these fields satisify the equations of the form:
$`(_x^2l(l1))A^I+Q_A^I`$ $`=`$ $`0`$ (35)
$`(_x^2(l+1)(l+2))\stackrel{~}{A}^I+\stackrel{~}{Q}_{\stackrel{~}{A}}^I=0`$ (37)
Here, $`A`$ and $`\stackrel{~}{A}`$ stand for $`T,S,U,V`$ and $`\stackrel{~}{T},\stackrel{~}{S},\stackrel{~}{U},\stackrel{~}{V}`$ respectively, and $`Q_A`$ and $`Q_{\stackrel{~}{A}}`$ are the second order corrections. Also, we note that the fields $`H_{\mu \nu }`$ are not independent degrees of freedom and they are completely determined by the equations,,
$`((_x^2+2l(l+1))H_{\mu \nu }^I2_{(\mu }^\lambda H_{\nu )\lambda }^I+(_\mu _\nu `$ $``$ $`g_{\mu \nu }(_x^2+1l(l+1)))H^I`$ (39)
$`+g_{\mu \nu }^\lambda ^\rho H_{\lambda \rho }^I+2(_\mu _\nu g_{\mu \nu }(_x^21{\displaystyle \frac{1}{2}}l(l+1)))\varphi ^I)`$ $`=`$ $`4g_{\mu \nu }l(l+1)b^I+(\mathrm{higher}\mathrm{order})`$ (40)
$`\left(^\nu H_{\mu \nu }^I_\mu H^I_\mu \varphi ^I\right)`$ $`=`$ $`4_\mu b^I+(\mathrm{higher}\mathrm{order})`$ (41)
$`\left(\left(_x^2+4\right)\varphi ^I+\left(_x^21l(l+1)\right)H^I^\mu ^\nu H_{\mu \nu }^I\right)`$ $`=`$ $`4l(l+1)b^I+(\mathrm{higher}\mathrm{order})`$ (42)
$`H^I`$ $`=`$ $`(\mathrm{quadratic}\mathrm{order}).`$ (43)
By making an ansatz, one can easily find the solution to the equations above:
$$H_{\mu \nu }^I=\frac{1}{l+1}\left(2l(l1)g_{\mu \nu }T^I+4_\mu _\nu T^I\right)+(\mathrm{higher}\mathrm{order}\mathrm{corrections})$$
(44)
Since we are interested only in the three-point function of chiral primary fields, only $`Q_T,Q_S,Q_U,Q_V`$ are of interest. We also put any non-chiral primary fields appearing in $`Q`$’s to be zero, and substitute (44) for $`H_{\mu \nu }`$, neglecting the higher order corrections. The detailed form of the $`Q`$’s are rather complicated and not interesting at this stage. They are of the generic form:
$$Q_A=\alpha ^\mu _\nu B_\mu _\nu C+\beta ^\mu B_\mu C+\gamma BC+\mathrm{},$$
(45)
where $`A`$, $`B`$, $`C`$ are chiral primary fields. We see that there are terms involving derivatives of the fields in $`Q`$’s. These terms can be removed by nonlinear redefinitions of the fields, which does not change the equations (37) at the linear level. It is easy to see that this can be done by redefining,
$$AA\frac{\alpha }{2c}(B)(C)\frac{1}{2c}(\alpha +(1+\mathrm{\Gamma })\beta )BC,$$
(46)
where $`\mathrm{\Gamma }\frac{1}{2}(l(l1)l_1(l_11)l_2(l_21))`$ with $`l,l_1,l_2`$ being the total angular momentum quantum numbers of $`A,B,C`$, respectively. After these redefinition, the linear term
$$(_x^2l(l1))A$$
(47)
generates additional term which removes the derivative terms, and $`Q_A`$ becomes:
$$Q_A=(\gamma +\mathrm{\Gamma }(\beta +(1+\mathrm{\Gamma })\alpha ))BC\mathrm{}$$
(48)
We then get
$`Q_T`$ $`=`$ $`{\displaystyle \frac{(l(l^21)+l_1(l_1^21)+l_2(l_2^21))}{2l(2l+1)(l1)(l_1+1)(l_2+1)}}\alpha \alpha _1\alpha _2(\mathrm{\Sigma }^21)\stackrel{~}{C}(I;I_1,I_2)T^2`$ (52)
$`{\displaystyle \frac{(l(l^21)l_1(l_1^21)l_2(l_2^21))\stackrel{~}{C}(I;I_1,I_2)}{2l(2l+1)(l1)}}\alpha \alpha _1\alpha _2(\mathrm{\Sigma }^21)S^2`$
$`+{\displaystyle \frac{(l_1+l_2)(l_11)(l_21)}{4l(2l+1)(l1)}}\alpha \alpha _1\alpha _2(\mathrm{\Sigma }^21)\stackrel{~}{C}(I;I_1,I_2)ϵ^{12ABCDEF}U_{CD}U_{EF}`$
$`+{\displaystyle \frac{(l_1+l_2)}{4l(2l+1)(l1)}}\alpha \alpha _1\alpha _2(\mathrm{\Sigma }^21)\stackrel{~}{C}(I;I_1,I_2)ϵ^{12ABCDEF}V_{CD}V_{EF}`$
$`Q_S`$ $`=`$ $`{\displaystyle \frac{(l(l^21)+l_1(l_1^21)l_2(l_2^21))}{2l(l^21)(2l+1)(l_2+1)}}\alpha \alpha _1\alpha _2(\mathrm{\Sigma }^21)\stackrel{~}{C}(I;I_1,I_2)S^{(1)}T^{(2)}`$ (55)
$`{\displaystyle \frac{(l_1l_2)(l_21)}{4l(l^21)(2l+1)}}\alpha \alpha _1\alpha _2(\mathrm{\Sigma }^21)\stackrel{~}{C}(I;I_1,I_2)V^{(1)}U^{(2)}`$
$`Q_U^{AB}`$ $`=`$ $`{\displaystyle \frac{(ll_2)}{2l(l1)(2l+1)}}\alpha \alpha _1\alpha _2(\mathrm{\Sigma }^21)\stackrel{~}{C}(I;I_1,I_2)S^{(1)}V^{(2)AB}`$ (59)
$`{\displaystyle \frac{l_11}{8l(l1)(2l+1)}}\alpha \alpha _1\alpha _2(\mathrm{\Sigma }^21)\stackrel{~}{C}(I;I_1,I_2)ϵ^{12ABCDEF}U_{CD}^{(1)}V_{EF}^{(2)}`$
$`+{\displaystyle \frac{(l+l_1)(l_11)}{2l(l1)(2l+1)(1+l_2)}}\alpha \alpha _1\alpha _2(\mathrm{\Sigma }^21)\stackrel{~}{C}(I;I_1,I_2)U^{(1)AB}T^{(2)}`$
$`Q_V^{AB}`$ $`=`$ $`{\displaystyle \frac{(l_2l)(l_21)}{2l(2l+1)}}\alpha \alpha _1\alpha _2(\mathrm{\Sigma }^21)\stackrel{~}{C}(I;I_1,I_2)S^{(1)}U^{(2)AB}`$ (64)
$`{\displaystyle \frac{(l_11)(l_21)}{16l(2l+1)}}\alpha \alpha _1\alpha _2(\mathrm{\Sigma }^21)\stackrel{~}{C}(I;I_1,I_2)ϵ^{12ABCDEF}U_{CD}U_{EF}`$
$`+{\displaystyle \frac{1}{16l(2l+1)}}\alpha \alpha _1\alpha _2(\mathrm{\Sigma }^21)\stackrel{~}{C}(I;I_1,I_2)ϵ^{12ABCDEF}V_{CD}V_{EF}`$
$`+{\displaystyle \frac{(l+l_1)}{2l(2l+1)(1+l_2)}}\alpha \alpha _1\alpha _2(\mathrm{\Sigma }^21)\stackrel{~}{C}(I;I_1,I_2)V^{(1)AB}T^{(2)}`$
at the quadratic level, where $`\alpha l_1+l_2l`$, $`\alpha _1l+l_2l_1`$, $`\alpha _2l+l_1l_2`$, $`\mathrm{\Sigma }l+l_1+l_2`$, and
$$\stackrel{~}{C}(I;I_1,I_2)Y_I^{}Y_{I_1}Y_{I_2}.$$
(66)
We also used the abbreviation $`A^{(i)}A^{I_i}`$ and $`A^2A^{I_1}A^{I_2}`$ for any field $`A`$.
## IV The effective action and the three-point function function
The equations in the previous section can be derived from the truncated Lagrangian
$``$ $`=`$ $`\left({\displaystyle \frac{1}{4}}eR(\omega ,e)+{\displaystyle \frac{1}{8}}eF_{\widehat{\mu }\widehat{\nu }}^{MN}(B)\stackrel{~}{H}_{MN}^{\widehat{\mu }\widehat{\nu }}(B,𝒱){\displaystyle \frac{1}{24}}eP_{\widehat{\mu }ABCD}\overline{P}^{\widehat{\mu }ABCD}\right)`$ (67)
$`=`$ $`{\displaystyle \frac{(2l+1)l(l^21)}{4}}T(_x^2l(l1))T+{\displaystyle \frac{(2l+1)l(l^21)}{2}}S(_x^2l(l1))S`$ (84)
$`+{\displaystyle \frac{l(2l+1)}{4}}V^{AB}(_x^2l(l1))V_{AB}+{\displaystyle \frac{l(l1)^2(2l+1)}{4}}U^{AB}(_x^2l(l1))U_{AB}`$
$`+{\displaystyle \frac{(l_1(l_1^21)+l_2(l_2^21)+l_3(l_3^21))}{3(l_1+1)(l_2+1)(l_3+1)}}\alpha _1\alpha _2\alpha _3(\mathrm{\Sigma }^21)C(I_1,I_2,I_3)T^3`$
$`+{\displaystyle \frac{(l_1(l_1^21)+l_2(l_2^21)l_3(l_3^21))}{(l_3+1)}}\alpha _1\alpha _2\alpha _3(\mathrm{\Sigma }^21)C(I_1,I_2,I_3)S^{(1)}S^{(2)}T^{(3)}`$
$`(l_1l_2)(l_21)\alpha _1\alpha _2\alpha _3(\mathrm{\Sigma }^21)C(I_1,I_2,I_3)V^{(1)AB}U_{AB}^{(2)}S^{(3)}`$
$`+{\displaystyle \frac{1}{24}}\alpha _1\alpha _2\alpha _3(\mathrm{\Sigma }^21)C(I_1,I_2,I_3)ϵ^{12ABCDEF}V_{AB}^{(1)}V_{CD}^{(2)}V_{EF}^{(3)}`$
$`+{\displaystyle \frac{l_1+l_2}{2(l_3+1)}}\alpha _1\alpha _2\alpha _3(\mathrm{\Sigma }^21)C(I_1,I_2,I_3)V^{(1)AB}V_{AB}^{(2)}T^{(3)}`$
$`{\displaystyle \frac{(l_11)(l_21)}{8}}\alpha _1\alpha _2\alpha _3(\mathrm{\Sigma }^21)C(I_1,I_2,I_3)ϵ^{12ABCDEF}U_{AB}^{(1)}U_{CD}^{(2)}V_{EF}^{(3)}`$
$`+{\displaystyle \frac{(l_11)(l_21)(l_1+l_2)}{2(l_3+1)}}\alpha _1\alpha _2\alpha _3(\mathrm{\Sigma }^21)C(I_1,I_2,I_3)U^{(1)AB}U_{AB}^{(2)}T^{(3)}.`$
with
$$C(I_1,I_2,I_3)Y^{I_1}Y^{I_2}Y^{I_3}.$$
(85)
Here the normalizations were fixed by directly substituting the expression (34) into the lagrangian (4) and evaluating the leading terms for some fields.
To compute from (84) the 2- and 3-point functions of chiral primary operators of the boundary theory, we apply the formulae derived, for instance, in ref.. From Eq.(17) and the correction factor in Eq.(95) of ref., we read off the tree-level two-point functions to beAlthough the following results are the correlation functions of the chiral primary operators of the boundary theory which couple to the bulk fields $`S,T,U,V`$ we will still denote them by $`S,T,U,V`$ for the simplicity of the notations.
$`T^{I_1}(x)T^{I_2}(y)`$ $`=`$ $`{\displaystyle \frac{(2l+1)(l^21)}{2}}{\displaystyle \frac{1}{\pi ^{1/2}}}{\displaystyle \frac{\mathrm{\Gamma }(l+1)}{\mathrm{\Gamma }(l1/2)}}(2l1){\displaystyle \frac{\delta ^{I_1I_2}}{|xy|^{2l}}}.`$ (86)
$`S^{I_1}(x)S^{I_2}(y)`$ $`=`$ $`(2l+1)(l^21){\displaystyle \frac{1}{\pi ^{1/2}}}{\displaystyle \frac{\mathrm{\Gamma }(l+1)}{\mathrm{\Gamma }(l1/2)}}(2l1){\displaystyle \frac{\delta ^{I_1I_2}}{|xy|^{2l}}}.`$ (87)
$`U^{(I_1)AB}(x)U^{(I_2)CD}(y)`$ $`=`$ $`(l1)^2(2l+1){\displaystyle \frac{1}{\pi ^{1/2}}}{\displaystyle \frac{\mathrm{\Gamma }(l+1)}{\mathrm{\Gamma }(l1/2)}}(2l1){\displaystyle \frac{(\delta _{AC}\delta _{BD}\delta _{AD}\delta _{BC})\delta ^{I_1I_2}}{|xy|^{2l}}}.`$ (88)
$`V^{(I_1)AB}(x)V^{(I_2)CD}(y)`$ $`=`$ $`(2l+1){\displaystyle \frac{1}{\pi ^{1/2}}}{\displaystyle \frac{\mathrm{\Gamma }(l+1)}{\mathrm{\Gamma }(l1/2)}}(2l1){\displaystyle \frac{(\delta _{AC}\delta _{BD}\delta _{AD}\delta _{BC})\delta ^{I_1I_2}}{|xy|^{2l}}}.`$ (89)
From Eq.(25) of the same paper we derive that
$`T^{(1)}(x)T^{(2)}(y)T^{(3)}(z)`$ $`=`$ $`{\displaystyle \frac{2^5}{\pi }}{\displaystyle \frac{\underset{i}{}\mathrm{\Gamma }(\frac{\alpha _i}{2}+1)\mathrm{\Gamma }(\frac{1}{2}\mathrm{\Sigma }+\frac{3}{2})(l_1(l_1^21)+l_2(l_2^21)+l_3(l_3^21))}{_i^3\left(\mathrm{\Gamma }(l_i1/2)(l_i+1)\right)|xy|^{\alpha _3}|yz|^{\alpha _1}|zx|^{\alpha _2}}}C(I_1,I_2,I_3)`$ (90)
$`S^{(1)}S^{(2)}T^{(3)}`$ $`=`$ $`{\displaystyle \frac{2^5}{\pi }}{\displaystyle \frac{\underset{i}{}\mathrm{\Gamma }(\frac{\alpha _i}{2}+1)\mathrm{\Gamma }(\frac{1}{2}\mathrm{\Sigma }+\frac{3}{2})(l_1(l_1^21)+l_2(l_2^21)+l_3(l_3^21))}{_i^3\left(\mathrm{\Gamma }(l_i1/2)\right)(l_3+1)|xy|^{\alpha _3}|yz|^{\alpha _1}|zx|^{\alpha _2}}}C(I_1,I_2,I_3)`$ (92)
$`V^{(1)AB}U^{(2)CD}S^{(3)}`$ $`=`$ $`{\displaystyle \frac{2^5}{\pi }}{\displaystyle \frac{\underset{i}{}\mathrm{\Gamma }(\frac{\alpha _i}{2}+1)\mathrm{\Gamma }(\frac{1}{2}\mathrm{\Sigma }+\frac{3}{2})(l_1l_2)(l_21)}{_i^3\left(\mathrm{\Gamma }(l_i1/2)\right)|xy|^{\alpha _3}|yz|^{\alpha _1}|zx|^{\alpha _2}}}`$ (95)
$`\times (\delta ^{AC}\delta ^{BD}\delta ^{AD}\delta ^{BC})C(I_1,I_2,I_3)`$
$`V_{AB}^{(1)}V_{CD}^{(2)}V_{EF}^{(3)}`$ $`=`$ $`{\displaystyle \frac{2^5}{\pi }}{\displaystyle \frac{\underset{i}{}\mathrm{\Gamma }(\frac{\alpha _i}{2}+1)\mathrm{\Gamma }(\frac{1}{2}\mathrm{\Sigma }+\frac{3}{2})}{_i^3\left(\mathrm{\Gamma }(l_i1/2)\right)|xy|^{\alpha _3}|yz|^{\alpha _1}|zx|^{\alpha _2}}}C(I_1,I_2,I_3)ϵ_{12ABCDEF}`$ (97)
$`V_{AB}^{(1)}V_{CD}^{(2)}T^{(3)}`$ $`=`$ $`{\displaystyle \frac{2^6}{\pi }}{\displaystyle \frac{(l_1+l_2)\underset{i}{}\mathrm{\Gamma }(\frac{\alpha _i}{2}+1)\mathrm{\Gamma }(\frac{1}{2}\mathrm{\Sigma }+\frac{3}{2})(\delta _{AC}\delta _{BD}\delta _{AD}\delta _{BC})}{(l_3+1)_i^3\left(\mathrm{\Gamma }(l_i1/2)\right)|xy|^{\alpha _3}|yz|^{\alpha _1}|zx|^{\alpha _2}}}C(I_1,I_2,I_3)`$ (99)
$`U_{AB}^{(1)}U_{CD}^{(2)}V_{EF}^{(3)}`$ $`=`$ $`{\displaystyle \frac{2^5}{\pi }}{\displaystyle \frac{(l_11)(l_21)\underset{i}{}\mathrm{\Gamma }(\frac{\alpha _i}{2}+1)\mathrm{\Gamma }(\frac{1}{2}\mathrm{\Sigma }+\frac{3}{2})}{\left(_i^3\mathrm{\Gamma }(l_i1/2)\right)|xy|^{\alpha _3}|yz|^{\alpha _1}|zx|^{\alpha _2}}}C(I_1,I_2,I_3)ϵ_{12ABCDEF}`$ (101)
$`U^{(1)AB}U^{(2)CD}T^{(3)}`$ $`=`$ $`{\displaystyle \frac{2^6}{\pi }}{\displaystyle \frac{(l_11)(l_21)(l_1+l_2)\underset{i}{}\mathrm{\Gamma }(\frac{\alpha _i}{2}+1)\mathrm{\Gamma }(\frac{1}{2}\mathrm{\Sigma }+\frac{3}{2})}{(l_3+1)\left(_i^3\mathrm{\Gamma }(l_i1/2)\right)|xy|^{\alpha _3}|yz|^{\alpha _1}|zx|^{\alpha _2}}}`$ (104)
$`\times (\delta ^{AC}\delta ^{BD}\delta ^{AD}\delta ^{BC})C(I_1,I_2,I_3).`$
The normalizations of the fields can be fixed by demanding that the two point functions are
$$A^{(I_1)}(x)A^{(I_2)}(y)=\frac{\delta ^{I_1I_2}}{|xy|^{2l}}.$$
(105)
for any two chiral primary fields $`A^I`$. After rescaling the fields in order to satisfy this condition, the normalized three-point function are given by:
$`T^{(1)}(x)T^{(2)}(y)T^{(3)}(z)`$ $`=`$ $`{\displaystyle \frac{2^{7/2}}{\pi ^{1/4}}}{\displaystyle \frac{\underset{i}{}\mathrm{\Gamma }(\frac{\alpha _i}{2}+1)\mathrm{\Gamma }(\frac{1}{2}\mathrm{\Sigma }+\frac{3}{2})}{\sqrt{_i^3\left(\mathrm{\Gamma }(l_i+3/2)\mathrm{\Gamma }(l_i+2)(l_i+1)(l_i^21)\right)}}}`$ (107)
$`\times {\displaystyle \frac{(l_1(l_1^21)+l_2(l_2^21)+l_3(l_3^21))}{|xy|^{\alpha _3}|yz|^{\alpha _1}|zx|^{\alpha _2}}}C(I_1,I_2,I_3)`$
$`S^{(1)}S^{(2)}T^{(3)}`$ $`=`$ $`{\displaystyle \frac{2^{5/2}}{\pi ^{1/4}}}{\displaystyle \frac{\underset{i}{}\mathrm{\Gamma }(\frac{\alpha _i}{2}+1)\mathrm{\Gamma }(\frac{1}{2}\mathrm{\Sigma }+\frac{3}{2})}{\sqrt{_i^3\left(\mathrm{\Gamma }(l_i+3/2)\mathrm{\Gamma }(l_i+1)(l_i^21)\right)(l_3+1)}}}`$ (110)
$`\times {\displaystyle \frac{(l_1(l_1^21)+l_2(l_2^21)+l_3(l_3^21))}{|xy|^{\alpha _3}|yz|^{\alpha _1}|zx|^{\alpha _2}}}C(I_1,I_2,I_3)`$
$`V^{(1)AB}U^{(2)CD}S^{(3)}`$ $`=`$ $`{\displaystyle \frac{2^2}{\pi ^{1/4}}}{\displaystyle \frac{\underset{i}{}\mathrm{\Gamma }(\frac{\alpha _i}{2}+1)\mathrm{\Gamma }(\frac{1}{2}\mathrm{\Sigma }+\frac{3}{2})(l_1l_2)C(I_1,I_2,I_3)(\delta ^{AC}\delta ^{BD}\delta ^{AD}\delta ^{BC})}{\sqrt{_i^3\left(\mathrm{\Gamma }(l_i+3/2)\mathrm{\Gamma }(l_i+1)\right)(l_3^21)}|xy|^{\alpha _3}|yz|^{\alpha _1}|zx|^{\alpha _2}}}`$ (112)
$`V_{AB}^{(1)}V_{CD}^{(2)}V_{EF}^{(3)}`$ $`=`$ $`{\displaystyle \frac{2^2}{\pi ^{1/4}}}{\displaystyle \frac{\underset{i}{}\mathrm{\Gamma }(\frac{\alpha _i}{2}+1)\mathrm{\Gamma }(\frac{1}{2}\mathrm{\Sigma }+\frac{3}{2})C(I_1,I_2,I_3)ϵ_{12ABCDEF}}{\sqrt{_i^3\left(\mathrm{\Gamma }(l_i+3/2)\mathrm{\Gamma }(l_i+1)\right)}|xy|^{\alpha _3}|yz|^{\alpha _1}|zx|^{\alpha _2}}}`$ (114)
$`V_{AB}^{(1)}V_{CD}^{(2)}T^{(3)}`$ $`=`$ $`{\displaystyle \frac{2^{7/2}}{\pi ^{1/4}}}{\displaystyle \frac{(l_1+l_2)\underset{i}{}\mathrm{\Gamma }(\frac{\alpha _i}{2}+1)\mathrm{\Gamma }(\frac{1}{2}\mathrm{\Sigma }+\frac{3}{2})}{(l_3+1)\sqrt{_i^3\left(\mathrm{\Gamma }(l_i+3/2)\mathrm{\Gamma }(l_i+1)\right)(l_3^21)}}}`$ (117)
$`\times {\displaystyle \frac{(\delta _{AC}\delta _{BD}\delta _{AD}\delta _{BC})C(I_1,I_2,I_3)}{|xy|^{\alpha _3}|yz|^{\alpha _1}|zx|^{\alpha _2}}}`$
$`U_{AB}^{(1)}U_{CD}^{(2)}V_{EF}^{(3)}`$ $`=`$ $`{\displaystyle \frac{2^2}{\pi ^{1/4}}}{\displaystyle \frac{\underset{i}{}\mathrm{\Gamma }(\frac{\alpha _i}{2}+1)\mathrm{\Gamma }(\frac{1}{2}\mathrm{\Sigma }+\frac{3}{2})C(I_1,I_2,I_3)ϵ_{12ABCDEF}}{\sqrt{\left(_i^3\mathrm{\Gamma }(l_i+3/2)\mathrm{\Gamma }(l_i+1)\right)}|xy|^{\alpha _3}|yz|^{\alpha _1}|zx|^{\alpha _2}}}`$ (119)
$`U^{(1)AB}U^{(2)CD}T^{(3)}`$ $`=`$ $`{\displaystyle \frac{2^{7/2}}{\pi ^{1/4}}}{\displaystyle \frac{(l_1+l_2)\underset{i}{}\mathrm{\Gamma }(\frac{\alpha _i}{2}+1)\mathrm{\Gamma }(\frac{1}{2}\mathrm{\Sigma }+\frac{3}{2})}{(l_3+1)\sqrt{\left(_i^3\mathrm{\Gamma }(l_i+3/2)\mathrm{\Gamma }(l_i+1)\right)(l_3^21)}}}`$ (122)
$`\times {\displaystyle \frac{(\delta ^{AC}\delta ^{BD}\delta ^{AD}\delta ^{BC})C(I_1,I_2,I_3)}{|xy|^{\alpha _3}|yz|^{\alpha _1}|zx|^{\alpha _2}}}.`$
## V Conclusions
In this paper, we gave the calculation of three point interactions of chiral primaries in 4D supergravity. In order to remove the derivatives from the interactions, we used nontrivial redefinitions of the fields which solved the linear equation of motion for chiral primaries. It is after removing the derivative terms that we deal with a standard field theory in two dimensions. These redefinitions were also needed for the cases of higher dimensional $`AdS`$ spaces,. Deeper reasons behind these redefinitions are still not clear. <sup>§</sup><sup>§</sup>§See also ref. for related discussions for the case of the $`D=11`$ supergravity on $`AdS_7\times S^4`$.
In section 4, we derived the three point interactions. The factors appear which are similar to the results in higher dimensional $`AdS`$ spacetimes. In order to make a useful statement on $`AdS_2/CFT_1`$ correspondence, similar computation has to be made in the boundary conformal quantum mechanics dual to this theory. There has been a conjecture that this dual quantum mechanics is given by a supersymmetric extension of the Calogero model,,. It would be interesting to see whether this is true. However, since a many-body quantum mechanics becomes a $`1+1`$ dimensional field theory after the second quantization, one might directly obtain this field theory starting from the Lagrangian (84) by appropriate truncation. This interesting issue, along with the evluation of the interaction terms for the fermions, are left for future studies.
Acknowledgements: I give special thanks to Antal Jevicki for many helpful discussions and suggestions. I also thank Tamiaki Yoneya and Sangmin Lee for useful discussions. This work was supported by JSPS through Institute of Physics, University of Tokyo.
Appendix 1 : Vector and scalar part of the lagrangian As mentioned in the text, the $`SU(8)`$ gauge symmetry can be fixed to the so called symmetric gauge where $`𝒱=\mathrm{exp}(X)`$ and
$$X=\left(\begin{array}{cc}0& W_{ABCD}\\ \overline{W}^{ABCD}& 0\end{array}\right)$$
(123)
Expanding in $`W`$, we get the expressions
$$𝒱=\left(\begin{array}{cc}\delta _{AB}^{}{}_{}{}^{CD}+\frac{1}{2}W_{ABEF}\overline{W}^{EFCD}+O(W^4)& W_{ABCD}+O(W^3)\\ \overline{W}^{ABCD}+O(W^3)& \delta _{}^{AB}{}_{CD}{}^{}+\frac{1}{2}\overline{W}^{ABEF}W_{EFCD}+O(W^4)\end{array}\right),$$
(124)
and
$$𝒱^1=\left(\begin{array}{cc}\delta _{AB}^{}{}_{}{}^{CD}+\frac{1}{2}W_{ABEF}\overline{W}^{EFCD}+O(W^4)& W_{ABCD}+O(W^3)\\ \overline{W}^{ABCD}+O(W^3)& \delta _{}^{AB}{}_{CD}{}^{}+\frac{1}{2}\overline{W}^{ABEF}W_{EFCD}+O(W^4)\end{array}\right).$$
(125)
From these expressions one easily gets
$$_{\widehat{\mu }}𝒱𝒱^1=\left(\begin{array}{cc}\frac{1}{2}_\mu (W\overline{W})(_\mu W)\overline{W}+O(W^4)& _\mu W+O(W^3)\\ _\mu \overline{W}+O(W^3)& \frac{1}{2}_\mu (\overline{W}W)(_\mu \overline{W})W+O(W^4)\end{array}\right)$$
(126)
Comparing with Eq.(1), we see that
$$P_{\widehat{\mu }ABCD}=_{\widehat{\mu }}W_{ABCD}+O(W^3),\overline{P}^{\widehat{\mu }ABCD}=_{\widehat{\mu }}W^{ABCD}+O(W^3).$$
(127)
$`\stackrel{~}{H}(B,𝒱)`$ is defined by the equation
$$\left(\begin{array}{c}G_{\widehat{\mu }\widehat{\nu }}+iH_{\widehat{\mu }\widehat{\nu }}\\ G_{\widehat{\mu }\widehat{\nu }}iH_{\widehat{\mu }\widehat{\nu }}\end{array}\right)=(𝒱^{}𝒱)^1\left(\begin{array}{c}i\stackrel{~}{G}_{\widehat{\mu }\widehat{\nu }}\stackrel{~}{H}_{\widehat{\mu }\widehat{\nu }}\\ i\stackrel{~}{G}_{\widehat{\mu }\widehat{\nu }}\stackrel{~}{H}_{\widehat{\mu }\widehat{\nu }}\end{array}\right),$$
(128)
where it is to be understood that the matrices are multiplied by contracting the $`SU(8)`$ indices, which I did not write explicitly. In this paper, I use only $`SU(8)`$ indices, in contrast to ref. where $`E_7`$ indices were also used. $`\stackrel{~}{G},\stackrel{~}{H}`$ denotes the dual fields,
$$\stackrel{~}{G}^{\widehat{\mu }\widehat{\nu }}=\frac{1}{2}ϵ^{\widehat{\mu }\widehat{\nu }\widehat{\rho }\widehat{\sigma }}G_{\widehat{\rho }\widehat{\sigma }},\stackrel{~}{H}^{\widehat{\mu }\widehat{\nu }}=\frac{1}{2}ϵ^{\widehat{\mu }\widehat{\nu }\widehat{\rho }\widehat{\sigma }}H_{\widehat{\rho }\widehat{\sigma }}$$
(129)
Expanding in $`W`$, we get
$$\left(\begin{array}{c}G_{\widehat{\mu }\widehat{\nu }}+iH_{\widehat{\mu }\widehat{\nu }}\\ G_{\widehat{\mu }\widehat{\nu }}iH_{\widehat{\mu }\widehat{\nu }}\end{array}\right)=\left(\begin{array}{cc}1+2W\overline{W}+O(W^4)& 2W\frac{4}{3}W\overline{W}W+O(W^5)\\ 2\overline{W}\frac{4}{3}\overline{W}W\overline{W}+O(W^5)& 1+2\overline{W}W+O(W^4)\end{array}\right)\left(\begin{array}{c}i\stackrel{~}{G}_{\widehat{\mu }\widehat{\nu }}\stackrel{~}{H}_{\widehat{\mu }\widehat{\nu }}\\ i\stackrel{~}{G}_{\widehat{\mu }\widehat{\nu }}\stackrel{~}{H}_{\widehat{\mu }\widehat{\nu }}\end{array}\right).$$
(130)
By adding the first and second row, we get
$`G`$ $`=`$ $`(1W\overline{W}+W\overline{W}+\overline{W}W{\displaystyle \frac{2}{3}}\overline{W}W\overline{W}{\displaystyle \frac{2}{3}}W\overline{W}W+O(W^5))\stackrel{~}{H}`$ (132)
$`+i(\overline{W}+W+W\overline{W}\overline{W}W{\displaystyle \frac{2}{3}}\overline{W}W\overline{W}+{\displaystyle \frac{2}{3}}W\overline{W}W)\stackrel{~}{G}.`$
Solving in terms of $`H`$, we finally get
$`G_{\widehat{\mu }\widehat{\nu }}^{AB}\stackrel{~}{H}_{\widehat{\mu }\widehat{\nu }AB}^{(B)}`$ $`=`$ $`G_{\widehat{\mu }\widehat{\nu }}^{AB}(1+W+\overline{W}+W^2+\overline{W}^2`$ (135)
$`{\displaystyle \frac{1}{3}}W\overline{W}W{\displaystyle \frac{1}{3}}\overline{W}W\overline{W}+W^3+\overline{W}^3)_{ABCD}G^{\widehat{\mu }\widehat{\nu }CD}`$
$`+iG_{\widehat{\mu }\widehat{\nu }}^{AB}(W\overline{W}+W^2\overline{W}^2+O(W^3))_{ABCD}\stackrel{~}{G}^{\widehat{\mu }\widehat{\nu }CD}+O(W^4),`$
. (136)
Appendix 2 : Spherical Harmonics
We list some formulas about spherical harmonics needed for the calculations in the text. The spherical hamonics on $`S^2`$ are very simple in that there are only scalar spherical harmonics, which we denote by $`Y_I=Y_{lm}`$. They are normalized so that
$$Y_{l_1m_1}Y_{l_2m_2}=\delta _{l_1l_2}\delta _{m_1m_2}.$$
(137)
The explicit form is given by
$$Y_{lm}\frac{1}{2}\sqrt{\frac{(2l+1)(lm)}{\pi (l+m)}}e^{im\varphi }P_l^m(cos\theta )$$
(138)
where $`P_l^m(x)`$ is the associated Legendre polynomial
$$P_l^m(x)\frac{(1)^m}{2^ll!}(1x^2)^{m/2}\frac{d^{m+l}}{dx^{m+l}}(x^21)^l,$$
(139)
although this explicit form was not really used in the text. More important is the formula for the integral of three spherical harmonics with derivatives, expressed in terms of $`C(I_1,I_2,I_3)`$. They are as follows: again, we use the abbreviation $`i`$ for $`I_i(l_i,m_i)`$
$`A(1;2,3)`$ $``$ $`{\displaystyle Y_1Y_2Y_3}`$ (141)
$`=`$ $`{\displaystyle \frac{1}{2}}(l_2(l_2+1)+l_3(l_3+1)l_1(l_1+1))C(1,2,3)`$ (142)
$`B(1;2,3)`$ $``$ $`{\displaystyle Y_1_\alpha _\beta Y_2^\alpha ^\beta Y_3}`$ (143)
$`=`$ $`{\displaystyle \frac{1}{4}}(l_1(l_1+1)l_2(l_2+1)l_3(l_3+1))`$ (145)
$`\times (2+l_1(l_1+1)l_2(l_2+1)l_3(l_3+1))C(1,2,3)`$
$`D(1;2,3)`$ $``$ $`{\displaystyle _\alpha _\beta Y_1^\alpha Y_2^\beta Y_3}`$ (146)
$`=`$ $`{\displaystyle \frac{1}{4}}(l_3(l_3+1)+l_1(l_1+1)l_2(l_2+1))`$ (148)
$`\times (l_3(l_3+1)+l_2(l_2+1)l_1(l_1+1))C(1,2,3)`$
$`G(1,2,3)`$ $``$ $`{\displaystyle _\alpha _\beta Y_1_\beta _\gamma Y_2^\gamma ^\alpha Y_3}`$ (149)
$`=`$ $`{\displaystyle \frac{1}{8}}(l_1^3(l_1+1)^3l_2^3(l_2+1)^3l_3^3(l_3+1)^3`$ (155)
$`2l_1^2(1_1+1)^22l_2^2(1_2+1)^22l_3^2(1_3+1)^2`$
$`+4l_1l_2(l_1+1)(l_2+1)+4l_2l_3(l_2+1)(l_3+1)+4l_3l_1(l_3+1)(l_1+1)`$
$`+l_1^2l_2(l_1+1)^2(l_2+1)+l_1l_2^2(l_1+1)(l_2+1)^2+l_2^2l_3(l_2+1)^2(l_3+1)`$
$`+l_2l_3^2(l_2+1)(l_3+1)^2+l_3^2l_1(l_3+1)^2(l_1+1)+l_3l_1^2(l_3+1)(l_1+1)^2`$
$`2l_1l_2l_3(l_1+1)(l_2+1)(l_3+1))`$
proof)
a) By integrating by parts, one gets
$`A(1;2,3)`$ $`=`$ $`{\displaystyle Y_1^2Y_2Y_3}{\displaystyle _\alpha Y_1^\alpha Y_2Y_3}`$ (156)
$`=`$ $`l_2(l_2+1)C(1,2,3)A(3;1,2).`$ (157)
By permuting the indices and adding the resulting equations, we get Eq.(142).
b) By integrating by parts, we get
$`B(1;2,3)`$ $`=`$ $`{\displaystyle Y_1_\alpha _\beta Y_2^\alpha ^\beta Y_3}`$ (158)
$`=`$ $`{\displaystyle _\alpha Y_1^\alpha ^\beta Y_2_\beta Y_3}{\displaystyle Y_1_\alpha _\beta ^\alpha Y_2^\beta Y_3}`$ (159)
$`=`$ $`{\displaystyle _\alpha Y_1^\alpha Y_2_y^2Y_3}+_\beta _\alpha Y_1^\alpha Y_2^\beta Y_3`$ (161)
$`Y_1^\gamma Y_2^\beta Y_3R_{\beta \gamma }Y_1^\beta _y^2Y_2_\beta Y_3`$
$`=`$ $`l_3(l_3+1)_\alpha Y_1^\alpha Y_2Y_3_\beta _\alpha ^\beta Y_1^\alpha Y_2Y_3_\beta _\alpha Y_1^\beta ^\alpha Y_2Y_3`$ (163)
$`Y_1^\beta Y_2_\beta Y_3+l_2(l_2+1)Y_1^\beta Y_2_\beta Y_3`$
$`=`$ $`l_3(l_3+1)A(3;1,2)R_{\alpha \gamma }^\alpha Y_1^\gamma Y_2Y_3`$ (166)
$`_\alpha _y^2Y_1^\alpha Y_2Y_3B(3;1,2)`$
$`+A(1;2,3)(l_2(l_2+1)1)`$
$`=`$ $`(l_1(l_1+1)l_3(l_3+1)1)A(3;1,2)B(3;1,2)+A(1;2,3)(l_2(l_2+1)1)`$ (167)
which gives
$$B(1;2,3)+B(3;1,2)=(l_2(l_2+1)1)A(1;2,3)+(l_1(l_1+1)l_3(l_3+1)1)A(3;1,2)$$
(168)
On the other hand, we have
$`B(1;2,3)`$ $`=`$ $`{\displaystyle ^\alpha Y_1_\alpha _\beta Y_2^\beta Y_3}{\displaystyle Y_1^\alpha _\beta _\alpha Y_2^\beta Y_3}`$ (169)
$`=`$ $`{\displaystyle ^\alpha ^\beta Y_1_\alpha _\beta Y_2Y_3}+{\displaystyle ^\alpha Y_1^\beta _\alpha _\beta Y_2Y_3}`$ (171)
$`R_{\beta \gamma }{\displaystyle Y_1^\gamma Y_2^\beta Y_3}Y_1^\beta _y^2Y_2_\beta Y_3`$
$`=`$ $`B(3;1,2)+R_{\alpha \gamma }^\alpha Y_1^\gamma Y_2Y_3`$ (173)
$`+_\alpha Y_1^\alpha _y^2Y_2Y_3Y_1^\beta Y_2_\beta Y_3+l_2(l_2+1)Y_1_\beta Y_2^\beta Y_3`$
$`=`$ $`B(3;1,2)+^\alpha Y_1_\alpha Y_2Y_3`$ (175)
$`l_2(l_2+1)_\alpha Y_1^\alpha Y_2Y_3Y_1^\beta Y_2_\beta Y_3+l_2(l_2+1)Y_1_\beta Y_2^\beta Y_3`$
$`=`$ $`B(3;1,2)+(1l_2(l_2+1))(A(3;1,2)A(1;2,3))`$ (176)
which is
$$B(1;2,3)B(3;1,2)=(1l_2(l_2+1))(A(3;1,2)A(1;2,3))$$
(177)
Adding (168) and (177) gives Eq.(145).
c)
$`D(1;2,3)`$ $`=`$ $`{\displaystyle _\alpha _\beta Y_1^\alpha Y_2^\beta Y_3}`$ (178)
$`=`$ $`{\displaystyle _\beta Y_1_y^2Y_2^\beta Y_3}^\beta Y_1^\alpha Y_2_\alpha _\beta Y_3`$ (179)
$`=`$ $`l_2(l_2+1){\displaystyle _\beta Y_1Y_2^\beta Y_3}+Y_1_\beta _\alpha Y_2^\beta ^\alpha Y_3+Y_1_\alpha Y_2^\beta ^\alpha _\beta Y_3`$ (180)
$`=`$ $`l_2(l_2+1){\displaystyle _\beta Y_1Y_2^\beta Y_3}+Y_1_\beta _\alpha Y_2^\beta ^\alpha Y_3`$ (182)
$`+R_{\alpha \gamma }Y_1^\alpha Y_2^\gamma Y_3+Y_1_\alpha Y_2^\alpha _y^2Y_3`$
$`=`$ $`l_2(l_2+1)A(2;3,1)+B(1;2,3)+(1l_3(l_3+1))A(1;2,3)`$ (183)
which is equivalent to Eq.(148).
d)
$`G(1,2,3)`$ $`=`$ $`{\displaystyle _\alpha _\beta Y_1_\beta _\gamma Y_2^\gamma ^\alpha Y_3}`$ (184)
$`=`$ $`{\displaystyle _\gamma _\alpha Y_1^\alpha ^\beta ^\gamma Y_2_\beta Y_3}(1l_1(l_1+1)){\displaystyle _\gamma Y_1^\beta ^\gamma Y_2_\beta Y_3}`$ (185)
$`=`$ $`{\displaystyle _\gamma _\alpha Y_1(R_{\alpha }^{}{}_{}{}^{\beta }{}_{\gamma }{}^{}{}_{}{}^{\delta }_\delta Y_2+^\beta _\alpha ^\gamma Y_2)_\beta Y_3}+(l_1(l_1+1)1){\displaystyle _\gamma Y_1^\beta ^\gamma Y_2_\beta Y_3}`$ (186)
$`=`$ $`(l_1(l_1+1)1)D(2;3,1){\displaystyle _y^2Y_1_\beta Y_2^\beta Y_3}`$ (188)
$`+{\displaystyle ^\beta ^\alpha Y_1_\alpha Y_2_\beta Y_3}{\displaystyle _\gamma _\alpha Y_1^\beta _\alpha ^\gamma Y_2_\beta Y_3}`$
$`=`$ $`(l_1(l_1+1)1)D(2;3,1)+l_1(l_1+1)A(1;2,3)+D(1;2,3)l_3(l_3+1)B(3;1,2)`$ (190)
$`+{\displaystyle ^\beta _\gamma _\alpha Y_1^\alpha ^\gamma Y_2_\beta Y_3}`$
$`=`$ $`(l_1(l_1+1)1)D(2;3,1)+l_1(l_1+1)A(1;2,3)+D(1;2,3)l_3(l_3+1)B(3;1,2)`$ (192)
$`+{\displaystyle (R_{}^{\beta }{}_{\alpha \gamma }{}^{}{}_{}{}^{\delta }_\delta Y_1+_\alpha ^\beta _\gamma Y_1)^\alpha ^\gamma Y_2_\beta Y_3}`$
$`=`$ $`l_1(l_1+1)D(2;3,1)+l_1(l_1+1)A(1;2,3)l_3(l_3+1)B(3;1,2)`$ (194)
$`+l_2(l_2+1)A(2;3,1)+l_2(l_2+1)D(1;2,3)`$
which is Eq.(155).
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# Search for Quark-Lepton Compositeness at Tevatron and LHC
## Abstract
We make a Monte Carlo study on compositeness of first generation quarks and leptons using the Drell-Yan distribution in the high dielectron mass region at the Tevatron and LHC energies. The current experimental lower limits on the compositeness scale, $`\mathrm{\Lambda }`$, vary from 2.5 to 6.1 TeV. In the present analysis, we assume that there will be no deviation of the dielectron mass spectrum from Standard Model prediction at center of mass energy 2 TeV (Tevatron) and 14 TeV (LHC). We then find that in the LL, RR, RL and LR chirality channels of the quark-electron currents, it is possible to extend the lower limits on $`\mathrm{\Lambda }`$ (at 95$`\%`$ CL) to a range of 6 to 10 TeV for $`2fb^1`$ and 9 to 19 TeV for $`30fb^1`$ of integrated luminosity at Tevatron. At LHC, the corresponding limits extend to a range of 16 to 25 TeV for $`10fb^1`$ and 20 to 36 TeV for $`100fb^1`$ of integrated luminosity.
PACS numbers: 12.60.Rc, 12.60.-i, 13.85.-t
The proliferation of quarks and leptons has inspired the speculation that they could be composite structures, i.e. bound states of more fundamental constituents often called preons . Below a characteristic energy scale called the compositeness scale, $`\mathrm{\Lambda }`$, the preon-binding interaction becomes strong and binds the constituents to form composite states like the quarks and leptons. With such a composite structure, there would be significant deviation from the Standard Model (SM) prediction of high energy cross sections. No such deviation has been observed so far. These null results have been used to put lower limits on quark-lepton compositeness scale $`\mathrm{\Lambda }`$. which varies from 2.5 to 6.1 TeV in the various chirality channels of the quark-lepton currents.
In this paper, we consider the effects of composite structure of first generation quarks and leptons on the Drell-Yan (DY) process $`q\overline{q}e^+e^{}`$ . If the compositeness scale, $`\mathrm{\Lambda }`$, is much greater than $`\sqrt{\widehat{s}}`$, the center of mass energy of the colliding partons, the quarks and electrons would appear to be point-like. The substructure coupling can then be approximated by a four-fermion contact interaction giving rise to the following effective lagrangian<sup>*</sup><sup>*</sup>*Here we have assumed that the contact interaction is color singlet and weak-isoscalar. :
$`_{ql}`$ $`=`$ $`{\displaystyle \frac{g_0^2}{\mathrm{\Lambda }^2}}\{\eta _{LL}(\overline{q}_L\gamma ^\mu q_L)(\overline{e}_L\gamma _\mu e_L)+\eta _{LR}(\overline{q}_L\gamma ^\mu q_L)(\overline{e}_R\gamma _\mu e_R)`$ (3)
$`+\eta _{RL}(\overline{u}_R\gamma _\mu u_R)(\overline{e}_L\gamma ^\mu e_L)+\eta _{RL}(\overline{d}_R\gamma _\mu d_R)(\overline{e}_L\gamma ^\mu e_L)`$
$`+\eta _{RR}(\overline{u}_R\gamma ^\mu u_R)(\overline{e}_R\gamma _\mu e_R)+\eta _{RR}(\overline{d}_R\gamma ^\mu d_R)(\overline{e}_R\gamma _\mu e_R)\}`$
where
$`q_L=\left[\begin{array}{c}u\\ d\end{array}\right]_L`$ (6)
is the left-handed quark doublet. $`u_R`$ and $`d_R`$ are the right-handed quark singlets. $`e_L`$ and $`e_R`$ are the left- and right-handed electrons respectively. The compositeness scale $`(\mathrm{\Lambda })`$ is chosen so that $`\frac{g_0^2}{4\pi }=1`$ and the largest $`\eta _{ij}=1`$, where $`g_0`$ is the coupling constant for the contact interaction and $`\eta _{ij}`$ is the interference term between the contact interaction and the SM lagrangian for the $`ij^{th}`$ channel, with $`i`$ and $`j`$ representing the helicities of the quark and the lepton currents. Including the above contact interaction (at $`\mathrm{\Lambda }>>\sqrt{\widehat{s}}`$), the DY cross section gets transformed as :
$`{\displaystyle \frac{d\sigma ^\mathrm{\Lambda }}{dm}}={\displaystyle \frac{d\sigma }{dm}}(DY)+\beta I+\beta ^2C,`$ (7)
where $`\beta =1/\mathrm{\Lambda }^2`$ and m is the dielectron invariant mass. In this expression, $`I`$ is due to the interference of DY and the contact term, and $`C`$ is the pure contact term contribution to the cross-section. The deviation in the dielectron production from SM expectations would be dominant in the high mass region above the Z pole. We have made separate studies for quark- electron compositeness for an integrated luminosity of 2 $`fb^1`$ (Run II) and 30 $`fb^1`$ (TEV33) with respect to the D$`Ø`$ detector at Tevatron and an integrated luminosity of 10 $`fb^1`$ and 100 $`fb^1`$ with respect to the CMS detector at LHC. However the results should be valid for the CDF detector at Tevatron and the ATLAS detector at LHC as well. We have simulated dielectron production through DY process alone in p$`\overline{p}`$ (pp) collisions at center of mass energy, $`\sqrt{s}`$, equal to 2 TeV (14 TeV) using PYTHIA . However since PYTHIA does not incorporate all the compositeness models, we have used a separate parton level Monte Carlo program to estimate dielectron production rates in the presence of compositeness. Assuming that the Tevatron and LHC data on dielectron production are consistent with DY predictions under SM, we extract limits on compositeness scale using Bayesian technique of statistical inference . We have considered four different models corresponding to the LL, RR, RL and LR chirality channels of equation 3 for quark-electron compositeness. The choice of $`\eta _{ij}`$ for the different models of compositeness is listed in Table I.
Exploring the lower limits on $`\mathrm{\Lambda }`$ at Tevatron
We simulate p$`\overline{p}`$ collisions using PYTHIA at 2 TeV and generate DY dielectron events between 95 GeV and 1.5 TeV of the dielectron invariant mass. The total number of dielectron events generated by PYTHIA, $`N_{gen}`$, gives the expected number of background subtracted dielectron events, $`N_{DY}`$, to be collected at Tevatron as :
$`N_{DY}=ϵ\times N_{gen}`$ (8)
where $`ϵ`$ is the detection efficiency of the dielectron. The detection efficiency, $`ϵ`$, of the dielectron involves contribution from the following terms:
* Energy smearing,
* Electron identification efficiency, $`ϵ_1`$, and
* Acceptance, $`ϵ_2`$.
The energy resolution of the electromagnetic calorimeter of the upgraded D$`Ø`$ detector is parameterized as :
$`({\displaystyle \frac{\sigma }{E}})^2=C^2+({\displaystyle \frac{a}{\sqrt{E}}})^2(EinGeV)`$ (9)
where the constant term, C, and the stochastic term, $`a`$ are taken to 2 $`\%`$ and 16 $`\%`$ respectively . We take the electron identification efficiency, $`ϵ_1`$, for a single electron to be 85$`\%`$. The identification efficiency for a dielectron is then $`ϵ_1^2`$. The acceptance, $`ϵ_2`$, of dielectron events in p$`\overline{p}`$ collisions is defined as the fraction of events in which the $`e^+e^{}`$ pair passes the fiducial and the kinematic cuts after taking into account the energy smearing. The fiducial and the kinematic cuts used are:
* $`\eta 2.5`$, where $`\eta `$ is the pseudorapidity ($`=ln[tan(\frac{\theta }{2})]`$). This ensures that the dielectron event selected is in the active detector region.
* A kinematic cut of $`p_T25GeV`$, where $`p_T`$ is the transverse momentum of the electron and the positron. This cut ensures an efficient trigger.This cut is based on the D$`Ø`$ Run I analysis of DY data at 1.8 TeV .
The dielectron detection efficiency, $`ϵ`$, is then :
$`ϵ=ϵ_{1}^{}{}_{}{}^{2}\times ϵ_2`$ (10)
We then generate the expected number of dielectron events, $`N_{exp}^\mathrm{\Lambda }`$, in various mass bins including the effect of the composite structure of quarks and electrons for various values of $`\mathrm{\Lambda }`$ using the parton level Monte Carlo. We calculate the cross section ($`\sigma ^\mathrm{\Lambda }`$) for the production of dielectrons including terms from the contact interaction lagrangian of equation 3 with the SM lagrangian. The LO cross section calculation is corrected for higher order QCD effects using a K-factor of 1.22This K-factor is the ratio of the NNLO DY cross section to the LO DY cross section at 1.8 TeV . We consider the same value for the K-factor for DY + compositeness at 2 TeV.. We checked the parton level MC calculation by comparing its
prediction with that from PYTHIA for the Drell-Yan process. Both calculations agree to within a few percent as shown in Fig. 1.
In order to obtain the lower limit on $`\mathrm{\Lambda }`$, we then use the Bayesian technique to compare the Drell Yan dielectron mass distribution (ie., $`N_{DY}`$) in the high mass region with the expected dielectron mass distribution for various values of $`\mathrm{\Lambda }`$ (ie., $`N_{exp}^\mathrm{\Lambda }`$). Limits are obtained independently for each separate channel of the contact interaction lagrangian: LL, RR, RL and LR with $`\eta _{ij}=\pm 1`$. Fig. 2 shows the cross section versus the dielectron invariant mass, in the high mass region between 50 GeV and 1.8 TeV in the LL channel for different values of $`\mathrm{\Lambda }`$ for $`\eta _{ij}=1`$ (constructive interference) and Fig. 3 shows the corresponding plot for $`\eta _{ij}=+1`$ (destructive interference).
Since the effect of compositeness is most pronounced in the high dielectron mass region we consider 10 different mass bins of variable width between 120 GeV and 1.5 TeV. The expected number of events at the compositeness scale, $`\mathrm{\Lambda }`$, in the $`k^{th}`$ mass bin is given as :
$`N_{exp}^{\mathrm{\Lambda },k}=ϵ^k(\sigma ^{\mathrm{\Lambda },k}\times L),`$ (11)
where $`\sigma ^{\mathrm{\Lambda },k}`$ is the cross section (including compositeness) for the $`k^{th}`$ mass bin and $`L`$ is the integrated luminosity. The posterior probability for the compositeness scale to be $`\mathrm{\Lambda }`$ given the expected DY dielectron data distribution, $`d_O`$, is:
$`P(\Lambda d_O)={\displaystyle \frac{1}{𝒵}}{\displaystyle \underset{k=1}{\overset{n}{}}}P^k(N_{DY}^kN_{exp}^{\mathrm{\Lambda },k})P(ϵ^k,L,\mathrm{\Lambda })`$ (12)
where $`𝒵`$ is the normalization constant. $`P^k(N_{DY}^kN_{exp}^{\mathrm{\Lambda },k})`$ is the likelihood function which follows a Poisson distribution for small $`N_{exp}^{\mathrm{\Lambda },k}`$:
$`P^k(N_{DY}^kN_{exp}^{\mathrm{\Lambda },k})={\displaystyle \frac{e^{N_{exp}^{\mathrm{\Lambda },k}}(N_{exp}^{\mathrm{\Lambda },k})^{N_{DY}^k}}{N_{DY}^k!}},(N_{exp}^{\mathrm{\Lambda },k}<10)`$ (13)
and a Gaussian distribution for large $`N_{exp}^{\mathrm{\Lambda },k}`$, with mean $`N_{exp}^{\mathrm{\Lambda },k}`$ and standard deviation, $`\sigma _1`$, ($`\sigma _1=\sqrt{N_{exp}^{\mathrm{\Lambda },k}}`$) :
$`P(N_{DY}^kN_{exp}^{\mathrm{\Lambda },k})={\displaystyle \frac{1}{\sqrt{2\pi }\sigma _1}}e^{\frac{(N_{DY}^kN_{exp}^{\mathrm{\Lambda },k})^2}{2\sigma _1^2}},(N_{exp}^{\mathrm{\Lambda },k}10).`$ (14)
$`P(ϵ^k,L,\mathrm{\Lambda })`$ is the joint $`prior`$ probability for the dielectron detection efficiency, $`ϵ^k`$, the integrated luminosity, $`L`$, and the compositeness scale, $`\mathrm{\Lambda }`$. Taking $`ϵ^k`$, $`L`$ and $`\mathrm{\Lambda }`$ to be independent,
$`P(ϵ^k,L,\mathrm{\Lambda })=P(ϵ^k)P(L)P(\mathrm{\Lambda }).`$ (15)
The $`\mathrm{𝑝𝑟𝑖𝑜𝑟}`$ probabilities of detection efficiency, $`ϵ^k`$, and integrated luminosity, $`L`$, are assumed to be Gaussian with their estimated value in each bin as the $`mean`$ and corresponding error as the $`width`$ of the Gaussian. The prior distribution $`P(\mathrm{\Lambda })`$ is chosen to be uniform in $`1/\mathrm{\Lambda }^2`$. This represents a prior essentially flat in cross section. The resulting posterior density $`P(\mathrm{\Lambda }d_O)`$ peaks at $`1/\mathrm{\Lambda }^2=0`$ and falls off monotonically with increasing $`1/\mathrm{\Lambda }^2`$. The $`95\%`$ CL lower limit on $`\mathrm{\Lambda }`$ is defined by:
$`{\displaystyle _{\mathrm{\Lambda }_{lim}}^{\mathrm{}}}𝑑\mathrm{\Lambda }^{}P(\mathrm{\Lambda }^{}d_O)=0.95.`$ (16)
The values of efficiency, $`ϵ^k`$, and the expected number of DY events, $`N_{DY}^k`$<sup>§</sup><sup>§</sup>§$`N_{DY}^k`$ is generated with a K-factor of 1.22 in PYTHIA., in individual mass bins are listed in Table II for an integrated luminosity of 2 $`fb^1`$ and 30 $`fb^1`$. The expected 95$`\%`$ CL lower limits on $`\mathrm{\Lambda }`$ for the LL, RR, RL and LR helicity channels of the quark- electron currents for both constructive and destructive interference are listed in Table III and Table IV for integrated luminosities of 2 $`fb^1`$ and 30 $`fb^1`$ respectively.
Exploring the lower limits on $`\mathrm{\Lambda }`$ at LHC
We have made a similar analysis of the DY process including the effect of quark-electron compositeness at 14 TeV. As before we have assumed that DY dielectron data that would be collected by the CMS detector at LHC would agree with SM prediction. We then use the Bayesian technique to obtain the lower limits on $`\mathrm{\Lambda }`$ at 14 TeV. We have made separate studies for $`10fb^1`$ of data and $`100fb^1`$ of data. A K-factor of 1.13 has been used as the NNLO correction factor. Fig. 4 shows the cross section versus the dielectron invariant mass, in the high mass region between 50 GeV and 2 TeV in the LL channel for different values of $`\mathrm{\Lambda }`$ for $`\eta _{ij}=1`$ (constructive interference) and Fig. 5 shows the corresponding plot for $`\eta _{ij}=+1`$ (destructive interference).
We generated DY events in the dielectron mass range of 150 GeV to 2 TeV. We then compared the expected number of DY events, $`N_{DY}`$, at $`\sqrt{s}`$ = 14 TeV with the expected number of dielectron events, $`N_{exp}^\mathrm{\Lambda }`$, at various values of $`\mathrm{\Lambda }`$ in the mass range of 500 GeV to 2 TeV where the deviation from SM predictions due to the composite structure of quarks and electrons is most pronounced at LHC. The electron identification efficiency, $`ϵ_1`$, is taken to be 95$`\%`$ . The constant and stochastic terms in the energy resolution of the electromagnetic calorimeter of the CMS detector are taken to be :
$`C`$ $`=`$ $`0.55\%,and`$ (17)
$`a`$ $`=`$ $`2.7\%,\eta 1.5`$ (19)
$`5.7\%,1.5<\eta 2.5`$
The fiducial and kinematic cuts selected are the same as for D$`Ø`$. The values of $`ϵ^k`$ and $`N_{DY}^k`$$`N_{DY}^k`$ is generated with a K-factor of 1.13 in PYTHIA., in individual mass bins are listed in Table V for integrated luminosities of 10 $`fb^1`$ and 100 $`fb^1`$.
The expected 95$`\%`$ CL lower limits on $`\mathrm{\Lambda }`$ for the LL, RR, RL and LR helicity channels of quark-electron currents for both constructive and destructive interference are listed in Table VI and Table VII for integrated luminosities of 10 $`fb^1`$ and 100 $`fb^1`$ respectively.
The discovery limits for $`\mathrm{\Lambda }`$ (defined as a deviation of 5$`\sigma `$ from SM prediction) for the various models have been listed for integrated luminosities of $`10fb^1`$, $`50fb^1`$, $`100fb^1`$, $`200fb^1`$ and $`500fb^1`$ in Table VIII for $`\eta _{ij}=1`$ and in Table IX for $`\eta _{ij}=+1`$.
Plots of the discovery limit versus the integrated luminosity for the various chirality channels are shown in Fig. 6 for $`\eta _{ij}=1`$ and in Fig. 7 for $`\eta _{ij}=+1`$.
To conclude, we have performed a Monte Carlo study of the dielectron invariant mass spectrum (DY + compositeness) for $`p\overline{p}`$ collisions at 2 TeV and $`pp`$ collisions at 14 TeV. We have considered the LL, RR, RL and LR chirality channels of the quark-electron currents. Assuming that Standard Model will describe the high mass DY dielectron data at 2 TeV and 14 TeV we have found that it is possible to extend the lower limits on the compositeness scale, $`\mathrm{\Lambda }`$, from the existing limits.
* For $`p\overline{p}`$ collisions at Tevatron we have made separate studies for integrated luminosities of $`2fb^1`$ and $`30fb^1`$ with respect to the D$`Ø`$ detector. The expected 95 $`\%`$ CL lower limits on $`\mathrm{\Lambda }`$ range between 6 to 10 TeV and 9 to 19 TeV for $`2fb^1`$ and $`30fb^1`$ of dielectron data, respectively. These limits are in agreement with similar limits on $`\mathrm{\Lambda }`$ quoted between 6 to 10 TeV for $`2fb^1`$ and 14 to 20 TeV for $`30fb^1`$ of data with respect to the CDF detector at Tevatron .
* For $`pp`$ collisions at LHC we have considered $`10fb^1`$ and $`100fb^1`$ of dielectron data with respect to the CMS detector. The expected 95 $`\%`$ CL lower limits on $`\mathrm{\Lambda }`$ range between 16 to 25 TeV for $`10fb^1`$ and between 20 to 36 TeV for $`100fb^1`$ of dielectron data.
* We have also explored the discovery potential for quark-electron compositeness (defined as a deviation of 5$`\sigma `$ from SM prediction) at LHC as a function of integrated luminosity.
The authors would like to thank Sreerup Raychaudhury and V.S.Narasimham for their advice, comments and stimulating questions. We would also like to thank D.P.Roy and Sudeshna Banerjee for several fruitful discussions.
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# On the hadron production from the quark–gluon plasma phase in ultra–relativistic heavy–ion collisions
## 1 Introduction
Recently we have suggested some kind of a coalescence model for the description of the hadronization from the quark–gluon plasma (QGP) phase of QCD considered as a thermalized quark–gluon system at high densities and temperature in which quarks, antiquarks and gluons being at the deconfined phase collide frequently each other. There is a belief that the QGP phase of the quark–gluon system can be realized in ultra–relativistic heavy–ion collision ($`E_{\mathrm{cms}}/\mathrm{nucleon}1\mathrm{GeV}`$) experiments.
At very high energies of heavy–ion collisions the quark–gluon system is composed from highly relativistic and very dense quarks, antiquarks and gluons. By virtue of the asymptotic freedom the particles are almost at liberty and due to high density collide themselves frequently that leads to an equilibrium state. If to consider such a state as a thermalized QGP phase of QCD, the probabilities of light massless quarks $`n_q(\stackrel{}{p})`$ and light massless antiquarks $`n_{\overline{q}}(\stackrel{}{p})`$, where $`q=u`$ or $`d`$, to have a momentum $`p`$ at a temperature $`T`$, can be described by the Fermi–Dirac distribution functions :
$`n_q(\stackrel{}{p})={\displaystyle \frac{1}{e^{\nu \left(T\right)+p/T}+1}},n_{\overline{q}}(\stackrel{}{p})={\displaystyle \frac{1}{e^{\nu \left(T\right)+p/T}+1}},`$ (1.1)
where a temperature $`T`$ is measured in $`\mathrm{MeV}`$, $`\nu (T)=\mu (T)/T`$, $`\mu (T)`$ is a chemical potential of the light massless quarks $`q=u,d`$, depending on a temperature $`T`$ . A chemical potential of light antiquarks amounts to $`\mu (T)`$. A positively defined $`\mu (T)`$ provides an abundance of light quarks with respect to light antiquarks for a thermalized state . A chemical potential $`\mu (T)`$ is a phenomenological parameter of the approach which we would fix below .
The probability for gluons to have a momentum $`\stackrel{}{p}`$ at a temperature $`T`$ is given by the Bose–Einstein distribution function
$`n_g(\stackrel{}{p})={\displaystyle \frac{1}{e^{p/T}1}}.`$ (1.2)
Since a strangeness of the colliding heavy–ions amounts to zero, the densities of strange quarks and antiquarks should be equal. The former implies a zero–value of a chemical potential $`\mu _s=\mu _{\overline{s}}=0`$. In this case the probabilities of strange quarks and antiquarks can be given by
$`n_s(\stackrel{}{p})=n_{\overline{s}}(\stackrel{}{p})={\displaystyle \frac{1}{e^{\sqrt{\stackrel{}{p}^{\mathrm{\hspace{0.17em}\hspace{0.17em}2}}+m_s^2}/T}+1}},`$ (1.3)
where $`m_s=135\mathrm{MeV}`$ is the mass of the strange quark and antiquark. The value of the current $`s`$–quark mass $`m_s=135\mathrm{MeV}`$ has been successfully applied to the calculation of chiral corrections to the amplitudes of low–energy interactions, form factors and mass spectra of low–lying hadrons and charmed heavy–light mesons . Unlike the massless antiquarks $`\overline{u}`$ and $`\overline{d}`$ for which the suppression is caused by a chemical potential $`\mu (T)`$, the strange quarks and antiquarks are suppressed by virtue of the non–zero mass $`m_s`$.
In Ref. we have supposed that a chemical potential $`\mu (T)`$, a phenomenological parameter of the description of the QGP state as a thermalized quark–gluon system at a temperature $`T`$, is an intrinsic characteristic of a thermalized quark–gluon system. Thereby, if the QGP is an excited state of the QCD vacuum, so a chemical potential should exist not only for ultra–relativistic heavy–ion collisions. Quark distribution functions of a thermalized quark–gluon system at a temperature $`T`$ should be characterized by a chemical potential $`\mu (T)`$ for any external state and any external conditions. Since any state of a thermalized system is closely related to external conditions, in order to obtain $`\mu (T)`$ we need only to specify the external conditions of a thermalized quark–gluon system the convenient for the determination of $`\mu (T)`$.
Indeed, it is well–known that the Helmholtz free energy $`F(T,V,N)`$ defining the partition function $`Z(T,V,N)`$, $`F(T,V,N)=T\mathrm{}nZ(T,V,N)`$ which plays a central role in studying thermalized systems, is nothing more than a work for an isothermic process. Therefore, by producing external conditions keeping $`T=const`$ and measuring a work one can get a full information about the Helmholtz free energy $`F_{\mathrm{exp}}(T,V,N)`$. Then, in terms of this Helmholtz free energy $`F_{\mathrm{exp}}(T,V,N)`$ one can obtain the partition function $`Z_{\mathrm{exp}}(T,V,N)`$ which can be applied to the description of the thermalized system at any $`T`$.
Following this idea we have fixed the chemical potential $`\mu (T)`$ in the form :
$`{\displaystyle \frac{\mu (T)}{\mu _0}}=\left[{\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2}}\sqrt{1+{\displaystyle \frac{4\pi ^6}{27}}\left({\displaystyle \frac{T}{\mu _0}}\right)^6}\right]^{1/3}\left[{\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2}}\sqrt{1+{\displaystyle \frac{4\pi ^6}{27}}\left({\displaystyle \frac{T}{\mu _0}}\right)^6}\right]^{1/3}.`$ (1.4)
In the low–temperature limit $`T0`$ we get
$`\mu (T)=\mu _0\left[1{\displaystyle \frac{\pi ^2}{3}}{\displaystyle \frac{T^2}{\mu _0^2}}+O\left(T^6\right)\right].`$ (1.5)
where $`\mu _0=\mu (0)=250\mathrm{MeV}`$ is a chemical potential at zero temperature . The $`T`$–dependence of a chemical potential given by Eq.(1.5) differs by a factor $`1/4`$ from the low–temperature behaviour of a chemical potential of a thermalized electron gas . The former is caused by the contribution of antiquarks.
In the high–temperature limit $`T\mathrm{}`$ a chemical potential $`\mu (T)`$ defined by Eq.(1.4) drops like $`T^2`$:
$`\mu (T)={\displaystyle \frac{\mu _0^3}{\pi ^2}}{\displaystyle \frac{1}{T^2}}+O\left(T^7\right).`$ (1.6)
A chemical potential drops very swiftly when a temperature increases. Indeed, at $`T=160\mathrm{MeV}`$ we obtain $`\mu (T)\mu _0/4`$, while at $`T=\mu _0`$ a value of a chemical potential makes up about tenth part of $`\mu _0`$, i.e. $`\mu (T)\mu _0/10`$. This implies that at very high temperatures the function $`\nu (T)=\mu (T)/T`$ becomes small and the contribution of a chemical potential of light quarks and antiquarks can be taken into account perturbatively. This assumes in particular that at temperatures $`T\mu _0=250\mathrm{MeV}`$ the number of light antiquarks will not be suppressed by a chemical potential relative to the number of light quarks.
In our approach the multiplicities of hadron production we define in terms of quark and anti–quark distributions functions in a way similar to a simple coalescence model but for correlated quarks and anti–quarks. Indeed, in a coalescence model quarks and anti–quarks are uncorrelated . This allows to introduce separately the number of light quarks $`q`$ and light anti–quarks $`\overline{q}`$ and the number of strange quarks $`s`$ and strange anti–quarks $`\overline{s}`$ . The subsequent calculation of multiplicities of hadrons in a simple coalescence model resembles a quark counting. In fact, multiplicities of hadrons are proportional to products of the number of quarks $`(q,s)`$ and anti–quarks $`(\overline{q},\overline{s})`$ in accord the naive quark structure of hadrons. Since quarks and anti–quarks do not correlate, so that the multiplicities of hadrons turn out to be independent on the momenta of hadrons <sup>1</sup><sup>1</sup>1The numbers of quarks $`(q,s)`$ and anti–quarks $`(\overline{q},\overline{s})`$ and the coefficients of proportionality are free parameters of a simple coalescence model. Therefore, a simple coalescence model contains seven free parameters. Five of them can be fixed from experimental data ..
In our approach the multiplicities of hadrons produced from the QGP phase are described by momentum integrals of quark and anti–quark distribution functions. Unlike a coalescence model these integrals depend explicitly on the momenta of hadrons, temperature $`T`$ and chemical potential $`\mu (T)`$.
For example, the multiplicities of the production of the $`K^\pm `$ mesons, $`N_{K^\pm }(\stackrel{}{q},T)`$, and $`\pi ^\pm `$ mesons, $`N_{\pi ^\pm }(\stackrel{}{q},T)`$, are defined in our approach as follows :
$`N_{K^+}(\stackrel{}{q},T)`$ $`=`$ $`3\times {\displaystyle \frac{C_M}{(M_KF_K)^{3/2}}}{\displaystyle \frac{d^3p}{(2\pi )^3}\frac{1}{e^{\nu \left(T\right)+\left|\stackrel{}{p}\stackrel{}{q}\right|/T}+1}\frac{1}{e^{\sqrt{\stackrel{}{p}^{\mathrm{\hspace{0.17em}\hspace{0.17em}2}}+m_s^2}/T}+1}},`$
$`N_K^{}(\stackrel{}{q},T)`$ $`=`$ $`3\times {\displaystyle \frac{C_M}{(M_KF_K)^{3/2}}}{\displaystyle \frac{d^3p}{(2\pi )^3}\frac{1}{e^{\nu \left(T\right)+\left|\stackrel{}{p}\stackrel{}{q}\right|/T}+1}\frac{1}{e^{\sqrt{\stackrel{}{p}^{\mathrm{\hspace{0.17em}\hspace{0.17em}2}}+m_s^2}/T}+1}},`$
$`N_{\pi ^\pm }(\stackrel{}{q},T)`$ $`=`$ $`3\times {\displaystyle \frac{C_M}{(M_\pi F_\pi )^{3/2}}}{\displaystyle \frac{d^3p}{(2\pi )^3}\frac{1}{e^{\nu \left(T\right)+\left|\stackrel{}{p}\stackrel{}{q}\right|/T}+1}\frac{1}{e^{\nu \left(T\right)+p/T}+1}},`$ (1.7)
where $`\stackrel{}{p}`$ is a relative momentum of quarks and anti–quarks coalesced in to a meson with a 3–momentum $`\stackrel{}{q}`$ at a temperature $`T`$. Since the main contribution to the integrals comes from the relative momenta of order $`pT`$, so that quark and anti–quarks coalesce at relative momenta of order $`pT`$. This agrees with the order of transversal momenta of hadrons, $`q_{}2÷3T`$, coupled in the center of mass frame of heavy–ion collisions. The factor $`3`$ corresponds the number of quark color degrees of freedom, $`M_K=500\mathrm{MeV}`$, $`F_K=160\mathrm{MeV}`$, $`M_\pi =140\mathrm{MeV}`$ and $`F_\pi =131\mathrm{MeV}`$ are the masses and the leptonic coupling constants of the $`K`$ and $`\pi `$ mesons, respectively . The dimensionless parameter $`C_M`$ is a free parameter of the approach. It is the same for all low–lying mesons.
The multiplicities of the vector meson production, for example, such as $`K^\pm `$ and $`\rho ^\pm `$ we define as
$`N_{K^+}(\stackrel{}{q},T)`$ $`=`$ $`3\times {\displaystyle \frac{C_M}{(M_K^{}F_K)^{3/2}}}{\displaystyle \frac{d^3p}{(2\pi )^3}\frac{1}{e^{\nu \left(T\right)+\left|\stackrel{}{p}\stackrel{}{q}\right|/T}+1}\frac{1}{e^{\sqrt{\stackrel{}{p}^{\mathrm{\hspace{0.17em}\hspace{0.17em}2}}+m_s^2}/T}+1}},`$
$`N_K^{}(\stackrel{}{q},T)`$ $`=`$ $`3\times {\displaystyle \frac{C_M}{(M_K^{}F_K)^{3/2}}}{\displaystyle \frac{d^3p}{(2\pi )^3}\frac{1}{e^{\nu \left(T\right)+\left|\stackrel{}{p}\stackrel{}{q}\right|/T}+1}\frac{1}{e^{\sqrt{\stackrel{}{p}^{\mathrm{\hspace{0.17em}\hspace{0.17em}2}}+m_s^2}/T}+1}},`$
$`N_{\rho ^\pm }(\stackrel{}{q},T)`$ $`=`$ $`3\times {\displaystyle \frac{C_M}{(M_\rho F_\pi )^{3/2}}}{\displaystyle \frac{d^3p}{(2\pi )^3}\frac{1}{e^{\nu \left(T\right)+\left|\stackrel{}{p}\stackrel{}{q}\right|/T}+1}\frac{1}{e^{\nu \left(T\right)+p/T}+1}},`$ (1.8)
where $`M_K^{}=892\mathrm{MeV}`$ and $`M_\rho =770\mathrm{MeV}`$ are the masses of the $`K^{}`$ and $`\rho `$ mesons, respectively .
In the case of baryons and antibaryons we suggest to define the multiplicities by using the diquark–quark picture of baryons and antibaryons. For example, the multiplicities of the proton ($`p`$) and antiproton ($`\overline{p}`$) we write in the form
$`N_p(\stackrel{}{q},T)`$ $`=`$ $`{\displaystyle \frac{3!}{3!}}\times {\displaystyle \frac{C_B}{(M_pF_\pi )^{3/2}}}`$
$`\times `$ $`{\displaystyle \frac{d^3p}{(2\pi )^3}\frac{1}{\left(e^{\nu \left(T\right)+p/T}+1\right)^2}\frac{1}{e^{\nu \left(T\right)+\left|\stackrel{}{p}\stackrel{}{q}\right|/T}+1}},`$
$`N_{\overline{p}}(\stackrel{}{q},T)`$ $`=`$ $`{\displaystyle \frac{3!}{3!}}\times {\displaystyle \frac{C_{\overline{B}}}{(M_pF_\pi )^{3/2}}}`$ (1.9)
$`\times `$ $`{\displaystyle \frac{d^3p}{(2\pi )^3}\frac{1}{\left(e^{\nu \left(T\right)+p/T}+1\right)^2}\frac{1}{e^{\nu \left(T\right)+\left|\stackrel{}{p}\stackrel{}{q}\right|/T}+1}},`$
where a momentum $`\stackrel{}{p}`$ has a meaning of a relative momentum of three–quark (three–anti–quark) system, $`M_p=940\mathrm{MeV}`$ is the mass of the proton and antiproton. As well as in the meson case the main contribution to the momentum integrals comes from the momenta of order $`pT`$ providing a coalescence of three quarks (three anti–quarks) into baryons (anti–baryons) at the momenta of order $`pT`$. That is again of order of transversal momenta of the produced hadrons, $`q_{}(2÷3)T`$, coupled in the center of mass frame of heavy–ion collisions. The factor $`3!`$ in the numerator is related to the quark colour degrees of freedom and defined by $`\epsilon _{ijk}\epsilon ^{ijk}=3!`$, where $`i,j`$ and $`k`$ are colour indices and run over $`i=1,2,3`$ each. In turn, in the denominator the factor $`3!`$ takes into account the identity of three light quarks $`(qqq)`$ and three antiquarks $`(\overline{q}\overline{q}\overline{q})`$. In the isotopical limit we do not distinguish $`u`$ and $`d`$ quarks as well as $`\overline{u}`$ and $`\overline{d}`$ antiquarks. The dimensionless parameters $`C_B`$ and $`C_{\overline{B}}`$ are free parameters of the approach. Each of them is equal for all components of octets of baryons and antibaryons, respectively, but $`C_BC_{\overline{B}}`$.
The paper is organized as follows. In Sect. 2 we calculate the theoretical values of multiplicities of the hadron production from the thermalized QGP phase. The theoretical predictions and experimental data are adduced in Table 1. In the Conclusion we discuss the obtained results. A possible estimate of the absolute values of our input parameters is discussed through the application of our approach to the calculation of the number of baryons and antibaryons relative to the number of photons at the early stage of the evolution of the Universe assuming that this evolution goes through the intermediate thermalized QGP phase.
## 2 Multiplicities of hadron production from the thermalized QGP phase
Now let us proceed to the evaluation of multiplicities of hadron production from the thermalized QGP phase of QCD. The theoretical predictions for the different ratios of hadron multiplicities we compare with experimental data adduced in Table I of Ref. . These are the data of various experiments given by NA44, NA49, NA50 and WA97 Collaborations on Pb+Pb collisions at 158 GeV/nucleon. Also we compare our results with the experimental data obtained by NA35 Collaboration on S+S collisions and NA38 Collaboration on O+U and S+U collisions at 200 GeV/nucleon. From Table 1 of this paper one can see that in the whole the experimental data for the hadron production are obtained for rapidities ranging over the region $`2.3y4.1`$. The relation between a 3–momentum $`q`$ and a rapidity $`y`$ reads
$`q=\sqrt{M^2\mathrm{sh}^2y+q_{}^2\mathrm{ch}^2y}M\mathrm{sh}y,`$ (2.1)
where $`M`$ and $`\stackrel{}{q}_{}`$ are the mass and the transversal momentum of the produced hadron. For rapidities $`y[2.3,4.1]`$ we get
$`q(5.0÷30.2)M.`$ (2.2)
Thus, for $`K`$ mesons and hadrons heavier than $`K`$ mesons typical momenta are of order of 2.5 GeV and greater. This gives a possibility to investigate the momentum integrals defining multiplicities of the hadron production at $`q\mathrm{}`$. As has been shown in Ref. the ratios of the multiplicities $`R_{K^+K^{}}(q,T)=N_{K^+}(\stackrel{}{q},T)/N_K^{}(\stackrel{}{q},T)`$ and $`R_{K^+\pi ^+}(q,T)=N_{K^+}(\stackrel{}{q},T)/N_{\pi ^+}(\stackrel{}{q},T)`$ are smooth functions of $`q`$, wobbling slightly around the asymptotic values obtained at $`q\mathrm{}`$, and describe good the experimental data at $`T=175\mathrm{MeV}`$. Below the theoretical results on the multiplicities of the hadron production we would compare with experimental data at $`T=175\mathrm{MeV}`$. In this case it is obvious that the typical momenta of the momentum integrals defining the multiplicities of hadron production are of order $`pT`$. Therefore, in our approach the momenta at which quarks and anti–quarks coalesces into hadrons should be of order $`pT`$. This agrees with the order of transversal momenta of hadrons, $`q_{}2÷3T`$, coupled in the center of mass frame of heavy–ion collisions.
Hence, for rapidities ranging over the region $`2.3y4.1`$ the typical total 3–momenta of hadrons are greater than 1 GeV. Thereby, multiplicities of the hadron production can be calculated in the asymptotic regime at $`qT`$.
We would like to accentuate that the relative momenta at which quarks and anti–quarks coalesce into hadrons are described effectively by the momentum of integration $`\stackrel{}{p}`$.
At $`qT`$ the multiplicities of the hadron production defined by Eqs.(1)–(1) can be represented in the following form
$`N_{\pi ^+}(\stackrel{}{q},T)`$ $`=`$ $`N_\pi ^{}(\stackrel{}{q},T)=N_{\pi ^0}(\stackrel{}{q},T)={\displaystyle \frac{3C_M}{(M_\pi F_\pi )^{3/2}}}e^{q/T}I_\pi (T),`$
$`N_{K^+}(\stackrel{}{q},T)`$ $`=`$ $`N_{K^0}(\stackrel{}{q},T)={\displaystyle \frac{3C_M}{(M_KF_K)^{3/2}}}e^{q/T}e^{+\nu \left(T\right)}I_K(T),`$
$`N_K^{}(\stackrel{}{q},T)`$ $`=`$ $`N_{\overline{K}^0}(\stackrel{}{q},T)={\displaystyle \frac{3C_M}{(M_KF_K)^{3/2}}}e^{q/T}e^{\nu \left(T\right)}I_{\overline{K}}(T),`$
$`N_{K_S^0}(\stackrel{}{q},T)`$ $`=`$ $`{\displaystyle \frac{1}{2}}N_{K^0}(\stackrel{}{q},T)+{\displaystyle \frac{1}{2}}N_{\overline{K}^0}(\stackrel{}{q},T)={\displaystyle \frac{1}{2}}{\displaystyle \frac{3C_M}{(M_KF_K)^{3/2}}}e^{q/T}e^{+\nu \left(T\right)}I_{K_S^0}(T),`$
$`N_\eta (\stackrel{}{q},T)`$ $`=`$ $`\mathrm{sin}^2\overline{\theta }{\displaystyle \frac{3C_M}{(M_\eta F_\pi )^{3/2}}}e^{q/T}I_\pi (T)+\mathrm{cos}^2\overline{\theta }{\displaystyle \frac{3C_M}{(M_\eta F_S)^{3/2}}}e^{q/T}I_\eta (T),`$
$`N_\varphi (\stackrel{}{q},T)`$ $`=`$ $`{\displaystyle \frac{3C_M}{(M_\varphi F_S)^{3/2}}}e^{q/T}I_\varphi (T),`$
$`N_p(\stackrel{}{q},T)`$ $`=`$ $`{\displaystyle \frac{C_B}{(M_pF_\pi )^{3/2}}}e^{q/T}e^{+\nu \left(T\right)}I_p(T),`$
$`N_\mathrm{\Lambda }(\stackrel{}{q},T)`$ $`=`$ $`{\displaystyle \frac{3C_B}{(M_\mathrm{\Lambda }F_K)^{3/2}}}e^{q/T}e^{+2\nu \left(T\right)}I_\mathrm{\Lambda }(T),`$
$`N_\mathrm{\Xi }(\stackrel{}{q},T)`$ $`=`$ $`{\displaystyle \frac{3C_B}{(M_\mathrm{\Xi }F_S)^{3/2}}}e^{q/T}e^{+\nu \left(T\right)}I_\mathrm{\Xi }(T),`$
$`N_\mathrm{\Omega }(\stackrel{}{q},T)`$ $`=`$ $`{\displaystyle \frac{C_B}{(M_\mathrm{\Omega }F_S)^{3/2}}}e^{q/T}I_\mathrm{\Omega }(T),`$
$`N_{\overline{p}}(\stackrel{}{q},T)`$ $`=`$ $`{\displaystyle \frac{C_{\overline{B}}}{(M_pF_\pi )^{3/2}}}e^{q/T}e^{\nu \left(T\right)}I_{\overline{p}}(T),`$
$`N_{\overline{\mathrm{\Lambda }}}(\stackrel{}{q},T)`$ $`=`$ $`{\displaystyle \frac{3C_{\overline{B}}}{(M_\mathrm{\Lambda }F_K)^{3/2}}}e^{q/T}e^{\mathrm{\hspace{0.17em}2}\nu \left(T\right)}I_{\overline{\mathrm{\Lambda }}}(T),`$
$`N_{\overline{\mathrm{\Xi }}}(\stackrel{}{q},T)`$ $`=`$ $`{\displaystyle \frac{3C_{\overline{B}}}{(M_\mathrm{\Xi }F_S)^{3/2}}}e^{q/T}e^{\nu \left(T\right)}I_{\overline{\mathrm{\Xi }}}(T),`$
$`N_{\overline{\mathrm{\Omega }}}(\stackrel{}{q},T)`$ $`=`$ $`{\displaystyle \frac{C_{\overline{B}}}{(M_\mathrm{\Omega }F_S)^{3/2}}}e^{q/T}I_{\overline{\mathrm{\Omega }}}(T),`$ (2.3)
where the structure functions $`I_i(T)(i=\pi ,K,\overline{K},\mathrm{})`$ are defined by
$`I_\pi (T)=e^{+\nu \left(T\right)}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dp}{4\pi ^2}}{\displaystyle \frac{p^2}{e^{+\nu \left(T\right)+p/T}+1}}+e^{\nu \left(T\right)}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dp}{4\pi ^2}}{\displaystyle \frac{p^2}{e^{\nu \left(T\right)+p/T}+1}}=`$ (2.4)
$`={\displaystyle \frac{T^3}{4\pi ^2}}\left[e^{+\nu \left(T\right)}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dxx^2}{e^{+\nu \left(T\right)+x}+1}}+e^{\nu \left(T\right)}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dxx^2}{e^{\nu \left(T\right)+x}+1}}\right]=3.153{\displaystyle \frac{T^3}{4\pi ^2}},`$
$`I_K(T)={\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dp}{4\pi ^2}}{\displaystyle \frac{p^2}{e^{\sqrt{p^{\mathrm{\hspace{0.17em}2}}+m_s^2}/T}+1}}+e^{\nu \left(T\right)}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dp}{4\pi ^2}}{\displaystyle \frac{p^2}{e^{\nu \left(T\right)+p/T}+1}}=`$
$`={\displaystyle \frac{T^3}{4\pi ^2}}\left[{\displaystyle \frac{m_s^3}{T^3}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dxx^2}{e^{\left(m_s/T\right)\sqrt{1+x^2}}+1}}+e^{\nu \left(T\right)}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dxx^2}{e^{\nu \left(T\right)+x}+1}}\right]=2.924{\displaystyle \frac{T^3}{4\pi ^2}},`$
$`I_{\overline{K}}(T)={\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dp}{4\pi ^2}}{\displaystyle \frac{p^2}{e^{\sqrt{p^{\mathrm{\hspace{0.17em}2}}+m_s^2}/T}+1}}+e^{+\nu \left(T\right)}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dp}{4\pi ^2}}{\displaystyle \frac{p^2}{e^{+\nu \left(T\right)+p/T}+1}}=`$
$`={\displaystyle \frac{T^3}{4\pi ^2}}\left[{\displaystyle \frac{m_s^3}{T^3}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dxx^2}{e^{\left(m_s/T\right)\sqrt{1+x^2}}+1}}+e^{+\nu \left(T\right)}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dxx^2}{e^{+\nu \left(T\right)+x}+1}}\right]=3.463{\displaystyle \frac{T^3}{4\pi ^2}},`$
$`I_\eta (T)=I_\varphi (T)=`$
$`={\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dp}{2\pi ^2}}{\displaystyle \frac{p^2}{e^{\sqrt{\stackrel{}{p}^{\mathrm{\hspace{0.17em}\hspace{0.17em}2}}+m_s^2}/T}+1}}={\displaystyle \frac{m_s^3}{2\pi ^2}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dx^2}{e^{\left(m_s/T\right)\sqrt{1+x^2}}+1}}=3.522{\displaystyle \frac{m_s^3}{2\pi ^2}},`$
$`I_p(T)={\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dp}{2\pi ^2}}{\displaystyle \frac{1}{\left(e^{\nu \left(T\right)+p/T}+1\right)^2}}={\displaystyle \frac{T^3}{2\pi ^2}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dxx^2}{\left(e^{\nu \left(T\right)+x}+1\right)^2}}=0.253{\displaystyle \frac{T^3}{2\pi ^2}},`$
$`I_\mathrm{\Lambda }(T)=e^{\nu \left(T\right)}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dp}{4\pi ^2}}{\displaystyle \frac{p^2}{e^{\sqrt{p^{\mathrm{\hspace{0.17em}2}}+m_s^2}/T}+1}}{\displaystyle \frac{1}{e^{\nu \left(T\right)+p/T}+1}}`$
$`+e^{2\nu \left(T\right)}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dp}{4\pi ^2}}{\displaystyle \frac{p^2}{\left(e^{\nu \left(T\right)+p/T}+1\right)^2}}=`$
$`={\displaystyle \frac{m_s^3}{4\pi ^2}}[e^{\nu \left(T\right)}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dxx^2}{e^{\left(m_s/T\right)\sqrt{1+x^2}}+1}}{\displaystyle \frac{1}{e^{\nu \left(T\right)+\left(m_s/T\right)x}+1}}`$
$`+e^{\mathrm{\hspace{0.17em}2}\nu \left(T\right)}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dxx^2}{\left(e^{\nu \left(T\right)+\left(m_s/T\right)x}+1\right)^2}}]=0.582{\displaystyle \frac{m_s^3}{4\pi ^2}},`$
$`I_\mathrm{\Xi }(T)={\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dp}{4\pi ^2}}{\displaystyle \frac{p^2}{\left(e^{\sqrt{p^{\mathrm{\hspace{0.17em}2}}+m_s^2}/T}+1\right)^2}}`$
$`+e^{\nu \left(T\right)}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dp}{4\pi ^2}}{\displaystyle \frac{p^2}{e^{\sqrt{p^{\mathrm{\hspace{0.17em}2}}+m_s^2}/T}+1}}{\displaystyle \frac{1}{e^{\nu \left(T\right)+p/T}+1}}=`$
$`={\displaystyle \frac{m_s^3}{4\pi ^2}}[{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dxx^2}{\left(e^{\left(m_s/T\right)\sqrt{1+x^2}}+1\right)^2}}+e^{\nu \left(T\right)}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dxx^2}{e^{\left(m_s/T\right)\sqrt{1+x^2}}+1}}`$
$`\times {\displaystyle \frac{1}{e^{\nu \left(T\right)+\left(m_s/T\right)x}+1}}]=0.528{\displaystyle \frac{m_s^3}{4\pi ^2}},`$
$`I_\mathrm{\Omega }(T)={\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dp}{2\pi ^2}}{\displaystyle \frac{p^2}{\left(e^{\sqrt{p^{\mathrm{\hspace{0.17em}2}}+m_s^2}/T}+1\right)^2}}={\displaystyle \frac{m_s^3}{2\pi ^2}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dxx^2}{\left(e^{\left(m_s/T\right)\sqrt{1+x^2}}+1\right)^2}}=`$
$`=0.251{\displaystyle \frac{m_s^3}{2\pi ^2}},`$
$`I_{\overline{p}}(T)={\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dp}{2\pi ^2}}{\displaystyle \frac{p^2}{\left(e^{\nu \left(T\right)+p/T}+1\right)^2}}={\displaystyle \frac{T^3}{2\pi ^2}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dxx^2}{\left(e^{\nu \left(T\right)+x}+1\right)^2}}=0.097{\displaystyle \frac{T^3}{2\pi ^2}},`$
$`\times {\displaystyle \frac{1}{e^{+\nu \left(T\right)+\left(m_s/T\right)x}+1}}]=0.558{\displaystyle \frac{m_s^3}{4\pi ^2}},`$
$`I_{\overline{\mathrm{\Lambda }}}(T)=e^{+\nu \left(T\right)}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dp}{4\pi ^2}}{\displaystyle \frac{p^2}{e^{\sqrt{p^{\mathrm{\hspace{0.17em}2}}+m_s^2}/T}+1}}{\displaystyle \frac{1}{e^{+\nu \left(T\right)+p/T}+1}}`$
$`+e^{+2\nu \left(T\right)}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dp}{4\pi ^2}}{\displaystyle \frac{p^2}{\left(e^{+\nu \left(T\right)+p/T}+1\right)^2}}=`$
$`=`$ $`{\displaystyle \frac{m_s^3}{4\pi ^2}}[e^{+\nu \left(T\right)}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dxx^2}{e^{\left(m_s/T\right)\sqrt{1+x^2}}+1}}{\displaystyle \frac{1}{e^{+\nu \left(T\right)+\left(m_s/T\right)x}+1}}`$
$`+e^{+\mathrm{\hspace{0.17em}2}\nu \left(T\right)}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dxx^2}{\left(e^{+\nu \left(T\right)+\left(m_s/T\right)x}+1\right)^2}}]=0.686{\displaystyle \frac{m_s^3}{4\pi ^2}},`$
$`I_{\overline{\mathrm{\Xi }}}(T)={\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dp}{4\pi ^2}}{\displaystyle \frac{p^2}{\left(e^{\sqrt{p^{\mathrm{\hspace{0.17em}2}}+m_s^2}/T}+1\right)^2}}`$
$`+e^{+\nu \left(T\right)}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dp}{4\pi ^2}}{\displaystyle \frac{p^2}{e^{\sqrt{p^{\mathrm{\hspace{0.17em}2}}+m_s^2}/T}+1}}{\displaystyle \frac{1}{e^{+\nu \left(T\right)+p/T}+1}}=`$
$`={\displaystyle \frac{m_s^3}{4\pi ^2}}[{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dxx^2}{\left(e^{\left(m_s/T\right)\sqrt{1+x^2}}+1\right)^2}}+e^{+\nu \left(T\right)}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dxx^2}{e^{\left(m_s/T\right)\sqrt{1+x^2}}+1}}`$
$`\times {\displaystyle \frac{1}{e^{+\nu \left(T\right)+\left(m_s/T\right)x}+1}}]=0.558{\displaystyle \frac{m_s^3}{4\pi ^2}},`$
$`I_{\overline{\mathrm{\Omega }}}(T)=I_\mathrm{\Omega }(T).`$
The numerical values of the integrals are obtained at $`m_s=135\mathrm{MeV}`$ and $`T=175\mathrm{MeV}`$.
The theoretical ratios of multiplicities of the hadron production which we compare with measured experimentally we define as follows
$`R_{K^+K^{}}(q,T)`$ $`=`$ $`{\displaystyle \frac{N_{K^+}(\stackrel{}{q},T)}{N_K^{}(\stackrel{}{q},T)}}=e^{+\mathrm{\hspace{0.17em}2}\nu \left(T\right)}{\displaystyle \frac{I_K(T)}{I_{\overline{K}}(T)}}=1.520,`$
$`R_{K^+\pi ^+}(q,T)`$ $`=`$ $`{\displaystyle \frac{N_{K^+}(\stackrel{}{q},T)}{N_{\pi ^+}(\stackrel{}{q},T)}}=\left({\displaystyle \frac{M_\pi F_\pi }{M_KF_K}}\right)^{3/2}e^{+\nu \left(T\right)}{\displaystyle \frac{I_K(T)}{I_\pi (T)}}=0.139,`$
$`R_{K^{}\pi ^{}}(q,T)`$ $`=`$ $`{\displaystyle \frac{N_K^{}(\stackrel{}{q},T)}{N_\pi ^{}(\stackrel{}{q},T)}}=\left({\displaystyle \frac{M_\pi F_\pi }{M_KF_K}}\right)^{3/2}e^{\nu \left(T\right)}{\displaystyle \frac{I_{\overline{K}}(T)}{I_\pi (T)}}=0.090,`$
$`R_{K_S^0\pi ^{}}(q,T)`$ $`=`$ $`{\displaystyle \frac{N_{K_S^0}(\stackrel{}{q},T)}{N_\pi ^{}(\stackrel{}{q},T)}}={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{M_\pi F_\pi }{M_KF_K}}\right)^{3/2}e^{+\nu \left(T\right)}{\displaystyle \frac{I_{K_S^0}(T)}{I_\pi ^{}(T)}}=0.113,`$
$`R_{\mathrm{\Xi }\mathrm{\Lambda }}(q,T)`$ $`=`$ $`{\displaystyle \frac{N_\mathrm{\Xi }(\stackrel{}{q},T)}{N_\mathrm{\Lambda }(\stackrel{}{q},T)}}=\left({\displaystyle \frac{M_\mathrm{\Lambda }F_K}{M_\mathrm{\Xi }F_S}}\right)^{3/2}e^{\nu \left(T\right)}{\displaystyle \frac{I_\mathrm{\Xi }(T)}{I_\mathrm{\Lambda }(T)}}=0.108,`$
$`R_{\mathrm{\Omega }\mathrm{\Xi }}(q,T)`$ $`=`$ $`{\displaystyle \frac{N_\mathrm{\Omega }(\stackrel{}{q},T)}{N_\mathrm{\Xi }(\stackrel{}{q},T)}}={\displaystyle \frac{1}{3}}\left({\displaystyle \frac{M_\mathrm{\Xi }}{M_\mathrm{\Omega }}}\right)^{3/2}e^{\nu \left(T\right)}{\displaystyle \frac{I_\mathrm{\Omega }(T)}{I_\mathrm{\Xi }(T)}}=0.166,`$
$`R_{\overline{\mathrm{\Lambda }}\overline{p}}(q,T)`$ $`=`$ $`{\displaystyle \frac{N_{\overline{\mathrm{\Lambda }}}(\stackrel{}{q},T)}{N_{\overline{p}}(\stackrel{}{q},T)}}=3\left({\displaystyle \frac{M_pF_\pi }{M_\mathrm{\Lambda }F_K}}\right)^{3/2}e^{\nu \left(T\right)}{\displaystyle \frac{I_{\overline{\mathrm{\Lambda }}}(T)}{I_{\overline{p}}(T)}}=2.081,`$
$`R_{\overline{\mathrm{\Xi }}\overline{\mathrm{\Lambda }}}(q,T)`$ $`=`$ $`{\displaystyle \frac{N_{\overline{\mathrm{\Xi }}}(\stackrel{}{q},T)}{N_{\overline{\mathrm{\Lambda }}}(\stackrel{}{q},T)}}=\left({\displaystyle \frac{M_\mathrm{\Lambda }F_K}{M_\mathrm{\Xi }F_S}}\right)^{3/2}e^{+\nu \left(T\right)}{\displaystyle \frac{I_{\overline{\mathrm{\Xi }}}(T)}{I_{\overline{\mathrm{\Lambda }}}(T)}}=0.173,`$
$`R_{\overline{\mathrm{\Omega }}\overline{\mathrm{\Xi }}}(q,T)`$ $`=`$ $`{\displaystyle \frac{N_{\overline{\mathrm{\Omega }}}(\stackrel{}{q},T)}{N_{\overline{\mathrm{\Xi }}}(\stackrel{}{q},T)}}={\displaystyle \frac{1}{3}}\left({\displaystyle \frac{M_\mathrm{\Xi }}{M_\mathrm{\Omega }}}\right)^{3/2}e^{+\nu \left(T\right)}{\displaystyle \frac{I_{\overline{\mathrm{\Omega }}}(T)}{I_{\overline{\mathrm{\Xi }}}(T)}}=0.282,`$
$`R_{\overline{\mathrm{\Omega }}\mathrm{\Omega }}(q,T)`$ $`=`$ $`{\displaystyle \frac{N_{\overline{\mathrm{\Omega }}}(\stackrel{}{q},T)}{N_\mathrm{\Omega }(\stackrel{}{q},T)}}={\displaystyle \frac{C_{\overline{B}}}{C_B}}=R_{\overline{\mathrm{\Omega }}\mathrm{\Omega }}^{\mathrm{exp}}=0.46\pm 0.15,`$
$`R_{\overline{\mathrm{\Lambda }}\mathrm{\Lambda }}(q,T)`$ $`=`$ $`{\displaystyle \frac{N_{\overline{\mathrm{\Lambda }}}(\stackrel{}{q},T)}{N_\mathrm{\Lambda }(\stackrel{}{q},T)}}={\displaystyle \frac{C_{\overline{B}}}{C_B}}\times e^{\mathrm{\hspace{0.17em}4}\nu \left(T\right)}{\displaystyle \frac{I_{\overline{\mathrm{\Lambda }}}(T)}{I_\mathrm{\Lambda }(T)}}=0.168\pm 0.055,`$
$`R_{\overline{\mathrm{\Xi }}\mathrm{\Xi }}(q,T)`$ $`=`$ $`{\displaystyle \frac{N_{\overline{\mathrm{\Xi }}}(\stackrel{}{q},T)}{N_\mathrm{\Xi }(\stackrel{}{q},T)}}={\displaystyle \frac{C_{\overline{B}}}{C_B}}\times e^{\mathrm{\hspace{0.17em}2}\nu \left(T\right)}{\displaystyle \frac{I_{\overline{\mathrm{\Xi }}}(T)}{I_\mathrm{\Xi }(T)}}=0.270\pm 0.088,`$
$`R_{\eta \pi ^0}(q,T)`$ $`=`$ $`{\displaystyle \frac{N_\eta (\stackrel{}{q},T)}{N_{\pi ^0}(\stackrel{}{q},T)}}=\mathrm{sin}^2\overline{\theta }\left({\displaystyle \frac{M_\pi }{M_\eta }}\right)^{3/2}+\mathrm{cos}^2\overline{\theta }\left({\displaystyle \frac{M_\pi F_\pi }{M_\eta F_S}}\right)^{3/2}{\displaystyle \frac{I_\eta (T)}{I_{\pi ^+}(T)}}=0.088,`$
$`R_{\varphi \pi }(q,T)`$ $`=`$ $`{\displaystyle \frac{N_\varphi (\stackrel{}{q},T)}{N_\pi (\stackrel{}{q},T)}}=\left({\displaystyle \frac{M_\pi F_\pi }{M_\varphi F_S}}\right)^{3/2}{\displaystyle \frac{I_\varphi (T)}{I_\pi (T)}}=7.84\times 10^3,`$
$`R_{\varphi (\rho ^0+\omega ^0)}(q,T)`$ $`=`$ $`{\displaystyle \frac{N_\varphi (\stackrel{}{q},T)}{N_{\rho ^0}(\stackrel{}{q},T)+N_{\omega ^0}(\stackrel{}{q},T)}}=\left({\displaystyle \frac{M_\rho F_\pi }{M_\varphi F_S}}\right)^{3/2}{\displaystyle \frac{I_\varphi (T)}{I_\pi (T)}}=0.103,`$
$`R_{\varphi K_S^0}(q,T)`$ $`=`$ $`{\displaystyle \frac{N_\varphi (\stackrel{}{q},T)}{N_{K_S^0}(\stackrel{}{q},T)}}=2\left({\displaystyle \frac{M_KF_K}{M_\varphi F_S}}\right)^{3/2}e^{\nu \left(T\right)}{\displaystyle \frac{I_\varphi (T)}{I_{K_S^0}(T)}}=0.071,`$
$`R_{\mathrm{\Lambda }K_S^0}(q,T)`$ $`=`$ $`{\displaystyle \frac{N_\mathrm{\Lambda }(\stackrel{}{q},T)}{N_{K_S^0}(\stackrel{}{q},T)}}={\displaystyle \frac{C_B}{C_M}}\times 2e^{+\nu \left(T\right)}{\displaystyle \frac{I_\mathrm{\Lambda }(T)}{I_{K_S^0}(T)}}={\displaystyle \frac{C_B}{C_M}}\times 0.148=`$
$`=`$ $`R_{\mathrm{\Lambda }K_S^0}^{\mathrm{exp}}=0.65\pm 0.11{\displaystyle \frac{C_B}{C_M}}=4.39\pm 0.74,`$
$`R_{pK^+}(q,T)`$ $`=`$ $`{\displaystyle \frac{N_p(\stackrel{}{q},T)}{N_{K^+}(\stackrel{}{q},T)}}={\displaystyle \frac{C_B}{C_M}}\times {\displaystyle \frac{1}{3}}\left({\displaystyle \frac{M_KF_K}{M_pF_\pi }}\right)^{3/2}{\displaystyle \frac{I_p(T)}{I_K(T)}}=0.136\pm 0.023,`$
$`R_{\overline{p}K^{}}(q,T)`$ $`=`$ $`{\displaystyle \frac{N_{\overline{p}}(\stackrel{}{q},T)}{N_K^{}(\stackrel{}{q},T)}}={\displaystyle \frac{C_{\overline{B}}}{C_M}}\times {\displaystyle \frac{1}{3}}\left({\displaystyle \frac{M_KF_K}{M_pF_\pi }}\right)^{3/2}{\displaystyle \frac{I_{\overline{p}}(T)}{I_{\overline{K}}(T)}}=0.020\pm 0.007.`$ (2.5)
The constant $`F_S=3.5F_\pi `$ is related to the leptonic constant of the pseudoscalar meson containing only $`s`$–quarks, $`s\overline{s}`$ . We have estimated $`F_S`$ in agreement with the experimental data on the $`\eta (550)/\pi ^0`$ and $`\varphi (1020)/\pi `$ production. For the description of the multiplicity of the $`\eta (550)`$ meson production we have taken into account that the low–energy meson phenomenology gives the following quark structure of the $`\eta (550)`$ meson:
$`\eta (550)=(q\overline{q})\mathrm{sin}\overline{\theta }+(s\overline{s})\mathrm{cos}\overline{\theta },`$ (2.6)
where $`\overline{\theta }=\vartheta _0\vartheta _P`$ with $`\vartheta _0=35.264^0`$, the ideal mixing angle, and $`\vartheta _P`$, the octet–singlet mixing angle. Recent analysis of the value of the octet–singlet mixing angle carried out by Bramon, Escribano and Scadron gives $`\vartheta _P=\mathrm{\hspace{0.17em}16.9}\pm 1.7^0`$. For the $`\varphi (1020)`$ meson we have supposed the $`s\overline{s}`$ quark structure .
## 3 Conclusion
The theoretical and experimental values of the ratios of hadron production are adduced in Table 1. From Table 1 one can see a good agreement between presently available set of hadron ratios measured for various experiments given by NA44, NA49, NA50 and WA97 Collaborations on Pb+Pb collisions at 158 GeV/nucleon, NA35 Collaboration on S+S collisions and NA38 Collaboration on O+U and S+U collisions at 200 GeV/nucleon and theoretical predictions for the ratios of multiplicities of hadron production from the thermalized QGP phase at a temperature $`T=175\mathrm{MeV}`$. Save the ratio $`\overline{\mathrm{\Lambda }}/\overline{p}`$, $`(\overline{\mathrm{\Lambda }}/\overline{p})_{\mathrm{th}}=2.081`$ and $`(\overline{\mathrm{\Lambda }}/\overline{p})_{\mathrm{exp}}=3\pm 1`$, the theoretical results agree with the experimental ones with accuracy better than 18$`\%`$.
In the our approach multiplicities of hadron production are defined by momentum integrals on quark (antiquark) distribution functions in accordance with the phenomenological quark structure of the hadron. For the analysis of the multiplicities of the baryon and antibaryon production in terms of the quark and antiquark distribution functions we have followed the diquark–quark picture for baryons and antibaryons. This has allowed to describe multiplicities of the baryon, antibaryon and meson production on the same footing.
For the explanation of experimental data on the hadron production in ultrarelativistic heavy–ion collisions we have used three input parameters $`C_{\overline{B}}/C_B`$, $`C_M/C_B`$ and $`F_S`$. These parameters are related to the spatial volumes of hadronization of the quarks and antiquarks from the thermalized QGP phase. The first two parameters have been fixed from the experimental data on the ratios $`(\overline{\mathrm{\Omega }}/\mathrm{\Omega })_{\mathrm{exp}}=0.46\pm 0.15`$, $`(\mathrm{\Lambda }/K_S^0)_{\mathrm{exp}}=0.65\pm 0.11`$. This gives $`C_{\overline{B}}/C_B=0.46\pm 0.15`$ and $`C_M/C_B=0.23\pm 0.04`$. In turn, the value of the parameter $`F_S=3.5F_\pi =458.5\mathrm{MeV}`$ is a result of a smooth fit of the ratios of hadrons containing the $`ss`$ and $`s\overline{s}`$ components in the quark structure. In the bulk our approach to the hadronization from the QGP phase of QCD has succeeded in describing 21 experimental data on ultrarelativistic heavy–ion collisions.
Unlike other approaches \[2–4\] the ratio of the $`\overline{\mathrm{\Omega }}`$ and $`\mathrm{\Omega }`$ baryon production is an input parameter in our model $`C_{\overline{B}}/C_B`$. The ratio $`\overline{\mathrm{\Omega }}/\mathrm{\Omega }`$ does not depend in our approach on both the momenta of baryons and the temperature. The former is due to the zero–value of the $`s`$–quark chemical potential, $`\mu _s=\mu _{\overline{s}}=0`$. As a result the ratio $`\overline{\mathrm{\Omega }}/\mathrm{\Omega }`$ can be only fitted in our approach. By fitting the ratio $`\overline{\mathrm{\Omega }}/\mathrm{\Omega }`$ from experimental data and applying this value to the description of other ratios of the baryon and antibaryon production for the thermalized QGP phase we have found a good agreement with experimental data. This confirms a self–consistency of the approach.
In turn, the distinction between the parameters $`C_B`$ and $`C_{\overline{B}}`$ can be related to a well–known factor of the baryon–antibaryon asymmetry in the Universe which one could put phenomenologically at the early stage of the evolution of the Universe whether the baryon synthesis in it goes through the intermediate QGP phase. Indeed, as has been stated by Börner : Within the standard big–bang model, however, there seems to be little chance of achieving a physical separation of baryon and antibaryon phases in an initially baryon–symmetric cosmological model. If the baryon number was exactly conserved – as it assumed to be the standard model – the small asymmetry necessary for our existence must be postulated initially. Grand unified theories offer the possibility of creating this small asymmetry from physical processes.
In our approach the baryon–antibaryon asymmetry at the hadronic level can be realized phenomenologically in terms of a different rate of hadronization of baryons and antibaryons caused by the input parameter $`C_{\overline{B}}/C_B=0.46\pm 0.15`$ fixed through the experimental data on the ratio $`\overline{\mathrm{\Omega }}/\mathrm{\Omega }`$ production in ultrarelativistic heavy–ion collisions. For the total number of antibaryons $`N_{\overline{B}}`$ relative to the total number of baryons $`N_B`$ produced for the baryon–antibaryon synthesis at the early stage of the evolution of the Universe at a temperature $`T=175\mathrm{MeV}`$ and gone through the intermediate thermalized QGP phase we predict
$`{\displaystyle \frac{N_{\overline{B}}}{N_B}}=0.41\times {\displaystyle \frac{C_{\overline{B}}}{C_B}}=0.19\pm 0.06.`$ (3.1)
This result can be supported by a trivial estimate in the approximation of the equilibrium baryon and anti–baryon gases. In such an approximation the ratio $`N_{\overline{B}}/N_B`$ is defined by
$`{\displaystyle \frac{N_{\overline{B}}}{N_B}}=e^{2\mu _B\left(T\right)/T}=e^{6\mu \left(T\right)/T}=0.17,`$ (3.2)
where a chemical potential $`\mu (T)`$ is given by Eq.(1.4) and $`T=175\mathrm{MeV}`$.
Thus, at the early stage of the evolution of the Universe the number of antibaryons should be of order of magnitude less compared with the number of baryons, $`N_{\overline{B}}0.2N_B`$. According to Börner it is more than enough for the existence of the life in the Universe. Recall that the standard approach predicts for every $`10^9`$ antibaryons only $`(10^9+1)`$ baryons. As has been stated by Börner: It is to that one part in $`10^9`$ excess of ordinary matter that we owe our existence! .
For the early Universe the total number of baryons and antibaryons was roughly equal to the number of photons $`N_{\mathrm{ph}}`$ :
$`{\displaystyle \frac{N_{\overline{B}}+N_B}{N_{\mathrm{ph}}}}1.`$ (3.3)
Since the density of photons is equal to
$`{\displaystyle \frac{N_{\mathrm{ph}}}{V}}={\displaystyle \frac{2.404}{\pi ^2}}T^3,`$ (3.4)
where $`V`$ is the volume of the early Universe, and the density of the total number of baryons and antibaryons $`N_{\overline{B}}+N_B`$ calculated in our approach at $`T=175\mathrm{MeV}`$ amounts to
$`{\displaystyle \frac{N_{\overline{B}}+N_B}{V}}=C_B\times 0.442\times T^3,`$ (3.5)
we can estimate the numerical value of the parameter $`C_B`$:
$`C_B0.55\pm 0.08.`$ (3.6)
This gives the estimate of other input parameters $`C_{\overline{B}}`$ and $`C_M`$:
$`C_{\overline{B}}`$ $``$ $`0.25\pm 0.08,`$
$`C_M`$ $``$ $`0.13\pm 0.03.`$ (3.7)
The analysis of the influence of the input parameter $`C_{\overline{B}}/C_B=0.46\pm 0.15`$ on the evolution of the baryon–antibaryon asymmetry in the Universe from the early Universe up to the present epoch and the formation of a dark and strange matter in the Universe we are planning to carry out in our forthcoming publications.
Now we would like to discuss in more details our approach when compared with a simple coalescence model . The main distinction of our approach from a simple coalescence model is in the correlation between quarks and anti–quarks coalescing into hadrons. In fact, in a simple coalescence model quarks and anti–quarks are uncorrelated . This has allowed to introduce separately the number of light quarks $`q`$ and light anti–quarks $`\overline{q}`$ and the number of strange quarks $`s`$ and strange anti–quarks $`\overline{s}`$ . Moreover, this has turned out to be of use in order to hide the dependence of quark and anti–quark distribution functions on a temperature $`T`$ and a chemical potential $`\mu (T)`$ in the number of light quarks $`q`$ and light anti–quarks $`\overline{q}`$. A non–vanishing chemical potential of strange quarks $`\mu _s(T)`$ is hidden in the number of strange quarks $`s`$ and strange anti–quarks $`\overline{s}`$. Then, the calculation of multiplicities of hadrons produced from the QGP phase in a simple coalescence model resembles a quark counting. In fact, multiplicities of hadrons are proportional to products of the number of quarks $`(q,s)`$ and anti–quarks $`(\overline{q},\overline{s})`$ in accord the naive quark structure of hadrons. Since quarks and anti–quarks do not correlate, so that the multiplicities of hadrons turn out to be independent on the momenta of hadrons. Then, the numbers of quarks $`(q,s)`$ and anti–quarks $`(\overline{q},\overline{s})`$ and the coefficients of proportionality, the coalescence coefficients $`C_\mathrm{p}`$, $`C_\mathrm{\Lambda }`$, $`C_\mathrm{\Xi }`$, $`C_\mathrm{\Omega }`$ and $`C_{\overline{\mathrm{p}}}`$, $`C_{\overline{\mathrm{\Lambda }}}`$, $`C_{\overline{\mathrm{\Xi }}}`$, $`C_{\overline{\mathrm{\Omega }}}`$, are free parameters of a simple coalescence model. Therefore, a total number free parameters appearing in a simple coalescence model for the description of baryon and anti–baryon production is equal to twelve. By the assumption $`C_\mathrm{p}/C_{\overline{\mathrm{p}}}=C_\mathrm{\Lambda }/C_{\overline{\mathrm{\Lambda }}}=C_\mathrm{\Xi }/C_\mathrm{\Xi }=C_\mathrm{\Omega }/C_{\overline{\mathrm{\Omega }}}=1`$ the number free parameters has been reduced up to five $`(q,s,\overline{q},\overline{s},C)`$, where $`C`$ is a common for baryons and anti–baryons coalescence coefficient. Two of these free parameters have been fixed from the fit of experimental data on the baryon and anti–baryon production: $`\overline{q}/q=0.41\pm 0.02`$ and $`\overline{s}/s=0.75\pm 0.06`$ . Thus, there are three free parameters left in a simple coalescence model applied to the description of baryon and anti–baryon production from the QGP phase. It is also important to note that a simple coalescence model explains only multiplicities of baryon and anti–baryon production. In fact, save the ratio of multiplicities of the K<sup>+</sup> and K<sup>-</sup> mesons none other multiplicities of pseudoscalar and vector mesons have been predicted within a simple coalescence model . Therefore, it is not completely clear how many free parameters would be added in a simple coalescence model for description of multiplicities of pseudoscalar and vector meson production.
In our approach, where quarks and anti–quarks coalescing into hadrons are correlated, we have six free parameters $`T`$, $`\mu (T)`$, $`C_M`$, $`C_B`$, $`C_{\overline{B}}`$ and $`F_S`$. Five of these parameters $`T=175\mathrm{MeV}`$, $`\mu (T)`$ given by Eq.(1.4), $`F_S=3.5F_\pi =458.5\mathrm{MeV}`$, $`C_M/C_B=0.23\pm 0.04`$ and $`C_{\overline{B}}/C_B=0.46\pm 0.15`$ are fixed. Therefore, only there is one free parameter left in the approach. Thus, if to take into account that within our approach we have described not only multiplicities of baryon and anti–baryon production from the QGP but also multiplicities of pseudoscalar and vector meson production, all together 21 experimental ratios, our approach to the thermalized QGP looks much more preferable with respect to a simple coalescence model. If to remind that in our approach due to correlations between quarks and anti–quarks we are able to follow a dependence of multiplicities of hadron production on the hadronic momenta, so an advantage of our approach with respect to a simple coalescence model becomes obvious.
## Acknowledgement
One of the authors (A.N. Ivanov) thanks Prof. T. S. Bir$`\stackrel{´}{\mathrm{o}}`$, Acad. J. Zim$`\stackrel{´}{\mathrm{a}}`$nyi and the members of the Quark–gluon plasma group of Theory Division of Research Institute of Particle and Nuclear Physics of Hungarian Academy of Sciences for helpful discussions during the seminar where this paper has been reported. He is also grateful to Prof. T. S. Bir$`\stackrel{´}{\mathrm{o}}`$ and Prof. V. Gogohia for warm hospitality extended to him during his visit to Budapest.
We are greatful to Prof. A. Rebhan for fruitful discussions. Discussions with Prof. I. N .Toptygin of the cosmological aspects of our approach are appreciated.
Table 1. The theoretical ratios of multiplicities of the hadron production are compared with the experimental data obtained by NA44, NA49, NA50 and WA97 Collaborations on Pb+Pb collisions at 158 GeV/nucleon, NA35 Collaboration on S+S collisions and NA38 Collaboration on O+U and S+U collisions at 200 GeV/nucleon. The theoretical multiplicities are calculated at the temperature $`T=175\mathrm{MeV}`$.
| N | Ratio | Model | Data | Coll. | Rapidity | Ref. |
| --- | --- | --- | --- | --- | --- | --- |
| 1 | $`\overline{p}/p`$ | 0.097(32) | 0.055(10) | NA44 | 2.3–2.9 | |
| | | 0.097(32) | 0.085(8) | NA49 | 2.5–3.3 | |
| 2 | $`\overline{\mathrm{\Lambda }}/\mathrm{\Lambda }`$ | 0.168(55) | 0.128(12) | WA97 | 2.4–3.4 | |
| 3 | $`\overline{\mathrm{\Xi }}/\mathrm{\Xi }`$ | 0.270(88) | 0.227(33) | NA49 | 3.1–3.85 | |
| | | 0.270(88) | 0.266(28) | WA97 | 2.4–3.4 | |
| 4 | $`\overline{\mathrm{\Omega }}/\mathrm{\Omega }`$ | fit | 0.46(15) | WA97 | 2.4–3.4 | |
| 5 | $`\mathrm{\Xi }/\mathrm{\Lambda }`$ | 0.108 | 0.127(11) | NA49 | 3.1–3.85 | |
| | | 0.108 | 0.093(7) | WA97 | 2.4–3.4 | |
| 6 | $`\mathrm{\Omega }/\mathrm{\Xi }`$ | 0.166 | 0.195(28) | WA97 | 2.4–3.4 | |
| 7 | $`\overline{\mathrm{\Xi }}/\overline{\mathrm{\Lambda }}`$ | 0.173 | 0.180(39) | NA49 | 3.1–3.85 | |
| | | 0.173 | 0.195(23) | WA97 | 2.4–3.4 | |
| 8 | $`\overline{\mathrm{\Lambda }}/\overline{p}`$ | 2.081 | 3(1) | NA49 | 3.1–3.85 | |
| 9 | $`\overline{\mathrm{\Omega }}/\overline{\mathrm{\Xi }}`$ | 0.282 | 0.27(6) | WA97 | 2.4–3.4 | |
| 10 | $`K^+/K^{}`$ | 1.520 | 1.85(9) | NA44 | 2.4–3.5 | |
| | | 1.520 | 1.8(1) | NA49 | all | |
| 11 | $`K^+/\pi ^+`$ | 0.139 | 0.137(8) | NA35 | all | |
| 12 | $`K^{}/\pi ^{}`$ | 0.090 | 0.076(5) | NA35 | all | |
| 13 | $`K_S^0/\pi ^{}`$ | 0.113 | 0.125(19) | NA49 | all | |
| 14 | $`\eta /\pi ^0`$ | 0.088 | 0.081(13) | WA98 | 2.3–2.9 | |
| 15 | $`2\varphi /(\pi ^++\pi ^{})`$ | $`7.8\times 10^3`$ | $`9.1(1.0)\times 10^3`$ | NA50 | 2.9–3.9 | |
| 16 | $`\varphi /(\rho ^0+\omega ^0)`$ | 0.103 | $`0.1`$ | NA38 | 2.8–4.1 | |
| 17 | $`\varphi /K_S^0`$ | 0.071 | 0.084(11) | NA49 | all | |
| 18 | $`\mathrm{\Lambda }/K_S^0`$ | fit | 0.65(11) | WA97 | 2.4–3.4 | |
| 19 | $`p/K^+`$ | 0.136(23) | | | | |
| 20 | $`\overline{p}/K^{}`$ | 0.020(7) | | | | |
| 21 | $`\overline{p}/pK^+/K^{}`$ | 0.147(57) | 0.102(19) | NA44 | 2.3–2.9 | |
| | | 0.147(57) | 0.153(17) | NA49 | 2.5–3.3 | |
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# 1 Introduction
## 1 Introduction
The study of Kaluza-Klein sphere reductions of supergravities has so far concentrated mostly on the examples where the theories admit vacuum solutions of the form AdS$`\times `$Sphere, which are the near-horizon structures of certain $`p`$-brane solutions of the theories.<sup>1</sup><sup>1</sup>1Throughout this paper we are concerned only with those “remarkable” Kaluza-Klein sphere reductions for which no known group-theoretic argument guaranteeing the consistency of the reduction exists. Consistent reductions on $`S^3`$, or indeed any group manifold $`G`$, can always be performed in the case where one truncates to the sector of singlets under the right action of the group $`G`$, but the consistency in such a case is guaranteed, and therefore is not of interest to us in the present paper. These include 11-dimensional supergravity, which has AdS$`{}_{4}{}^{}\times S^7`$ and AdS$`{}_{7}{}^{}\times S^4`$ vacuum solutions, and type IIB supergravity, which has an AdS$`{}_{5}{}^{}\times S^5`$ vacuum solution. The $`S^7`$ and $`S^4`$ reduction Anätze for 11-dimensional supergravity were presented in . For type IIB supergravity the $`S^5`$ reduction of its $`SL(2,\text{I}\mathrm{R})`$-singlet subsector, which gives rise to five-dimensional fields including the entire set of $`SO(6)`$ gauge bosons, was given in . Explicit reduction Ansätze for various subsectors of these supergravity reductions can be found in .
In general, vacuum supergravity solutions with non-trivial field-strength fluxes are of the form of warped products of a certain spacetime geometry and internal spheres. The consistency of sphere reductions in such cases have been much less fully studied. The first example of this type was the consistent warped $`S^4`$ reduction of massive type IIA supergravity, to give rise to the massive $`SU(2)`$-gauged supergravity in $`D=6`$. The vacuum AdS<sub>6</sub> solution can be viewed as the near-horizon structure of an intersecting D4-D8 brane .
Further examples of consistent sphere reductions were obtained in , where the resulting theories admit “vacuum solutions” that are domain walls rather than AdS spacetimes. In , a necessary condition for the consistency of a sphere reduction of a theory was given. Namely, if a theory can be consistently reduced on $`S^n`$, with a massless truncation that retains all the $`SO(n+1)`$ Yang-Mills gauge fields, then a necessary requirement is that if instead a toroidal reduction on $`T^n`$ is performed, this must give rise to the same content of massless fields.<sup>2</sup><sup>2</sup>2This is because by turning off the gauge-coupling parameter $`g`$, by sending the radius of the $`n`$-sphere to infinity, we must recover the same massless field content as would result from a flat (toroidal) reduction. Furthermore, the $`T^n`$-reduced theory must have at least an $`SO(n+1)`$ global symmetry, with sufficiently many abelian vector fields to supply at least those of the adjoint representation of $`SO(n+1)`$. These conditions are very restrictive, and only in limited cases can a consistent sphere reduction that retains all the Yang-Mills fields occur.
In the type IIA and type IIB theories, the NS-NS branes and D-branes have near-horizon structures of the form (Domain wall)$`\times S^n`$, for various values of $`n`$. It is easy to verify in each case that if a Kaluza-Klein reduction of the theory on $`T^n`$ is performed, this yields a scalar coset $`SL(n+1,\text{I}\mathrm{R})/SO(n+1)`$ in the lower dimension, since this theory can instead be obtained from a $`T^{n+1}`$ reduction from $`D=11`$. In particular, therefore, this theory has an $`SO(n+1)`$ global symmetry subgroup. It was noted that there exists a “dual frame” for each D-brane, in which the near-horizon structure of the brane generically becomes AdS$`{}_{10n}{}^{}\times S^n`$, while instead it becomes Minkowski$`\times S^3`$ when $`n=3`$. This leads to the conjectured Domain Wall/QFT correspondence , generalising the notion of the AdS/CFT correspondence . This correspondence was generalised to lower dimensions in .
Further evidence for Domain Wall/QFT correspondence was obtained in , where it was shown that it is consistent to reduce to the subsector scalars associated with the Cartan generators in the scalar coset $`SL(n+1,\text{I}\mathrm{R})/SO(n+1)`$. The multi-parameter domain-wall solutions of these lower-dimensional theories can then be lifted back to the higher dimension, where they correspond to certain ellipsoidal distributions of the $`p`$-branes, thus implying that these domain wall geometries correspond to the Coulomb branch of the quantum field theory, and generalising the results on the Coulomb branch in the AdS/CFT correspondence, discussed in . Interestingly, the wave equations for minimally-coupled scalar fluctuations in the lower-dimensional domain-wall backgrounds depend only on the dimension of the internal sphere used in the reduction .
The consistency of the Kaluza-Klein reduction in this Cartan subsector of scalar fields leads one to believe further that it is consistent to reduce the type IIA and type IIB theories on the relevant $`n`$-spheres, while retaining all the massless fields. For non-trivial vacuum NS-NS flux, both the type IIA and type IIB theories can be expected to be consistently reducible on $`S^3`$ and on $`S^7`$, and indeed, for $`N=1`$ supergravity, the consistency has been demonstrated, and the corresponding gauged supergravities in $`D=7`$ and $`D=3`$ were obtained in . For non-trivial vacuum R-R flux, we expect that it is consistent to reduce the type IIA theory on $`S^n`$ with $`n=2,4,6,8`$ and the type IIB theory on $`S^n`$ with $`n=1,3,5,7`$. Indeed, the $`S^4`$ and $`S^7`$ reductions of type IIA, where $`S^4`$ and $`S^7`$ are associated with the D4-brane and NS-NS string, can be established from the $`S^1`$ reduction of the corresponding $`S^4`$ or $`S^7`$ reduction of eleven-dimensional supergravity. In section 5, we carry this out explicitly for the $`S^4`$ reduction of the type IIA theory.
First, we demonstrate explicitly that it is consistent to perform an $`S^3`$ reduction of the type IIA theory, while retaining all the massless fields, including in particular the entire set of $`SO(4)`$ Yang-Mills gauge fields. This case is of particular interest because $`S^3`$ is itself the group manifold $`SU(2)`$, and strings propagating in group-manifold backgrounds have been extensively studied in the past. It should be emphasised though that typically when Kaluza-Klein reductions on a group manifold $`G`$ have been discussed in the literature, a truncation is performed in which only those fields that are singlets under the left action $`G_L`$ of the $`G_L\times G_R`$ isometry group are retained. For example, the $`S^3`$ reduction of $`N=1`$ supergravity in $`D=10`$, retaining only one $`SU(2)`$ Yang-Mills fields, was performed to give rise to gauged simple supergravity in $`D=7`$ with a domain wall vacuum solution . Such a truncation guarantees that a consistent reduction can be performed, but it fails to exploit the much more remarkable fact that in this $`S^3`$ case a reduction that retains all the $`SO(4)SU(2)_L\times SU(2)_R`$ Yang-Mills fields, and not merely those of $`SU(2)_L`$, is possible. A further reason for wishing to include all the gauge fields of $`SO(4)`$ is that only then do we obtain a seven-dimensional theory with maximal supersymmetry.
We obtain the consistent $`S^3`$ reduction of type IIA supergravity by taking a singular limit of the $`S^4`$ reduction of eleven-dimensional supergravity, in which the $`S^4`$ degenerates to $`\text{I}\mathrm{R}\times S^3`$. In order to do this, we begin in section 2 by reviewing the $`S^4`$ reduction from $`D=11`$, first obtained in . By substituting this into the eleven-dimensional Bianchi identity and equation of motion for the 4-form, we obtain complete and explicit seven-dimensional equations of motion, and the Lagrangian that generates them. We also discuss the “ungauging limit” in which the radius of the 4-sphere is sent to infinity. (An analogous limit was also considered in , in the context of the $`U(1)^2`$ subgroup of $`SO(4)`$.) In particular, we clarify certain aspects of this limiting process, showing that the limit is smooth in the seven-dimensional equations of motion, but pathological at the level of the gauged-supergravity Lagrangian.
In section 3 we take a different singular limit of the seven-dimensional $`SO(5)`$-gauged supergravity, in which an $`SO(4)`$ gauging remains. Again, this is a smooth limit of the equations of motion, but not of the gauged supergravity Lagrangian. In section 4 we apply this limiting procedure to the $`S^4`$ Kaluza-Klein reduction Ansatz of eleven-dimensional supergravity, showing that it corresponds to a degeneration of the 4-sphere to $`\text{I}\mathrm{R}\times S^3`$. The reduction can then be viewed as an initial reduction to give type IIA supergravity in $`D=10`$, followed by a reduction on $`S^3`$. By this means, we arrive at the consistent $`S^3`$ reduction Ansatz for type IIA supergravity.
In section 5 we construct instead the consistent $`S^4`$ Kaluza-Klein reduction of type IIA supergravity. This can again be obtained from the starting point of the $`S^4`$ reduction of the eleven-dimensional theory. In this case we do not need to take any singular limit in the internal directions, but rather, we perform a standard Kaluza-Klein $`S^1`$ reduction of the original seven-dimensional theory coming from $`D=11`$, and show how this can be reinterpreted as an $`S^4`$ reduction of type IIA supergravity. The paper ends with concluding remarks in section 6.
## 2 The $`S^4`$ reduction of eleven-dimensional supergravity
### 2.1 Metric and 4-form Ansatz
The complete Ansatz for the $`S^4`$ reduction of eleven-dimensional supergravity was obtained in , using a formalism based on an analysis of the supersymmetry transformation rules. One may also study the reduction from a purely bosonic standpoint, by verifying that if the Ansatz is substituted into the eleven-dimensional equations of motion, it consistently yields the equations of motion of the seven-dimensional gauged $`SO(5)`$ supergravity. We shall carry out this procedure here, in order to establish notation, and to obtain the complete system of seven-dimensional bosonic equations of motion, which we shall need in the later part of the paper.
After some manipulation, the Kaluza-Klein $`S^4`$ reduction Ansatz obtained in for eleven-dimensional supergravity can be expressed as follows:
$`d\widehat{s}_{11}^2`$ $`=`$ $`\mathrm{\Delta }^{1/3}ds_7^2+{\displaystyle \frac{1}{g^2}}\mathrm{\Delta }^{2/3}T_{ij}^1D\mu ^iD\mu ^j,`$ (1)
$`\widehat{F}_{\left(4\right)}`$ $`=`$ $`{\displaystyle \frac{1}{4!}}ϵ_{i_1\mathrm{}i_5}[{\displaystyle \frac{1}{g^3}}U\mathrm{\Delta }^2\mu ^{i_1}D\mu ^{i_2}\mathrm{}D\mu ^{i_5}`$ (2)
$`+{\displaystyle \frac{4}{g^3}}\mathrm{\Delta }^2T^{i_1m}DT^{i_2n}\mu ^m\mu ^nD\mu ^{i_3}\mathrm{}D\mu ^{i_5}`$
$`+{\displaystyle \frac{6}{g^2}}\mathrm{\Delta }^1F_{\left(2\right)}^{i_1i_2}D\mu ^{i_3}D\mu ^{i_4}T^{i_5j}\mu ^j]T_{ij}S_{\left(3\right)}^i\mu ^j+{\displaystyle \frac{1}{g}}S_{\left(3\right)}^iD\mu ^i,`$
where
$`U2T_{ij}T_{jk}\mu ^i\mu ^k\mathrm{\Delta }T_{ii},\mathrm{\Delta }T_{ij}\mu ^i\mu ^j,`$
$`F_{\left(2\right)}^{ij}dA_{\left(1\right)}^{ij}+gA_{\left(1\right)}^{ik}A_{\left(1\right)}^{kj},D\mu ^id\mu ^i+gA_{\left(1\right)}^{ij}\mu ^j,`$
$`DT_{ij}dT_{ij}+gA_{\left(1\right)}^{ik}T_{kj}+gA_{\left(1\right)}^{jk}T_{ik},\mu ^i\mu ^i1,`$ (3)
where the symmetric matrix $`T_{ij}`$, which parameterises the scalar coset $`SL(6,\text{I}\mathrm{R})/SO(6)`$, is unimodular.
### 2.2 Derivation of the seven-dimensional equations of motion
Consider first the Bianchi identity $`d\widehat{F}_{\left(4\right)}=0`$. Substituting (2) into this, we obtain the following seven-dimensional equations:
$`D(T_{ij}S_{\left(3\right)}^j)`$ $`=`$ $`F_{\left(2\right)}^{ij}S_{\left(3\right)}^j,`$ (4)
$`H_{\left(4\right)}^i`$ $`=`$ $`gT_{ij}S_{\left(3\right)}^j+{\displaystyle \frac{1}{8}}ϵ_{ij_1\mathrm{}j_4}F_{\left(2\right)}^{j_1j_2}F_{\left(2\right)}^{j_3j_4},`$ (5)
where we define
$$H_{\left(4\right)}^iDS_{\left(3\right)}^i=dS_{\left(3\right)}^i+gA_{\left(1\right)}^{ij}S_{\left(3\right)}^j.$$
(6)
Next, we substitute the Ansatz into the $`D=11`$ field equation $`d\widehat{}\widehat{F}_{\left(4\right)}=\frac{1}{2}\widehat{F}_{\left(4\right)}\widehat{F}_{\left(4\right)}`$. To do this, we need the eleven-dimensional Hodge dual $`\widehat{}\widehat{F}_{\left(4\right)}`$, which we find is given by
$`\widehat{}\widehat{F}_{\left(4\right)}`$ $`=`$ $`gUϵ_{\left(7\right)}{\displaystyle \frac{1}{g}}T_{ij}^1DT^{ik}\mu _kD\mu ^j+{\displaystyle \frac{1}{2g^2}}T_{ik}^1T_j\mathrm{}^1F_{\left(2\right)}^{ij}D\mu ^kD\mu ^{\mathrm{}}`$
$`+{\displaystyle \frac{1}{g^4}}\mathrm{\Delta }^1T_{ij}S_{\left(3\right)}^i\mu ^jW{\displaystyle \frac{1}{6g^3}}\mathrm{\Delta }^1ϵ_{ij\mathrm{}_1\mathrm{}_2\mathrm{}_3}S_{\left(3\right)}^mT_{im}T_{jk}\mu ^kD\mu ^\mathrm{}_1D\mu ^\mathrm{}_2D\mu ^\mathrm{}_3,`$
where
$$W\frac{1}{24}ϵ_{i_1\mathrm{}i_5}\mu ^{i_1}D\mu ^{i_2}\mathrm{}D\mu ^{i_5}.$$
(8)
The field equation for $`\widehat{F}_{\left(4\right)}`$ then implies
$`D\left(T_{ik}^1T_j\mathrm{}^1F_{\left(2\right)}^{ij}\right)`$ $`=`$ $`2gT_{i[k}^1DT_{\mathrm{}]i}{\displaystyle \frac{1}{2g}}ϵ_{i_1i_2i_3k\mathrm{}}F_2^{i_1i_2}H_{\left(4\right)}^{i_3}`$ (9)
$`+{\displaystyle \frac{3}{2g}}\delta _{i_1i_2k\mathrm{}}^{j_1j_2j_3j_4}F_{\left(2\right)}^{i_1i_2}F_{\left(2\right)}^{j_1j_2}F_{\left(2\right)}^{j_3j_4}S_{\left(3\right)}^kS_{\left(3\right)}^{\mathrm{}}.`$
$`D\left(T_{ik}^1D(T_{kj})\right)`$ $`=`$ $`2g^2(2T_{ik}T_{kj}T_{kk}T_{ij})ϵ_{\left(7\right)}+T_{im}^1T_k\mathrm{}^1F_{\left(2\right)}^m\mathrm{}F_{\left(2\right)}^{kj}`$ (10)
$`+T_{jk}S_{\left(3\right)}^kS_{\left(3\right)}^i\frac{1}{5}\delta _{ij}[2g^2(2T_{ik}T_{ik}2(T_{ii})^2)ϵ_{\left(7\right)}`$
$`+T_{nm}^1T_k\mathrm{}^1F_{\left(2\right)}^m\mathrm{}F_{\left(2\right)}^{kn}+T_k\mathrm{}S_{\left(3\right)}^kS_{\left(3\right)}^{\mathrm{}}],`$
for the Yang-Mills and scalar equations of motion in $`D=7`$.<sup>3</sup><sup>3</sup>3Note from (9) that it would be inconsistent to set the Yang-Mills fields to zero while retaining the scalars $`T_{ij}`$, since the currents $`T_{i[k}^1DT_{\mathrm{}]i}`$ act as sources for them. A truncation where the Yang-Mills fields are set to zero is consistent, however, if the scalars are also truncated to the diagonal subsector $`T_{ij}=\mathrm{diag}(X_1,X_2,\mathrm{},X_6)`$, as in the consistent reductions constructed in .
We find that all the equations of motion can be derived from the following seven-dimensional Lagrangian
$`_7`$ $`=`$ $`R\text{1}\mathrm{l}\frac{1}{4}T_{ij}^1DT_{jk}T_k\mathrm{}^1DT_\mathrm{}i\frac{1}{4}T_{ik}^1T_j\mathrm{}^1F_{\left(2\right)}^{ij}F_{\left(2\right)}^k\mathrm{}\frac{1}{2}T_{ij}S_{\left(3\right)}^iS_{\left(3\right)}^j`$ (11)
$`+{\displaystyle \frac{1}{2g}}S_{\left(3\right)}^iH_{\left(4\right)}^i{\displaystyle \frac{1}{8g}}ϵ_{ij_1\mathrm{}j_4}S_{\left(3\right)}^iF_{\left(2\right)}^{j_1j_2}F_{\left(2\right)}^{j_3j_4}+{\displaystyle \frac{1}{g}}\mathrm{\Omega }_{\left(7\right)}V\text{1}\mathrm{l},`$
where $`H_{\left(4\right)}^i`$ are given by (6) and the potential $`V`$ is given by
$$V=\frac{1}{2}g^2\left(2T_{ij}T_{ij}(T_{ii})^2\right),$$
(12)
and $`\mathrm{\Omega }_{\left(7\right)}`$ is a Chern-Simons type of term built from the Yang-Mills fields, which has the property that its variation with respect to $`A_{\left(1\right)}^{ij}`$ gives
$$\delta \mathrm{\Omega }_{\left(7\right)}=\frac{3}{4}\delta _{i_1i_2k\mathrm{}}^{j_1j_2j_3j_4}F_{\left(2\right)}^{i_1i_2}F_{\left(2\right)}^{j_1j_2}F_{\left(2\right)}^{j_3j_4}\delta A_{\left(1\right)}^k\mathrm{}.$$
(13)
Note that the $`S_{\left(3\right)}^i`$ are viewed as fundamental fields in the Lagrangian, and that (5) is their first-order equation. In fact (11) is precisely the bosonic sector of the Lagrangian describing maximal gauged seven-dimensional supergravity that was derived in . An explicit expression for the 7-form $`\mathrm{\Omega }_{\left(7\right)}`$ can be found there.
Although we have fully checked the eleven-dimensional Bianchi identity and field equation for $`\widehat{F}_{\left(4\right)}`$ here, we have not completed the task of substituting the Ansatz into the eleven-dimensional Einstein equations. This would be an extremely complicated calculation, on account of the Yang-Mills gauge fields. However, various complete consistency checks, including the higher-dimensional Einstein equation, have been performed in various truncations of the full $`N=4`$ maximal supergravity embedding, including the $`N=2`$ gauged theory in , and the non-supersymmetric truncation in where the gauge fields are set to zero and only the diagonal scalars in $`T_{ij}`$ are retained. All the evidence points to the full consistency of the reduction.<sup>4</sup><sup>4</sup>4The original demonstration in , based on the reduction of the eleven-dimensional supersymmetry transformation rules, also provides extremely compelling evidence. Strictly speaking the arguments presented there also fall short of a complete and rigorous proof, since they involve an approximation in which the quartic fermion terms in the theory are neglected.
It is perhaps worth making a few further remarks on the nature of the reduction Ansatz. One might wonder whether the Ansatz (2) on the 4-form field strength $`\widehat{F}_{\left(4\right)}`$ could be re-expressed as an Ansatz on its potential $`\widehat{A}_{\left(3\right)}`$. As it stands, (2) only satisfies the Bianchi identity $`d\widehat{F}_{\left(4\right)}=0`$ by virtue of the lower-dimensional equations (4) and (5). However, if (5) is substituted into (2), we obtain an expression that satisfies $`d\widehat{F}_{\left(4\right)}=0`$ without the use of any lower-dimensional equations. However, one does still have to make use of the fact that the $`\mu ^i`$ coordinates satisfy the constraint $`\mu ^i\mu ^i=1`$, and this prevents one from writing an explicit Ansatz for $`\widehat{A}_{\left(3\right)}`$ that has a manifest $`SO(5)`$ symmetry. One could solve for one of the $`\mu ^i`$ in terms of the others, but this would break the manifest local symmetry from $`SO(5)`$ to $`SO(4)`$. In principle though, this could be done, and then one could presumably substitute the resulting Ansatz directly into the eleven-dimensional Lagrangian. After integrating out the internal 4-sphere directions, one could then in principle obtain a seven-dimensional Lagrangian in which, after re-organising terms, the local $`SO(5)`$ symmetry could again become manifest.
It should, of course, be emphasised that merely substituting an Ansatz into a Lagrangian and integrating out the internal directions to obtain a lower-dimensional Lagrangian is justifiable only if one already has an independent proof of the consistency of the proposed reduction Ansatz.<sup>5</sup><sup>5</sup>5A classic illustration is provided by the example of 5-dimensional pure gravity with an (inconsistent) Kaluza-Klein reduction in which the scalar dilaton is omitted. Substituting this into the 5-dimensional Einstein-Hilbert action yields the perfectly self-consistent Einstein-Maxwell action in $`D=4`$, but fails to reveal that setting the scalar to zero is inconsistent with the internal component of the 5-dimensional Einstein equation. If one is in any case going to work with the higher-dimensional field equations in order to prove the consistency, it is not clear that there would be any significant benefit to be derived from then re-expressing the Ansatz in a form where it could be substituted into the Lagrangian.
### 2.3 Ungauging: the $`g0`$ limit
It is interesting to observe that one cannot take the limit $`g0`$ in the Lagrangian (11), on account of the terms proportional to $`g^1`$ in the second line. We know, on the other hand, that it must be possible to recover the ungauged $`D=7`$ theory by turning off the gauge coupling constant. In fact the problem is associated with a pathology in taking the limit at the level of the Lagrangian, rather than in the equations of motion. This can be seen by looking instead at the seven-dimensional equations of motion, which were given earlier. The only apparent obstacle to taking the limit $`g0`$ is in the Yang-Mills equations (9), but in fact this illusory. If we substitute the first-order equation (5) into (9) it gives
$$D\left(T_{ik}^1T_j\mathrm{}^1F_{\left(2\right)}^{ij}\right)=2gT_{i[k}^1DT_{\mathrm{}]i}\frac{1}{2}ϵ_{i_1i_2i_3k\mathrm{}}F_2^{i_1i_2}T_{ij}S_{\left(3\right)}^jS_{\left(3\right)}^kS_{\left(3\right)}^{\mathrm{}},$$
(14)
which has a perfectly smooth $`g0`$ limit. It is clear that equations of motion (5) and (10) and the Einstein equations of motion also have a smooth limit. (The reason why the Einstein equations have the smooth limit is because the $`1/g`$ terms in the Lagrangian (11) do not involve the metric, and thus they give no contribution.)
Unlike in the gauged theory, we should not treat the $`S_{\left(3\right)}^i`$ fields as fundamental variables in a Lagrangian formulation in the ungauged limit. This is because once the gauge coupling $`g`$ is sent to zero, the fields $`S_{\left(3\right)}^i`$ behave like 3-form field strengths. This can be seen from the first-order equation of motion (5), which in the limit $`g0`$ becomes
$$dS_{\left(3\right)}^i=\frac{1}{8}ϵ_{ij_1\mathrm{}j_4}dA_{\left(1\right)}^{j_1j_2}dA_{\left(1\right)}^{j_3j_4},$$
(15)
and should now be interpreted as a Bianchi identity. This can be solved by introducing 2-form gauge potentials $`A_{\left(2\right)}^i`$, with the $`S_{\left(3\right)}^i`$ given by
$$S_{\left(3\right)}^i=dA_{\left(2\right)}^i+\frac{1}{8}ϵ_{ij_1\mathrm{}j_4}A_{\left(1\right)}^{j_1j_2}dA_{\left(1\right)}^{j_3j_4}.$$
(16)
In terms of these 2-form potentials, the equations of motion can now be obtained from the Lagrangian
$`_7^0`$ $`=`$ $`R\text{1}\mathrm{l}\frac{1}{4}T_{ij}^1dT_{jk}T_k\mathrm{}^1dT_\mathrm{}i\frac{1}{4}T_{ik}^1T_j\mathrm{}^1F_{\left(2\right)}^{ij}F_{\left(2\right)}^k\mathrm{}\frac{1}{2}T_{ij}S_{\left(3\right)}^iS_{\left(3\right)}^j`$ (17)
$`+\frac{1}{2}A_{\left(1\right)}^{ij}S_{\left(3\right)}^iS_{\left(3\right)}^j2S_{\left(3\right)}^iA_{\left(2\right)}^jdA_{\left(1\right)}^{ij},`$
where $`S_{\left(3\right)}^i`$ is given by (16). This is precisely the bosonic Lagrangian of the ungauged maximal supergravity in $`D=7`$.
It is worth exploring in a little more detail why it is possible to take a smooth $`g0`$ limit in the seven-dimensional equations of motion, but not in the Lagrangian. We note that in this limit the Lagrangian (11) can be expressed as
$$_7=\frac{1}{g}L+𝒪(1),$$
(18)
where
$$L=\frac{1}{2}S_{\left(3\right)}^idS_{\left(3\right)}^i\frac{1}{8}ϵ_{ij_1\mathrm{}j_4}S_{\left(3\right)}^iF_{\left(2\right)}^{j_1j_2}F_{\left(2\right)}^{j_3j_4}+\mathrm{\Omega }_{\left(7\right)}.$$
(19)
The term $`L/g`$, which diverges in the $`g0`$ limit, clearly emphasises that the Lagrangian (17) is not merely the $`g0`$ limit of (11). However if we make use of the equations of motion, we find that in the $`g0`$ limit the $`S_{\left(3\right)}^i`$ can be solved by (16). Substituting this into (19), we find that in this limit it becomes
$$L=\frac{1}{16}ϵ_{ij_1\mathrm{}j_4}dA_{\left(2\right)}^idA_{\left(1\right)}^{j_1j_2}dA_{\left(1\right)}^{j_3j_4}+𝒪(g),$$
(20)
and so the singular terms in $`L/g`$ form a total derivative and hence can be subtracted from the Lagrangian. This analysis explains why it is possible to take a smooth $`g0`$ limit in the equations of motion, but not in the Lagrangian.
## 3 The gauged $`SO(4)`$ limit of maximal $`D=7`$ supergravity
Here we examine, at the level of the seven-dimensional theory itself, how to take a limit in which the $`SO(5)`$ gauged sector is broken down to $`SO(4)`$. In a later section, we shall show how this can be interpreted as an $`S^3`$ reduction of type IIA supergravity. We shall do that by showing how to take a limit in which the internal $`S^4`$ in the original reduction from $`D=11`$ becomes $`\text{I}\mathrm{R}\times S^3`$. For now, however, we shall examine the $`SO(4)`$-gauged limit entirely from the perspective of the seven-dimensional theory itself.
To take the limit, we break the $`SO(5)`$ covariance by splitting the $`\underset{¯}{5}`$ index $`i`$ as
$$i=(0,\alpha ),$$
(21)
where $`1\alpha 4`$. We also introduce a constant parameter $`\lambda `$, which will be sent to zero as the limit is taken. We find that the various seven-dimensional fields, and the $`SO(5)`$ gauge-coupling constant, should be scaled as follows:
$`g=\lambda ^2\stackrel{~}{g},A_{\left(1\right)}^{0\alpha }=\lambda ^3\stackrel{~}{A}_{\left(1\right)}^{0\alpha },A_{\left(1\right)}^{\alpha \beta }=\lambda ^2\stackrel{~}{A}_{\left(1\right)}^{\alpha \beta },`$
$`S_{\left(3\right)}^0=\lambda ^4\stackrel{~}{S}_{\left(3\right)}^0,S_{\left(3\right)}^\alpha =\lambda \stackrel{~}{S}_{\left(3\right)}^\alpha .`$ (22)
$`T_{ij}^1=\left(\begin{array}{cc}\lambda ^8\mathrm{\Phi }& \lambda ^3\mathrm{\Phi }\chi _\alpha \\ \lambda ^3\mathrm{\Phi }\chi _\alpha & \lambda ^2M_{\alpha \beta }^1+\lambda ^2\mathrm{\Phi }\chi _\alpha \chi _\beta \end{array}\right).`$
As we show in the next section, this rescaling corresponds to a degeneration of $`S^4`$ to $`R\times S^3`$. Note that in this rescaling, we have also performed a decomposition of the scalar matrix $`T_{ij}^1`$ that is of the form of a Kaluza-Klein metric decomposition. It is useful also to present the consequent decomposition for $`T_{ij}`$, which turns out to be
$$T_{ij}=\left(\begin{array}{cc}\lambda ^8\mathrm{\Phi }^1+\lambda ^8\chi _\gamma \chi ^\gamma & \lambda ^3\chi ^\alpha \\ \lambda ^3\chi ^\alpha & \lambda ^2M_{\alpha \beta }\end{array}\right).$$
(23)
Calculating the determinant, we get
$$det(T_{ij})=\mathrm{\Phi }^1det(M_{\alpha \beta }).$$
(24)
Since we know that $`det(T_{ij})=1`$, it follows that
$$\mathrm{\Phi }=det(M_{\alpha \beta }).$$
(25)
The fields $`\chi _\alpha `$ are “axionic” scalars. Note that we shall also have
$`H_{\left(4\right)}^0=\lambda ^4\stackrel{~}{H}_{\left(4\right)}^0,H_{\left(4\right)}^\alpha =\lambda \stackrel{~}{H}_{\left(4\right)}^\alpha ,`$
$`\stackrel{~}{H}_{\left(4\right)}^0=d\stackrel{~}{S}_{\left(3\right)}^0,\stackrel{~}{H}_{\left(4\right)}^\alpha =\stackrel{~}{D}\stackrel{~}{S}_{\left(3\right)}^\alpha \stackrel{~}{g}\stackrel{~}{A}_{\left(1\right)}^{0\alpha }\stackrel{~}{S}_{\left(3\right)}^0.`$ (26)
We have defined an $`SO(4)`$-covariant exterior derivative $`\stackrel{~}{D}`$, which acts on quantities with $`SO(4)`$ indices $`\alpha ,\beta ,\mathrm{}`$ in the obvious way:
$$\stackrel{~}{D}X_\alpha =dX_\alpha +\stackrel{~}{g}\stackrel{~}{A}_{\left(1\right)}^{a\beta }X_\beta ,$$
(27)
etc. It is helpful also to make the following further field redefinitions:
$`G_{\left(2\right)}^\alpha `$ $``$ $`\stackrel{~}{F}_{\left(2\right)}^{0\alpha }+\chi _\beta \stackrel{~}{F}_{\left(2\right)}^{\beta \alpha },`$
$`G_{\left(3\right)}^\alpha `$ $``$ $`\stackrel{~}{S}_{\left(3\right)}^\alpha \chi _\alpha \stackrel{~}{S}_{\left(3\right)}^0,`$ (28)
$`G_{\left(1\right)}^\alpha `$ $``$ $`\stackrel{~}{D}\chi _\alpha \stackrel{~}{g}\stackrel{~}{A}_{\left(1\right)}^{0\alpha },`$
where $`\stackrel{~}{F}_{\left(2\right)}^{0\alpha }\stackrel{~}{D}\stackrel{~}{A}_{\left(1\right)}^{0\alpha }`$.
We may now subsititute these redefined fields into the seven-dimensional equations of motion. We find that a smooth limit in which $`\lambda `$ is sent to zero exists, leading to an $`SO(4)`$-gauged seven-dimensional theory. Our results for the seven-dimensional equations of motion are as follows. The fields $`H_{\left(4\right)}^i`$ become
$$\stackrel{~}{H}_{\left(4\right)}^0=d\stackrel{~}{S}_{\left(3\right)}^0,\stackrel{~}{H}_{\left(4\right)}^\alpha =\stackrel{~}{D}G_{\left(3\right)}^\alpha +G_{\left(1\right)}^\alpha \stackrel{~}{S}_{\left(3\right)}^0+\chi _\alpha dS_{\left(3\right)}^0.$$
(29)
The first-order equations (6) give
$`\stackrel{~}{H}_{\left(4\right)}^0`$ $`=`$ $`\frac{1}{8}ϵ_{\alpha _1\mathrm{}\alpha _4}\stackrel{~}{F}_{\left(2\right)}^{\alpha _1\alpha _2}\stackrel{~}{F}_{\left(2\right)}^{\alpha _3\alpha _4},`$
$`\stackrel{~}{F}_{\left(4\right)}^\alpha `$ $`=`$ $`\stackrel{~}{g}M_{\alpha \beta }G_{\left(3\right)}^\beta \frac{1}{2}ϵ_{\alpha \beta \gamma \delta }G_{\left(2\right)}^\beta \stackrel{~}{F}_{\left(2\right)}^{\gamma \delta }G_{\left(1\right)}^\alpha \stackrel{~}{S}_{\left(3\right)}^0,`$ (30)
where we have defined
$$F_{\left(4\right)}^\alpha \stackrel{~}{D}G_{\left(3\right)}^\alpha .$$
(31)
The second-order equations (4) (which are nothing but Bianchi identities following from (6)) become
$`d(\mathrm{\Phi }^1\stackrel{~}{S}_{\left(3\right)}^0)`$ $`=`$ $`M_{\alpha \beta }G_{\left(3\right)}^\alpha G_{\left(1\right)}^\beta +G_{\left(2\right)}^\alpha G_{\left(3\right)}^\alpha ,`$
$`\stackrel{~}{D}(M_{\alpha \beta }G_{\left(3\right)}^\beta )`$ $`=`$ $`\stackrel{~}{F}_{\left(2\right)}^{\alpha \beta }G_{\left(3\right)}^\beta G_{\left(2\right)}^\alpha \stackrel{~}{S}_{\left(3\right)}^0.`$ (32)
The Yang-Mills equations (9) become
$`\stackrel{~}{D}(\mathrm{\Phi }M_{\alpha \beta }^1G_{\left(2\right)}^\beta )`$ $`=`$ $`\stackrel{~}{g}\mathrm{\Phi }M_{\alpha \beta }G_{\left(1\right)}^\beta \stackrel{~}{S}_{\left(3\right)}^0G_{\left(3\right)}^\alpha \frac{1}{2}ϵ_{\alpha \beta _1\beta _2\beta _3}M_{\beta _3\gamma }\stackrel{~}{F}_{\left(2\right)}^{\beta _1\beta _2}G_{\left(3\right)}^\gamma ,`$
$`\stackrel{~}{D}\left[M_{\gamma \alpha }^1M_{\delta \beta }^1\stackrel{~}{F}_{\left(2\right)}^{\gamma \delta }\right]`$ $`=`$ $`2\stackrel{~}{g}M_{\gamma [\alpha }^1\stackrel{~}{D}M_{\beta ]\gamma }G_{\left(3\right)}^\alpha G_{\left(3\right)}^\beta +\mathrm{\Phi }M_{\alpha \gamma }^1G_{\left(1\right)}^\beta G_{\left(2\right)}^\gamma `$
$`\mathrm{\Phi }M_{\beta \gamma }^1G_{\left(1\right)}^\alpha G_{\left(2\right)}^\gamma ϵ_{\alpha \beta \gamma \delta }M_{\delta \lambda }G_{\left(2\right)}^\gamma G_{\left(3\right)}^\lambda \frac{1}{2}\mathrm{\Phi }^1ϵ_{\alpha \beta \gamma \delta }\stackrel{~}{F}_{\left(2\right)}^{\gamma \delta }\stackrel{~}{S}_{\left(3\right)}^0.`$
Finally, the scalar field equations (10) give the following:
$`d(\mathrm{\Phi }^1d\mathrm{\Phi })`$ $`=`$ $`\mathrm{\Phi }M_{\alpha \beta }G_1^\alpha G_{\left(1\right)}^\beta +\mathrm{\Phi }M_{\alpha \beta }^1G_{\left(2\right)}^\alpha G_{\left(2\right)}^\beta \mathrm{\Phi }^1\stackrel{~}{S}_{\left(3\right)}^0\stackrel{~}{S}_{\left(3\right)}^0+\frac{1}{5}Q,`$
$`\stackrel{~}{D}(\mathrm{\Phi }M_{\alpha \beta }G_{\left(1\right)}^\beta )`$ $`=`$ $`\mathrm{\Phi }M_{\beta \gamma }^1G_{\left(2\right)}^\gamma \stackrel{~}{F}_{\left(2\right)}^{\alpha \beta }M_{\alpha \beta }G_{\left(3\right)}^\beta \stackrel{~}{S}_{\left(3\right)}^0,`$ (34)
$`\stackrel{~}{D}(M_{\alpha \gamma }^1\stackrel{~}{D}M_{\gamma \beta })`$ $`=`$ $`\mathrm{\Phi }M_{\beta \gamma }G_{\left(1\right)}^\gamma G_{\left(1\right)}^\alpha +M_{\beta \gamma }G_{\left(3\right)}^\gamma G_{\left(3\right)}^\alpha \mathrm{\Phi }M_{\alpha \gamma }^1G_{\left(2\right)}^\gamma G_{\left(2\right)}^\beta `$
$`+M_{\alpha \gamma }^1M_{\lambda \delta }^1\stackrel{~}{F}_{\left(2\right)}^{\gamma \delta }\stackrel{~}{F}_{\left(2\right)}^{\lambda \beta }+2\stackrel{~}{g}^2(2M_{\alpha \gamma }M_{\gamma \beta }M_{\gamma \gamma }M_{\alpha \beta })ϵ_{\left(7\right)}\frac{1}{5}\delta _{\alpha \beta }Q.`$
In these equations, the quantity $`Q`$ is the limit of the trace term multiplying $`\delta _{ij}`$ in (10), and is given by
$`Q`$ $`=`$ $`2\stackrel{~}{g}^2\left(2M_{\alpha \beta }M_{\alpha \beta }(M_{\alpha \alpha })^2\right)ϵ_{\left(7\right)}M_{\alpha \gamma }^1M_{\beta \delta }^1\stackrel{~}{F}_{\left(2\right)}^{\alpha \beta }\stackrel{~}{F}_{\left(2\right)}^{\gamma \delta }`$ (35)
$`+\mathrm{\Phi }^1\stackrel{~}{S}_{\left(3\right)}^0\stackrel{~}{S}_{\left(3\right)}^02\mathrm{\Phi }M_{\alpha \beta }^1G_{\left(2\right)}^\alpha G_{\left(2\right)}^\beta +M_{\alpha \beta }G_{\left(3\right)}^\alpha G_{\left(3\right)}^\beta .`$
Having obtained the seven-dimensional equations of motion for the $`SO(4)`$-gauged limit, we can now seek a Lagrangian from which they can be generated. A crucial point is that the equations involving $`\stackrel{~}{H}_{\left(4\right)}^0`$ in (29) and (30) give
$$d\stackrel{~}{S}_{\left(3\right)}^0=\frac{1}{8}ϵ_{\alpha _1\mathrm{}\alpha _4}\stackrel{~}{F}_{\left(2\right)}^{\alpha _1\alpha _2}\stackrel{~}{F}_{\left(2\right)}^{\alpha _3\alpha _4},$$
(36)
which allows us to strip off the exterior derivative by writing
$$\stackrel{~}{S}_{\left(3\right)}^0=dA_{\left(2\right)}+\omega _{\left(3\right)},$$
(37)
where $`\stackrel{~}{S}_{\left(3\right)}^0`$ is now viewed as a field strength with 2-form potential $`A_{\left(2\right)}`$, and
$$\omega _{\left(3\right)}\frac{1}{8}ϵ_{\alpha _1\mathrm{}\alpha _4}(\stackrel{~}{F}_{\left(2\right)}^{\alpha _1\alpha _2}\stackrel{~}{A}_{\left(1\right)}^{\alpha _3\alpha _4}\frac{1}{3}\stackrel{~}{g}\stackrel{~}{A}_{\left(1\right)}^{\alpha _1\alpha _2}\stackrel{~}{A}_{\left(1\right)}^{\alpha _3\beta }\stackrel{~}{A}_{\left(1\right)}^{\beta \alpha _4}).$$
(38)
We can now see that the equations of motion can be derived from the following seven-dimensional Lagrangian, in which $`A_{\left(2\right)}`$, and not its field strength $`\stackrel{~}{S}_{\left(3\right)}^0dA_{\left(2\right)}+\omega _{\left(3\right)}`$, is viewed as a fundamental field:
$`_7`$ $`=`$ $`R\text{1}\mathrm{l}\frac{5}{16}\mathrm{\Phi }^2d\mathrm{\Phi }d\mathrm{\Phi }\frac{1}{4}M_{\alpha \beta }^1\stackrel{~}{D}M_{\beta \gamma }M_{\gamma \delta }^1\stackrel{~}{D}M_{\delta \alpha }\frac{1}{2}\mathrm{\Phi }^1\stackrel{~}{S}_{\left(3\right)}^0\stackrel{~}{S}_{\left(3\right)}^0`$
$`\frac{1}{4}M_{\alpha \gamma }^1M_{\beta \delta }^1\stackrel{~}{F}_{\left(2\right)}^{\alpha \beta }\stackrel{~}{F}_{\left(2\right)}^{\gamma \delta }\frac{1}{2}\mathrm{\Phi }M_{\alpha \beta }^1G_{\left(2\right)}^\alpha G_{\left(2\right)}^\beta \frac{1}{2}\mathrm{\Phi }M_{\alpha \beta }G_{\left(1\right)}^\alpha G_{\left(1\right)}^\beta `$
$`\frac{1}{2}M_{\alpha \beta }G_{\left(3\right)}^\alpha G_{\left(3\right)}^\beta \stackrel{~}{V}\text{1}\mathrm{l}+{\displaystyle \frac{1}{2\stackrel{~}{g}}}\stackrel{~}{D}\stackrel{~}{S}_{\left(3\right)}^\alpha \stackrel{~}{S}_{\left(3\right)}^\alpha +\stackrel{~}{S}_{\left(3\right)}^\alpha \stackrel{~}{S}_{\left(3\right)}^0A_{\left(1\right)}^{0\alpha }`$
$`+{\displaystyle \frac{1}{2\stackrel{~}{g}}}ϵ_{\alpha \beta \gamma \delta }\stackrel{~}{S}_{\left(3\right)}^\alpha \stackrel{~}{F}_{\left(2\right)}^{0\beta }\stackrel{~}{F}_{\left(2\right)}^{\gamma \delta }+\frac{1}{4}ϵ_{\alpha _1\mathrm{}\alpha _4}\stackrel{~}{S}_{\left(3\right)}^0\stackrel{~}{F}_{\left(2\right)}^{\alpha _1\alpha _2}\stackrel{~}{A}_{\left(1\right)}^{0\alpha _3}\stackrel{~}{A}_{\left(1\right)}^{0\alpha _4}+{\displaystyle \frac{1}{\stackrel{~}{g}}}\stackrel{~}{\mathrm{\Omega }}_{\left(7\right)},`$
where $`\stackrel{~}{\mathrm{\Omega }}_{\left(7\right)}`$ is built purely from $`\stackrel{~}{A}_{\left(1\right)}^{\alpha \beta }`$ and $`\stackrel{~}{A}_{\left(1\right)}^{0\alpha }`$. It is defined by the requirement that its variations with respect to $`\stackrel{~}{A}_{\left(1\right)}^{\alpha \beta }`$ and $`\stackrel{~}{A}_{\left(1\right)}^{0\alpha }`$ should produce the necessary terms in the equations of motion for these fields. Since it has a rather complicated structure, we shall not present it here. Note that the fields that are treated as fundamental in this Lagrangian are the metric and the scalars $`(\mathrm{\Phi },M_{\alpha \beta })`$, together with $`(\chi _\alpha ,\stackrel{~}{A}_{\left(1\right)}^{\alpha \beta },\stackrel{~}{A}_{\left(1\right)}^{0\alpha },\stackrel{~}{S}_{\left(3\right)}^\alpha ,A_{\left(2\right)})`$, but it should be borne in mind that $`\mathrm{\Phi }`$ is not independent of $`M_{\alpha \beta }`$, because of the relation (25). It can be useful, therefore, to define the unimodular matrix $`\stackrel{~}{M}_{\alpha \beta }\mathrm{\Phi }^{1/4}M_{\alpha \beta }`$, so that $`\stackrel{~}{M}_{\alpha \beta }`$ and $`\mathrm{\Phi }`$ are independent fields. The scalar part of the Lagrangian (3) then becomes
$$_{\mathrm{scal}}=\frac{1}{4}\mathrm{\Phi }^2d\mathrm{\Phi }d\mathrm{\Phi }\frac{1}{4}\stackrel{~}{M}_{\alpha \beta }^1\stackrel{~}{D}\stackrel{~}{M}_{\beta \gamma }\stackrel{~}{M}_{\gamma \delta }^1\stackrel{~}{D}\stackrel{~}{M}_{\delta \alpha }.$$
(40)
## 4 $`\text{I}\mathrm{R}\times S^3`$ limit of the $`S^4`$ reduction
In the previous section, we obtained a scaling limit of the gauged $`SO(5)`$ theory in seven dimensions, in which an $`SO(4)`$ gauging survives. In this section, we apply this scaling procedure to the $`S^4`$ reduction Ansatz of section 2, showing that it leads to a degeneration in which the 4-sphere becomes $`\text{I}\mathrm{R}\times S^3`$. We can then reinterpret the reduction from $`D=11`$ as an initial “ordinary” Kaluza-Klein reduction step from $`D=11`$ to give the type IIA supergravity in $`D=10`$, followed by a non-trivial reduction of the type IIA theory on $`S^3`$, in which the entire $`SO(4)`$ isometry group is gauged.<sup>6</sup><sup>6</sup>6The $`S^3`$ reduction of type IIA supergravity discussed in , giving a seven-dimensional theory with just an $`SU(2)`$ gauging, was rederived in as a singular limit of the $`S^4`$ reduction of $`D=11`$ supergravity that was obtained in . Since the $`S^3`$ reduction in retains only the left-acting $`SU(2)`$ of the $`SO(4)SU(2)_L\times SU(2)_R`$ of gauge fields, the consistency of that reduction is guaranteed by group-theoretic arguments, based on the fact that all the retained fields are singlets under the right-acting $`SU(2)_R`$. The subtleties of the consistency of the $`S^4`$ reduction in are therefore lost in the singular limit to $`\text{I}\mathrm{R}\times S^3`$ discussed in , since a truncation to the $`SU(2)_L`$ subgroup of the $`SO(4)`$ gauge group is made. By contrast, the $`\text{I}\mathrm{R}\times S^3`$ singular limit that we consider here retains all the fields of the $`S^4`$ reduction in , and the proof of the consistency of the resulting $`S^3`$ reduction of the type IIA theory follows from the non-trivial consistency of the reduction in , and has no simple group-theoretic explanation.
### 4.1 The $`\text{I}\mathrm{R}\times S^3`$ reduction Ansatz
To take this limit, we combine the scalings of seven-dimensional quantities derived in the previous section with an appropriately-matched rescaling of the coordinates $`\mu ^i`$ defined on the internal 4-sphere. As in , we see that after splitting the $`\mu ^i`$ into $`\mu ^0`$ and $`\mu ^\alpha `$, these additional scalings should take the form
$$\mu ^0=\lambda ^5\stackrel{~}{\mu }^0,\mu ^\alpha =\stackrel{~}{\mu }^\alpha .$$
(41)
In the limit where $`\lambda `$ goes to zero, we see that the original constraint $`\mu ^i\mu ^i=1`$ becomes
$$\stackrel{~}{\mu }^\alpha \stackrel{~}{\mu }^\alpha =1,$$
(42)
implying that the $`\stackrel{~}{\mu }^\alpha `$ coordinates define a 3-sphere, while the coordinate $`\stackrel{~}{\mu }^0`$ is now unconstrained and ranges over the real line $`\text{I}\mathrm{R}`$.
Combining this with the rescalings of the previous section, we find that the $`S^4`$ metric reduction Ansatz (1) becomes
$$d\widehat{s}_{11}^2=\lambda ^{2/3}\left[\stackrel{~}{\mathrm{\Delta }}^{1/3}ds_7^2+\frac{1}{\stackrel{~}{g}^2}\stackrel{~}{\mathrm{\Delta }}^{2/3}M_{\alpha \beta }^1\stackrel{~}{D}\stackrel{~}{\mu }^\alpha \stackrel{~}{D}\stackrel{~}{\mu }^\beta +\frac{1}{\stackrel{~}{g}^2}\stackrel{~}{\mathrm{\Delta }}^{2/3}\mathrm{\Phi }(d\stackrel{~}{\mu }_0+\stackrel{~}{g}\stackrel{~}{A}_{\left(1\right)}^{0\alpha }\stackrel{~}{\mu }^\alpha +\chi _\alpha \stackrel{~}{D}\stackrel{~}{\mu }^\alpha )^2\right],$$
(43)
where
$$\stackrel{~}{\mathrm{\Delta }}M_{\alpha \beta }\stackrel{~}{\mu }^\alpha \stackrel{~}{\mu }^\beta .$$
(44)
Thus $`\stackrel{~}{\mu }_0`$ can be interpreted as the “extra” coordinate of a standard type of Kaluza-Klein reduction from $`D=11`$ to $`D=10`$, with
$$d\widehat{s}_{11}^2=e^{\frac{1}{6}\varphi }ds_{10}^2+e^{\frac{4}{3}\varphi }(d\stackrel{~}{\mu }_0+𝒜_{\left(1\right)})^2.$$
(45)
By comparing (45) with (43), we can read off the $`S^3`$ reduction Ansatz for the ten-dimensional fields. Thus we find that the ten-dimensional metric is reduced according to
$$ds_{10}^2=\mathrm{\Phi }^{1/8}\left[\stackrel{~}{\mathrm{\Delta }}^{1/4}ds_7^2+\frac{1}{\stackrel{~}{g}^2}\stackrel{~}{\mathrm{\Delta }}^{3/4}M_{\alpha \beta }^1\stackrel{~}{D}\stackrel{~}{\mu }^\alpha \stackrel{~}{D}\stackrel{~}{\mu }^\beta \right],$$
(46)
while the Ansatz for the dilaton $`\varphi `$ of the ten-dimensional theory is
$$e^{2\varphi }=\stackrel{~}{\mathrm{\Delta }}^1\mathrm{\Phi }^{3/2}.$$
(47)
Finally, the reduction Ansatz for the ten-dimensional Kaluza-Klein vector is
$$𝒜_{\left(1\right)}=\stackrel{~}{g}\stackrel{~}{A}_{\left(1\right)}^{0\alpha }\stackrel{~}{\mu }^\alpha +\chi _\alpha \stackrel{~}{D}\stackrel{~}{\mu }^\alpha .$$
(48)
These results for the $`S^3`$ reduction of the ten-dimensional metric and dilaton agree precisely with the results obtained in . (Note that the field $`\mathrm{\Phi }`$ is called $`Y`$ there, and our $`M_{\alpha \beta }`$ is called $`T_{ij}`$ there.) Note that the field strength $`_{\left(2\right)}=d𝒜_{\left(1\right)}`$ following from (48) has the simple expression
$$_{\left(2\right)}=\stackrel{~}{g}G_{\left(2\right)}^\alpha \stackrel{~}{\mu }^\alpha +G_{\left(1\right)}^\alpha \stackrel{~}{D}\stackrel{~}{\mu }^\alpha .$$
(49)
So far, we have read off the reduction Ansätze for those fields of ten-dimensional type IIA supergravity that come from the reduction of the eleven-dimensional metric. The remaining type IIA fields come from the reduction of the eleven-dimensional 4-form. Under the standard Kaluza-Klein procedure, this reduces as follows:
$$\widehat{F}_{\left(4\right)}=F_{\left(4\right)}+F_{\left(3\right)}(d\stackrel{~}{\mu }^0+𝒜_{\left(1\right)}).$$
(50)
By applying the $`\lambda `$-rescaling derived previously to the $`S^4`$ reduction Ansatz (2) for the eleven-dimensional 4-form, and comparing with (50), we obtain the following expressions for the $`S^3`$ reduction Ansätze for the ten-dimensional 4-form and 3-form fields:
$`F_{\left(4\right)}`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}^1}{\stackrel{~}{g}^3}}M_{\alpha \beta }G_{\left(1\right)}^\alpha \stackrel{~}{\mu }^\beta \stackrel{~}{W}+{\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}^1}{2\stackrel{~}{g}^2}}ϵ_{\alpha _1\mathrm{}\alpha _4}M_{\alpha _4\beta }\stackrel{~}{\mu }^\beta G_{\left(2\right)}^{\alpha _1}\stackrel{~}{D}\stackrel{~}{\mu }^{\alpha _2}\stackrel{~}{D}\stackrel{~}{\mu }^{\alpha _3}`$ (51)
$`M_{\alpha \beta }G_{\left(3\right)}^\alpha \stackrel{~}{\mu }^\beta +{\displaystyle \frac{1}{\stackrel{~}{g}}}G_{\left(3\right)}^\alpha \stackrel{~}{D}\stackrel{~}{\mu }^\alpha ,`$
$`F_{\left(3\right)}`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{U}\stackrel{~}{\mathrm{\Delta }}^2}{\stackrel{~}{g}^3}}\stackrel{~}{W}+{\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}^2}{2\stackrel{~}{g}^3}}ϵ_{\alpha _1\mathrm{}\alpha _4}M_{\alpha _1\beta }\stackrel{~}{\mu }^\beta \stackrel{~}{D}M_{\alpha _2\gamma }\stackrel{~}{\mu }^\gamma \stackrel{~}{D}\stackrel{~}{\mu }^{\alpha _3}\stackrel{~}{D}\stackrel{~}{\mu }^{\alpha _4}`$ (52)
$`+{\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}^1}{2\stackrel{~}{g}^2}}ϵ_{\alpha _1\mathrm{}\alpha _4}M_{\alpha _1\beta }\stackrel{~}{\mu }^\beta \stackrel{~}{F}_{\left(2\right)}^{\alpha _2\alpha _3}\stackrel{~}{D}\stackrel{~}{\mu }^{\alpha _4}+{\displaystyle \frac{1}{\stackrel{~}{g}}}\stackrel{~}{S}_{\left(3\right)}^0,`$
where
$$\stackrel{~}{W}\frac{1}{6}ϵ_{\alpha _1\mathrm{}\alpha _4}\stackrel{~}{\mu }^{\alpha _1}\stackrel{~}{D}\stackrel{~}{\mu }^{\alpha _2}\stackrel{~}{D}\stackrel{~}{\mu }^{\alpha _3}\stackrel{~}{D}\stackrel{~}{\mu }^{\alpha _4}.$$
(53)
It is also useful to present the expressions for the ten-dimensional Hodge duals of the field strengths:
$`\mathrm{e}^{\frac{1}{2}\varphi }\overline{}F_{\left(4\right)}`$ $`=`$ $`{\displaystyle \frac{1}{\stackrel{~}{g}}}\mathrm{\Phi }M_{\alpha \beta }G_{\left(1\right)}^\alpha \stackrel{~}{\mu }^\beta {\displaystyle \frac{1}{\stackrel{~}{g}^2}}\mathrm{\Phi }M_{\alpha \beta }^1G_{\left(2\right)}^\alpha \stackrel{~}{D}\stackrel{~}{\mu }^\beta +{\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}^1}{\stackrel{~}{g}^4}}M_{\alpha \beta }G_{\left(3\right)}^\alpha \stackrel{~}{\mu }^\beta \stackrel{~}{W}`$ (54)
$`+{\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}^1}{2\stackrel{~}{g}^3}}ϵ_{\alpha _1\mathrm{}\alpha _4}M_{\alpha _1\beta }\stackrel{~}{\mu }^\beta M_{\alpha _2\gamma }G_{\left(3\right)}^\gamma \stackrel{~}{D}\stackrel{~}{\mu }^{\alpha _3}\stackrel{~}{D}\stackrel{~}{\mu }^{\alpha _4},`$
$`\mathrm{e}^\varphi \overline{}F_{\left(3\right)}`$ $`=`$ $`\stackrel{~}{g}\stackrel{~}{U}ϵ_{\left(7\right)}{\displaystyle \frac{1}{\stackrel{~}{g}^3}}\mathrm{\Phi }^1\stackrel{~}{S}_{\left(3\right)}^0\stackrel{~}{W}`$
$`+{\displaystyle \frac{1}{2\stackrel{~}{g}^2}}M_{\alpha \gamma }^1M_{\beta \delta }^1\stackrel{~}{F}_{\left(2\right)}^{\alpha \beta }\stackrel{~}{D}\stackrel{~}{\mu }^\gamma \stackrel{~}{D}\stackrel{~}{\mu }^\delta {\displaystyle \frac{1}{\stackrel{~}{g}}}M_{\alpha \beta }^1\stackrel{~}{D}M_{\alpha \gamma }\stackrel{~}{\mu }^\gamma \stackrel{~}{D}\stackrel{~}{\mu }^\beta ,`$
$`e^{\frac{3}{2}\varphi }\overline{}_{\left(2\right)}`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}^1\mathrm{\Phi }}{\stackrel{~}{g}^5}}G_{\left(2\right)}^\alpha \stackrel{~}{\mu }^\alpha \stackrel{~}{W}+{\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}^1\mathrm{\Phi }}{2\stackrel{~}{g}^4}}ϵ_{\alpha _1\mathrm{}\alpha _4}M_{\alpha _1\beta }\stackrel{~}{\mu }^\beta M_{\alpha _2\gamma }G_{\left(1\right)}^\gamma \stackrel{~}{D}\stackrel{~}{\mu }^{\alpha _3}\stackrel{~}{D}\stackrel{~}{\mu }^{\alpha _4}.`$
(Here we are using $`\overline{}`$ to denote a Hodge dualisation in the ten-dimensional metric $`ds_{10}^2`$, to distinguish it from $``$ which denotes the seven-dimensional Hodge dual in the metric $`ds_7^2`$. )
### 4.2 Verification of the reduction Ansatz
The consistency of the $`S^3`$ reduction of the type IIA theory using the Ansatz that we obtained in the previous subsection is guaranteed by virtue of the consistency of the $`S^4`$ reduction from $`D=11`$. It is still useful, however, to examine the reduction directly, by substituting the Ansatz into the equations of motion of type IIA supergravity. By this means we can obtain an explicit verification of the validity of the limiting procedures that we applied in obtaining the $`S^3`$ reduction Ansatz.
The bosonic Lagrangian for type IIA supergravity can be written as
$`_{10}`$ $`=`$ $`R\overline{}\text{1}\mathrm{l}\frac{1}{2}\overline{}d\varphi d\varphi \frac{1}{2}\mathrm{e}^{\frac{3}{2}\varphi }\overline{}_{\left(2\right)}_{\left(2\right)}\frac{1}{2}\mathrm{e}^{\frac{1}{2}\varphi }\overline{}F_{\left(4\right)}F_{\left(4\right)}`$ (55)
$`\frac{1}{2}\mathrm{e}^\varphi \overline{}F_{\left(3\right)}F_{\left(3\right)}+\frac{1}{2}dA_{\left(3\right)}dA_{\left(3\right)}A_{\left(2\right)},`$
where
$$F_{\left(4\right)}=dA_{\left(3\right)}dA_{\left(2\right)}𝒜_{\left(1\right)},F_{\left(3\right)}=dA_{\left(2\right)},_{\left(2\right)}=d𝒜_{\left(1\right)}.$$
(56)
(In this subsection, we use a bar where necessary to indicate ten-dimensional quantities.)
The equations of motion derived from the above Lagrangian are
$`d\overline{}d\varphi `$ $`=`$ $`\frac{1}{2}e^\varphi \overline{}F_{\left(3\right)}F_{\left(3\right)}\frac{3}{4}e^{\frac{3}{2}\varphi }\overline{}_{\left(2\right)}_{\left(2\right)}\frac{1}{4}\mathrm{e}^{\frac{1}{2}\varphi }\overline{}F_{\left(4\right)}F_{\left(4\right)},`$
$`d(e^{\frac{1}{2}\varphi }\overline{}F_{\left(4\right)})`$ $`=`$ $`F_{\left(4\right)}F_{\left(3\right)},`$
$`d(e^{\frac{3}{2}\varphi }\overline{}_{\left(2\right)})`$ $`=`$ $`e^{\frac{1}{2}\varphi }\overline{}F_{\left(4\right)}F_{\left(3\right)},`$
$`d(e^\varphi \overline{}F_{\left(3\right)})`$ $`=`$ $`\frac{1}{2}F_{\left(4\right)}F_{\left(4\right)}e^{\frac{1}{2}\varphi }\overline{}F_{\left(4\right)}_{\left(2\right)}.`$ (57)
Note that it is consistent to truncate the theory to the NS-NS sector, namely the subsector comprising the metric, the dilaton and the 3-form field strength. This implies that it is possible also to perform an $`S^3`$ reduction of the NS-NS sector alone, which was indeed demonstrated in . On the other hand it is not consistent to truncate the theory to a subsector comprising only the metric, the dilaton and the 4-form field strength, which again is in agreement with the conclusion in that it is not consistent to perform an $`S^4`$ reduction on such a subsector. However, as we show in section 5, there is a consistent $`S^4`$ reduction if we include all the fields of the type IIA theory.
The reduction Ansatz obtained in section (4.1) can now be substituted into the type IIA equations of motion, to verify that it indeed leads to the equations of motion for the $`SO(4)`$-gauged seven-dimensional theory constructed in section 3.
## 5 $`S^4`$ reduction of type IIA supergravity
We can also derive the Ansatz for the consistent $`S^4`$ reduction of type IIA supergravity from the $`S^4`$ reduction Ansatz of eleven-dimensional supergravity. In this case we do not need to take any singular limit of the internal 4-sphere, but rather, we extract the “extra” coordinate from the seven-dimensional spacetime of the original eleven-dimensional supergravity reduction Ansatz. The resulting six-dimensional $`SO(5)`$ gauged maximal supergravity can be obtained from the Kaluza-Klein reduction of seven-dimensional gauged maximal supergravity on a circle.
We begin, therefore, by making a standard $`S^1`$ Kaluza-Klein reduction of the seven-dimensional metric:
$$ds_7^2=e^{2\alpha \phi }ds_6^2+e^{8\alpha \phi }(dz+\overline{𝒜}_{\left(1\right)})^2,$$
(58)
where $`\alpha =1/\sqrt{40}`$. With this parameterisation the metric reduction preserves the Einstein frame, and the dilatonic scalar $`\phi `$ has the canonical normalisation for its kinetic term in six dimensions.<sup>7</sup><sup>7</sup>7We use a bar to denote six-dimensional fields, in cases where this is necessary to avoid an ambiguity. Substituting (58) into the original metric reduction Ansatz (1), we obtain
$$d\widehat{s}_{11}^2=\mathrm{\Delta }^{1/3}e^{2\alpha \phi }ds_6^2+\frac{1}{g^2}\mathrm{\Delta }^{2/3}T_{ij}^1D\mu ^iD\mu ^j+\mathrm{\Delta }^{1/3}e^{8\alpha \phi }(dz+\overline{𝒜}_{\left(1\right)})^2.$$
(59)
In order to extract the Ansatz for the $`S^4`$ reduction of type IIA supergravity, we must first rewrite (59) in the form
$$d\widehat{s}_{11}^2=e^{\frac{1}{6}\varphi }ds_{10}^2+e^{\frac{4}{3}\varphi }(dz+𝒜_{\left(1\right)})^2,$$
(60)
which is a canonical $`S^1`$ reduction from $`D=11`$ to $`D=10`$. It is not immediately obvious that this can easily be done, since the Yang-Mills fields $`A_{\left(1\right)}^{ij}`$ appearing in the covariant differentials $`D\mu ^i`$ in (59) must themselves be reduced according to standard Kaluza-Klein rules,
$$A_{\left(1\right)}^{ij}=\overline{A}_{\left(1\right)}^{ij}+\chi ^{ij}(dz+\overline{𝒜}_{\left(1\right)}),$$
(61)
where $`\overline{A}_{\left(1\right)}^{ij}`$ are the $`SO(5)`$ gauge potentials in six dimensions, and $`\chi ^{ij}`$ are six-dimensional axions. Thus we have
$$D\mu ^i=\overline{D}\mu ^i+g\chi ^{ij}\mu ^j(dz+\overline{𝒜}_{\left(1\right)}),$$
(62)
where
$$\overline{D}\mu ^id\mu ^i+g\overline{A}_{\left(1\right)}^{ij}\mu ^j.$$
(63)
This means that the differential $`dz`$ actually appears in a much more complicated way in (59) than is apparent at first sight. Nonetheless, we find that one can in fact “miraculously” complete the square, and thereby rewrite (59) in the form of (60).
To present the result, it is useful to make the following definitions:
$`\mathrm{\Omega }`$ $``$ $`\mathrm{\Delta }^{1/3}e^{8\alpha \phi }+\mathrm{\Delta }^{2/3}T_{ij}^1\chi ^{ik}\chi ^j\mathrm{}\mu ^k\mu ^{\mathrm{}},`$
$`Z_{ij}`$ $``$ $`T_{ij}^1\mathrm{\Omega }^1\mathrm{\Delta }^{2/3}T_{ik}^1T_j\mathrm{}^1\chi ^{km}\chi ^\mathrm{}n\mu ^m\mu ^n,`$ (64)
In terms of these, we find after some algebra that we can rewrite (59) as
$$d\widehat{s}_{11}^2=\mathrm{\Delta }^{1/3}e^{2\alpha \phi }ds_6^2+\frac{1}{g^2}\mathrm{\Delta }^{2/3}Z_{ij}\overline{D}\mu ^i\overline{D}\mu ^j+\mathrm{\Omega }(dz+𝒜_{\left(1\right)})^2,$$
(65)
where the ten-dimensional potential $`𝒜_{\left(1\right)}`$ is given in terms of six-dimensional fields by
$$𝒜_{\left(1\right)}=\overline{𝒜}_{\left(1\right)}+\frac{1}{g}\mathrm{\Omega }^1\mathrm{\Delta }^{2/3}T_{ij}^1\chi ^{jk}\mu ^k\overline{D}\mu ^i.$$
(66)
This is therefore the Kaluza-Klein $`S^4`$ reduction Ansatz for the 1-form $`𝒜_{\left(1\right)}`$ of the type IIA theory. Comparing (65) with (60), we see that the Kaluza-Klein reduction Ansätze for the metric $`ds_{10}^2`$ and dilaton $`\varphi `$ of the type IIA theory are given by
$`ds_{10}^2`$ $`=`$ $`\mathrm{\Omega }^{1/8}\mathrm{\Delta }^{1/3}e^{2\alpha \phi }ds_6^2+{\displaystyle \frac{1}{g^2}}\mathrm{\Omega }^{1/8}\mathrm{\Delta }^{2/3}Z_{ij}\overline{D}\mu ^i\overline{D}\mu ^j,`$
$`e^{\frac{4}{3}\varphi }`$ $`=`$ $`\mathrm{\Omega }.`$ (67)
The $`S^4`$ reduction Ansatz for the R-R 4-form $`F_{\left(4\right)}`$ of the type IIA theory is obtained in a similar manner, by first implementing a standard $`S^1`$ Kaluza-Klein reduction on the various seven-dimensional fields appearing in the $`S^4`$ reduction Ansatz (2) for the eleven-dimensional 4-form $`\widehat{F}_{\left(4\right)}`$, and then matching this to a standard $`S^1`$ reduction of $`\widehat{F}_{\left(4\right)}`$ from $`D=11`$ to $`D=10`$:
$$\widehat{F}_{\left(4\right)}=F_{\left(4\right)}+F_{\left(3\right)}(dz+𝒜_{\left(1\right)}).$$
(68)
Note that in doing this, it is appropriate to treat the 3-form fields $`S_{\left(3\right)}^i`$ of the seven-dimensional theory as field strengths for the purpose of the $`S^1`$ reduction to $`D=6`$, viz.
$$S_{\left(3\right)}^i=\overline{S}_{\left(3\right)}^i+\overline{S}_{\left(2\right)}^i(dz+\overline{𝒜}_{\left(1\right)}).$$
(69)
It is worth noting also that this implies that the reduction of the seven-dimensional Hodge duals $`S_{\left(3\right)}^i`$ will be given by
$$S_{\left(3\right)}^i=e^{4\alpha \phi }\overline{}\overline{S}_{\left(3\right)}^i(dz+\overline{𝒜}_{\left(1\right)})+e^{6\alpha \phi }\overline{}\overline{S}_{\left(2\right)}^i,$$
(70)
where $`\overline{}`$ denotes a Hodge dualisation in the six-dimensional metric $`ds_6^2`$.
With these preliminaries, it is now a mechanical, albeit somewhat uninspiring, exercise to make the necessary substitutions into (2), and, by comparing with (68), read off the expressions for $`F_{\left(4\right)}`$ and $`F_{\left(3\right)}`$. These give the Kaluza-Klein $`S^4`$ reductions Ansätze for the 4-form and 3-form field strengths of type IIA supergravity. We shall not present the results explicitly here, since they are rather complicated, and are easily written down “by inspection” if required. For these purposes, the following identities are useful:
$`(dz+\overline{𝒜}_{\left(1\right)})`$ $`=`$ $`(dz+𝒜_{\left(1\right)}){\displaystyle \frac{1}{g}}\mathrm{\Omega }^1\mathrm{\Delta }^{2/3}T_{ij}^1\chi ^{jk}\mu ^k\overline{D}\mu ^i,`$
$`D\mu ^i`$ $`=`$ $`T_{ij}Z_{jk}\overline{D}\mu ^k+g\chi ^{ij}\mu ^j(dz+𝒜_{\left(1\right)}),`$ (71)
$`DX_i`$ $`=`$ $`\overline{D}X_i\mathrm{\Omega }^1\mathrm{\Delta }^{2/3}\chi ^{ij}X_jT_k\mathrm{}^1\chi ^\mathrm{}m\mu ^m\overline{D}\mu ^k+g\chi ^{ij}(dz+𝒜_{\left(1\right)}),`$
where in the last line $`X_i`$ represents any six-dimensional field in the vector representation of $`SO(5)`$, and the covariant derivative generalises to higher-rank $`SO(5)`$ tensors in the obvious way.
If we substitute the $`S^4`$ reduction Ansätze given for the ten-dimensional dilaton, metric and 1-form in (67), and (66), together with those for $`F_{\left(4\right)}`$ and $`F_{\left(3\right)}`$ as described above, into the equations of motion of type IIA supergravity, we shall obtain a consistent reduction to six dimensions. This six-dimensional theory will be precisely the one that follows by performing an ordinary $`S^1`$ Kaluza-Klein reduction on the $`SO(5)`$-gauged maximal supergravity in $`D=7`$, whose bosonic Lagrangian is given in (3).
It is perhaps worth remarking that the expression (67) for the Kaluza-Klein $`S^4`$ reduction of the type IIA supergravity metric illustrates a point that has been observed previously (for example in ), namely that the Ansatz becomes much more complicated when axions or pseudoscalars are involved. Although the axions $`\chi ^{ij}`$ would not be seen in the metric Ansatz in a linearised analysis, they make an appearance in a rather complicated way in the full non-linear Ansatz that we have obtained here, for example in the quantities $`\mathrm{\Omega }`$ and $`Z_{ij}`$ defined in (64). They will also, of course, appear in the Ansätze for the $`F_{\left(4\right)}`$ and $`F_{\left(3\right)}`$ field strengths. It may be that the results we are finding here could be useful in other contexts, for providing clues as to how the axionic scalars should appear in the Kaluza-Klein reduction Ansatz.
## 6 Conclusions
In this paper, we have obtained a consistent 3-sphere reduction of type IIA supergravity, in which all the massless $`SO(4)`$ gauge bosons associated with the isometry group of the 3-sphere are retained. The resulting seven-dimensional gauged supergravity will, accordingly, be maximally supersymmetric. It is, however, not a theory that admits an AdS<sub>7</sub> vacuum solution, but rather, it allows a domain wall as its “most symmetric” ground state. Since the 3-sphere is isomorphic to $`SU(2)`$ our construction can be set in the context of a string propagating in a group-manifold background. However, the reduction of fields that we considered here goes beyond what is customarily included in such cases, since we can retain the entire set of $`SO(4)SU(2)_L\times SU(2)_R`$ Yang-Mills fields, and not merely those of either the left-acting or right-acting $`SU(2)`$.
It is perhaps worth emphasising that although we can interpret the $`\text{I}\mathrm{R}\times S^3`$ limit of the $`S^4`$ reduction from $`D=11`$ as an $`S^3`$ reduction of the type IIA theory, we cannot reverse the roles of the $`\text{I}\mathrm{R}`$ and $`S^3`$ factors and interpret the limit as an $`S^3`$ reduction of eleven-dimensional supergravity to give an $`SO(4)`$-gauged supergravity in $`D=8`$, which then undergoes a further reduction to $`D=7`$. The reason for this is that when the limiting procedure is applied to the $`\mu ^i`$ coordinates of $`S^4`$, as in (41), the original coordinate $`\mu ^0`$ is set to zero, and so all fields necessarily become independent of the rescaled coordinate $`\stackrel{~}{\mu }^0`$ on the $`\text{I}\mathrm{R}`$ factor. This means that the consistent reduction involving $`S^3`$ in the limit works only if the fields are all assumed to be independent of the coordinate $`\stackrel{~}{\mu }^0`$ as well, and so there would be no possibility of extracting an eight-dimensional covariant theory by just considering the $`S^3`$ factor in the $`\text{I}\mathrm{R}\times S^3`$ reduction.
The consistent $`S^3`$ reduction of type IIA supergravity that we have constructed in this paper represents another element in the accumulating body of examples where “remarkable” Kaluza-Klein sphere reductions exist, even though there is no known group-theoretic explanation for their consistency. What is still lacking is a deeper understanding of why they should work. One might be tempted to think that supersymmetry could provide the key, but this evidently cannot in general be the answer, since there are examples such as the consistent $`S^3`$ and $`S^{D3}`$ reductions of the $`D`$-dimensional low-energy limit of the bosonic string (in arbitrary dimension $`D`$) which are obviously unrelated to supersymmetry.
As we discussed in introduction, we expect further examples of consistent sphere reduction in type IIA and type IIB supergravities. In particular, for non-trivial vacuum NS-NS flux, we expect that it is consistent to reduce both type IIA and type IIB on $`S^3`$ and $`S^7`$. For non-trivial vacuum R-R flux, we expect that it is consistent to reduce the type IIA theory on $`S^n`$ with $`n=2,4,6,8`$ and for the type IIB theory on $`S^n`$ with $`n=1,3,5,7`$. The resulting maximal gauged supergravities in the lower dimensions in general have domain-walls rather than AdS as vacuum solutions, except in the case $`n=5`$ for type IIB. We constructed two such examples in this paper, namely the $`S^3`$ and $`S^4`$ reductions of the type IIA theory. These domain-wall supergravities provide useful tools with which to explore the Domain Wall/QFT correspondence.
## Acknowledgements
C.N.P. is grateful for hospitality at the University of Pennsylvania during the early stages of this work, and at the Caltech-USC Center for Theoretical Physics during its completion.
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# I. Introduction
## I. Introduction
The pioneering work of Connes, Douglas, and Schwarz (CDS) revealing the equivalence between noncommutative Yang-Mills theory living on the noncommutative torus and toroidally compactified IKKT(and also BFSS) M(atrix) theory with the constant 3-form background field has spurred various works on noncommutative geometry and M/string theory since then. It has soon been known that the T-duality of M(atrix) theory can be understood in terms of Morita equivalence of the vector bundles over noncommutative tori .
Many of these works have been related to the torus compactification and not much has been addressed to the noncommutative orbifold case. Recently, Konechny and Schwarz worked out the compactification of M(atrix) theory on the $`_2`$ orbifold of the noncommutative two torus. However, physically more relevant compactification on the $`_2`$ orbifold of noncommutative 4-torus, a singular $`K_3`$ surface, has not been worked out so far. In the commutative case, systems of D0-branes on the commutative orbifold $`𝕋^4/_2`$ were studied in , and it is our main objective to extend the result of to the noncommutative case.
We consider the compactification in the context of IKKT M(atrix) model on the orbifold $`𝕋^4/_2`$ where $`_2`$ acts as a central symmetry $`xx`$. Thus, we need to find a Hilbert space $``$ and unitary representations of $`^4`$ and $`_2`$ on $``$ and Hermitian operators $`X`$ such that
$`U_iX_jU_i^1=X_j+2\pi \delta _i^jR_i`$ (1)
$`U_iX_\nu U_i^1=X_\nu `$ (2)
$`\mathrm{\Omega }X_i\mathrm{\Omega }=X_i`$ (3)
$`\mathrm{\Omega }X_\nu \mathrm{\Omega }=X_\nu ,\text{ }\nu =0,5,\mathrm{},9,`$ (4)
Following the description of and we can find operator relations compatible with the quotient conditions (1)-(4):
$`U_iU_j=e^{2\pi i\theta _{ij}}U_jU_i,`$ (5)
$`\mathrm{\Omega }U_i\mathrm{\Omega }=U_i^1,\text{ }\mathrm{\Omega }^2=1.`$ (6)
When $`\theta =0`$, the relations (5), (6) describe a $`_2`$ equivariant vector bundle on the $`_2`$ space $`𝕋^4`$ and $`X_i`$ specify an equivariant connection on the bundle. Now the equivariant version of the Serre-Swan theorem indicates that there is a one-to-one correspondence between $`_2`$ equivariant vector bundles on the $`_2`$ space $`𝕋^4`$ and finitely generated projective modules over the crossed product $`C^{}`$-algebra $`C(𝕋^4)_\alpha _2`$. As a noncommutative analogue we see that the relations (5), (6) imply that the Hilbert space $``$ is simply a module over the crossed product algebra $`C(𝕋_\theta ^4)_\alpha _2`$ or $`𝒜_\theta _\alpha _2`$, where $`\alpha `$ denotes the action of $`_2`$ on $`𝒜_\theta `$ by involution. The crossed product $`𝒜_\theta _\alpha _2`$ is the $`C^{}`$-completion of the linear space of $`𝒜_\theta `$-valued functions on $`_2`$. Thus a general element of $`𝒜_\theta _\alpha _2`$ is a formal linear combinations of elements of the form $`_iU_i^{n_i}\mathrm{\Omega }^{ϵ_i}`$, where $`ϵ_i\{0,1\}`$. As noted in , a $`𝒜_\theta `$-module is a finitely generated projective module if and only if its corresponding module over $`𝒜_\theta _\alpha _2`$ is finitely generated projective. Thus, bundles on a NC torus $`𝕋_\theta ^4`$ is closely related with bundles on the noncommutative torodial orbifold $`𝕋_\theta ^4/_2`$.
In this paper, we find a projective module solution to the quotient conditions (1)-(4). First we calculate a CDS type solution of M(atrix) theory compactified on the noncommutative 4-torus. There, we also show explicitly that the dual tori are actually related to each other through SO(4,4$`|`$) transformations. From this solution we discuss that the moduli space of constant curvature connections can be identified with ordinary 4-torus. Based on such explicit CDS type solution on noncommutative $`𝕋^4`$, we find its $`_2`$ orbifold solutions extending the result of to the noncommutative torodial orbifold $`𝕋_\theta ^4/_2`$.
In Section II, we review the projective modules over noncommutative torus. In Section III, we construct a projective module on noncommutative 4-torus a la Rieffel explicitly, and find a CDS type solution of M(atrix) theory compactified on the noncommutative 4-torus. It is also shown that the dual torus is actually related via SO(4,4$`|`$) transformation. In Section IV, we find a solution for the noncommutative toroidal orbifold. From this solution we study the moduli space of equivariant constant curvature connections. We conclude in Section V.
## II. Noncommutative vector bundles over noncommutative torus
In this section we review noncommutative vector bundles over NC $`d`$-torus $`𝕋_\theta ^d`$, following the lines of . Recall that $`𝕋_\theta ^d`$ is the deformed algebra of the algebra of smooth functions on the torus $`𝕋^d`$ with the deformation parameter $`\theta `$, which is a real $`d\times d`$ anti-symmetric matrix. This algebra is generated by operators $`U_1,\mathrm{},U_d`$ obeying the following relations
$`U_iU_j=e^{2\pi i\theta _{ij}}U_jU_i\text{ and }U_i^{}U_i=U_iU_i^{}=1,\text{ }i,j=1,\mathrm{},d.`$
The above relations define the presentation of the involutive algebra
$$𝒜_\theta ^d=\{a_{i_1\mathrm{}i_d}U_1^{i_1}\mathrm{}U_d^{i_d}a=(a_{i_1\mathrm{}i_d})𝒮(^d)\}$$
where $`𝒮(^d)`$ is the Schwartz space of sequences with rapid decay. According to the dictionary in , the construction of a noncommutative vector bundle over $`𝕋_\theta ^d`$ corresponds to the construction of finitely generated projective modules over $`𝒜_\theta ^d`$. It was proved in that every projective module over a smooth algebra $`𝒜_\theta ^d`$ can be represented by a direct sum of modules of the form $`𝒮(^p\times ^q\times F)`$, the linear space of Schwartz functions on $`^p\times ^q\times F`$, where $`2p+q=d`$ and $`F`$ is a finite abelian group. The module action is specified by operators on $`𝒮(^p\times ^q\times F)`$ and the commutation relation of these operators should be matched with that of elements in $`𝒜_\theta ^d`$.
On such bundles or modules there are notions of connections and the Chern character . Recall that there is the dual action of the torus group $`𝕋^d`$ on $`𝒜_\theta ^d`$ which gives a Lie group homomorphism of $`𝕋^d`$ into the group of automorphisms of $`𝒜_\theta ^d`$. Its infinitesimal form generates a homomorphism of Lie algebra $`L`$ of $`𝕋^d`$ into Lie algebra of derivations of $`𝒜_\theta ^d`$. Note that the Lie algebra $`L`$ is abelian and is isomorphic to $`^d`$. Let $`\delta :L\mathrm{Der}(𝒜_\theta ^d)`$ be the homomorphism. For each $`XL`$, $`\delta (X):=\delta _X`$ is a derivation i.e., for $`u,v𝒜_\theta ^d`$,
$`\delta _X(uv)=\delta _X(u)v+u\delta _X(v).`$
Derivations corresponding to the generators $`\{e_1,\mathrm{},e_d\}`$ of $`L`$ will be denoted by $`\delta _1,\mathrm{},\delta _d`$. For the generators $`U_i`$’s of $`𝕋_\theta ^d`$, it has the following property
$`\delta _i(U_j)=2\pi i\delta _{ij}U_j.`$
If $`E`$ is a projective $`𝒜_\theta ^d`$-module, a connection $``$ on $`E`$ is a linear map from $`E`$ to $`EL^{}`$ such that for all $`XL`$,
$`_X(\xi u)=(_X\xi )u+\xi \delta _X(u),\xi E,u𝒜_\theta ^d.`$
It is easy to see that
$`[_i,U_j]=2\pi i\delta _{ij}U_j.`$
Furthermore, for an $`𝒜_\theta ^d`$-valued inner product $`,`$ on $`E`$, if $``$ has the property that
$`_X\xi ,\eta +\xi ,_X\eta =\delta _X(\xi ,\eta ),`$
then it is called a Hermitian connection. The curvature $`_{}`$ of a connection $``$ is a 2-form on $`L`$ with values in the algebra of endomorphisms of $`E`$. That is, for $`X,YL`$,
$$_{}(X,Y):=[_X,_Y]_{[X,Y]}.$$
Since $`L`$ is abelian, we simply have $`_{}(X,Y)=[_X,_Y]`$. Denote by $`=\mathrm{End}_{𝒜_\theta }(E)`$ the algebra of endomorphisms of $`E`$. Note that if $``$ and $`^{}`$ are two Hermitian connections, then $`_X_X^{}`$ belongs to the algebra $``$. Thus once we have fixed a connection $``$, then every other connections is of the form $`+A`$, here $`A`$ is a linear map $`L`$ into $``$. In other words, the space of Hermitian connections is an affine space with vector space consisting of the linear maps from $`L`$ to $``$ and also the algebra is related with a moduli space of a certain connections.
We now consider the endomorphisms algebra of a module over $`𝒜_\theta ^d`$. Let $`\mathrm{\Lambda }`$ be a lattice in $`H=M\times \widehat{M}`$, where $`M=^p\times ^q\times F`$ and $`\widehat{M}`$ is its dual. Let $`T`$ be the corresponding embedding map in the sense of . Thus $`\mathrm{\Lambda }`$ is the image of $`^d`$ under the map $`T`$ and this determines a projective module which will be denoted by $`E_\mathrm{\Lambda }`$. Consider the lattice
$`\mathrm{\Lambda }^{}:=\{(m,\widehat{s})M\times \widehat{M}\theta ((m,\widehat{s}),(n,\widehat{t}))=\widehat{t}(m)\widehat{s}(n),\text{ for all }(n,\widehat{t})\mathrm{\Lambda }\}.`$
From the definition, it is easy to see that every operator of the form
$$𝒰_{(m,\widehat{s})}=(n)=e^{2\pi i\widehat{s}(n)}f(n+m)$$
for $`(m,\widehat{s})\mathrm{\Lambda }^{}`$, commutes with all operators $`𝒰_{(n,\widehat{t})}`$, $`(n,\widehat{t})\mathrm{\Lambda }`$. In fact one can show that the algebra of endomorphisms on $`E_\mathrm{\Lambda }`$, denoted by $`\text{End}_{𝒜_\theta }(E_\mathrm{\Lambda })`$, is a $`C^{}`$-algebra which is obtained by $`C^{}`$-completion of the space spanned by operators $`𝒰_{(m,\widehat{s})}`$, $`(m,\widehat{s})\mathrm{\Lambda }^{}`$. As shown in , the algebra $`\text{End}_{𝒜_\theta }(E_\mathrm{\Lambda })`$ can be identified with a noncommutative torus $`𝒜_{\widehat{\theta }}`$, here $`\widehat{\theta }`$ is a bilinear form on $`\mathrm{\Lambda }^{}`$,i.e., $`𝒜_{\widehat{\theta }}`$ is Morita equivalent to $`𝒜_\theta `$. Recall that a $`C^{}`$-algebra $`A`$ is said to be (strongly) Morita equivalent to $`A^{}`$ if $`A^{}\text{End}_A(E)`$ for some finite projective module $`E`$. In general, as was proved in , a NC torus $`𝒜_{\stackrel{~}{\theta }}`$ is Morita equivalent to $`𝒜_\theta `$ if $`\theta `$ and $`\stackrel{~}{\theta }`$ are related by $`\stackrel{~}{\theta }=(A\theta +B)(C\theta +D)^1`$, where $`\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)\text{ SO}(d,d|).`$
We shall now turn to the description of the Chern character. In general $`K_0(𝒜_\theta ^d)`$ classifies projective modules over $`𝒜_\theta ^d`$. In fact the positive cone $`K_0^+(𝒜_\theta ^d)`$ corresponds to genuine projective modules and if $`\theta `$ is not rational, $`K_0^+(𝒜_\theta ^d)`$ consists exactly of its elements of strictly positive trace. The Chern character of a gauge bundle on a noncommutative torus is an element in the Grassmann algebra $`^{}(L^{})`$, where $`L`$ denotes the Lie algebra of $`𝕋^d`$ and $`L^{}`$ is the dual vector space of $`L`$. Since there is a lattice $`D`$ in $`L`$, we see that there are elements of $`^{}D^{}`$ which are integral. Now the Chern character is the map $`\mathrm{Ch}:K_0(𝒜_\theta ^d)^{\mathrm{ev}}(L^{})`$ defined by
$`\mathrm{Ch}(E):=\widehat{\tau }(e^{\frac{}{2\pi i}})={\displaystyle \underset{k=0}{}}{\displaystyle \frac{1}{(2\pi i)^k}}{\displaystyle \frac{\widehat{\tau }(^k)}{k!}},`$
where $`E`$ is any gauge bundle and $``$ is a curvature of an arbitrary connection on $`E`$ and $`\widehat{\tau }`$ is a trace on the algebra of endomorphisms. In general the Chern character is integral in the commutative case. This is no longer true for the noncommutative case. However, in the case of noncommutative torus, there is an integral element related to the Chern character by the formula
$`\mathrm{Ch}(E)=e^{i(\theta )}\mu (E).`$ (7)
Here $`i(\theta )`$ denotes the contraction with the deform parameter $`\theta `$ regarded as an element of $`^2L`$. The formula (7) can be realized as a noncommutative generalization of Mukai vector. In particular, $`\mu (E)=e^{i(\theta )}\mathrm{Ch}(E)`$ is an integral element of $`^{}(L^{})`$ which is related with the Chern character on the classical torus. Also once we fix the deformation parameter, then the Chern character $`\mathrm{Ch}(E)`$ is completely determined by its integral part $`\mu (E)`$. Note that if the 0th component of the Chern character or the trace is strictly positive, then the gauge bundle $`E`$ belongs to the positive cone of $`K_0(𝒜_\theta ^d)`$ and hence it can be written as a direct sum of the form $`𝒮(^p\times ^q\times F)`$, .
## III. Compactification on noncommutative $`𝕋^4`$.
In this section we study the compactification solutions on a noncommutative 4-torus $`𝕋_\theta ^4`$ for the case $`e^{2\pi i\theta _{ij}}1`$, following the guide line in . After we fix $`U_1,U_2,U_3`$ and $`U_4`$, or a projective module, the general solution has the form of $`X_i=\overline{X}_i+A_i`$, where $`\overline{X}_i`$ are particular solutions and $`A_i`$ are operators commuting with $`U_i`$. Here we consider a projective module of the form $`𝒮(^p\times ^q)𝒮(F)`$, where $`2p+q=4`$. Thus there are three types of modules over $`𝒜_\theta `$ according to $`p=0,1,2`$. When $`p=0`$, it is a free module. The other two types are of the form $`𝒮(\times ^2)𝒮(F)`$ and $`𝒮(^2)𝒮(F)`$. As is discussed in Section II, a gauge bundle on $`𝕋_\theta ^4`$ correspond to an element of positive trace which is the 0th component of the Chern character and the Chern character is determined by its integral part $`\mu `$. Thus it is natural to start with the construction on $`𝒮(F)`$ to describe projective modules. Here we will only consider the case when $`p=2`$ which is related with (4220)-systems with a constant curvature considered in . Let $`F=_{M_1}\times _{M_2}`$, where $`_{M_i}=/M_i`$, ($`i=1,2`$) and consider the space $`^{M_1}^{M_2}`$ as the space of functions on $`C(_{M_1}\times _{M_2})`$. For all $`M_i`$ and $`N_i/M_i`$ such that $`M_i`$ and $`N_i`$ are relatively prime, define operators $`W_i`$ on $`C(_{M_1}\times _{M_2})`$ by
$`(W_1f)(k_1,k_2)`$ $`=f(k_1N_1,k_2)`$
$`(W_2f)(k_1,k_2)`$ $`=\mathrm{exp}({\displaystyle \frac{2\pi ik_1}{M_1}})f(k_1,k_2)`$
$`(W_3f)(k_1,k_2)`$ $`=f(k_1,k_2N_2)`$
$`(W_4f)(k_1,k_2)`$ $`=\mathrm{exp}({\displaystyle \frac{2\pi ik_2}{M_2}})f(k_1,k_2).`$
The operators satisfy the commutation relation
$`W_1W_2`$ $`=\mathrm{exp}(2\pi i{\displaystyle \frac{N_1}{M_1}})W_2W_1`$
$`W_3W_4`$ $`=\mathrm{exp}(2\pi i{\displaystyle \frac{N_2}{M_2}})W_4W_3,`$
otherwise commuting. If we write $`W_iW_j=\mathrm{exp}(2\pi i\psi _{ij})W_jW_i`$, then the antisymmetric $`4\times 4`$ matrix $`\psi =(\psi _{ij})`$ is of the form
$`\psi =\left(\begin{array}{cccc}0& \frac{N_1}{M_1}& 0& 0\\ \frac{N_1}{M_1}& 0& 0& 0\\ 0& 0& 0& \frac{N_2}{M_2}\\ 0& 0& \frac{N_2}{M_2}& 0\end{array}\right).`$ (8)
Let $`T:^4^2\times ^2`$ be an embedding map. Thus its matrix representation $`T=\left(\begin{array}{c}x_{ij}\end{array}\right)`$, $`i,j=1,\mathrm{},4`$, has nonzero determinant and satisfies $`(^2T^{})(\omega )=\gamma `$ where $`\omega =e_3e_1+e_4e_2^2(^4)`$ and $`e_i`$ are standard basis for $`^4`$. Equivalently, if we consider the Heisenberg representation of $`^4`$ in a Hilbert space, the desired operators acting on the space of smooth functions on $`^2`$ are defined by the following form:
$`(V_if)(s_1,s_2)=(V_{e_i}f)(s_1,s_2):=\mathrm{exp}(2\pi i(s_1x_{3i}+s_2x_{4i}))f(s_1+x_{1i},s_2+x_{2i}).`$
These operators obey the commutation relation
$`V_iV_j=e^{2\pi i\gamma _{ij}}V_jV_i,`$
where
$`\gamma _{ij}=\left|\begin{array}{cc}x_{1i}& x_{1j}\\ x_{3i}& x_{3j}\end{array}\right|+\left|\begin{array}{cc}x_{2i}& x_{2j}\\ x_{4i}& x_{4j}\end{array}\right|.`$
Since $`\gamma `$ is a real matrix, the operators $`V_i`$ act on the Schwartz space $`𝒮(^2)`$. Now we define operators $`U_i=V_iW_i`$ acting on the space $`E_T:=𝒮(^2)^{M_1}^{M_2}`$ as follows
$`(U_1f)(s_1,s_2,k_1,k_2)`$ $`=e^{2\pi i(s_1x_{31}+s_2x_{41})}f(s_1+x_{11},s_2+x_{21},k_1N_1,k_2)`$
$`(U_2f)(s_1,s_2,k_1,k_2)`$ $`=e^{2\pi i(s_1x_{32}+s_2x_{42})}e^{\frac{2\pi ik_1}{M_1}}f(s_1+x_{12},s_2+x_{22},k_1,k_2)`$
$`(U_3f)(s_1,s_2,k_1,k_2)`$ $`=e^{2\pi i(s_1x_{33}+s_2x_{43})}f(s_1+x_{13},s_2+x_{23},k_1,k_2N_2)`$
$`(U_4f)(s_1,s_2,k_1,k_2)`$ $`=e^{2\pi i(s_1x_{34}+s_2x_{44})}e^{\frac{2\pi ik_2}{M_2}}f(s_1+x_{14},s_2+x_{24},k_1,k_2).`$
Then it is easy to see that they satisfy
$`U_iU_j=\mathrm{exp}(2\pi i\gamma _{ij}+2\pi i\psi _{ij})U_jU_i.`$
Thus we have solution of (5) if $`\gamma =\psi \theta `$.
Consider operators $`\overline{X}_i`$ acting on $`E_T=𝒮(^2)^{M_1}^{M_2}`$ given by
$`(\overline{X}_if)(s_1,s_2,k_1,k_2)`$ $`=2\pi iA_i^1s_1f(s_1,s_2,k_1,k_2)+2\pi iA_i^2s_2f(s_1,s_2,k_1,k_2)`$
$`A_i^3{\displaystyle \frac{f(s_1,s_2,k_1,k_2)}{s_1}}A_i^4{\displaystyle \frac{f(s_1,s_2,k_1,k_2)}{s_2}},`$ (9)
where $`A_i^k`$ are any real numbers yet to be determined. From the definition of $`U_i`$ and $`\overline{X}_i`$, it is easy to see that the operators $`W_i`$ are commute with $`\overline{X}_i`$. Suppose that the operators $`\overline{X}_i`$ satisfy the equation (1), i.e.,
$$U_i\overline{X}_jU_i^1=\overline{X}_j+2\pi \delta _i^jR_i.$$
By a straightforward calculation, the constant matrix $`(A_i^j)`$ in (9) can be obtained as in the following form:
$$\left(\begin{array}{c}R_iA_i^j\end{array}\right)T=i\text{ Id}.$$
Since the inverse matrix of $`T`$ can be written as
$$T^1=\frac{1}{detT}\left(\begin{array}{c}(1)^{i+j}B_{ji}\end{array}\right),$$
where $`B_{ij}`$ is the $`(ij)`$-minor of the matrix $`T`$, we see that
$`A_i^k=(1)^{i+k}{\displaystyle \frac{R_i}{i}}{\displaystyle \frac{1}{detT}}B_{ki},`$ (10)
and this gives a particular solution to the equations (2) and (3). It is easy to check that the commutator has of the form
$`[\overline{X}_i,\overline{X}_j]=2\pi i\left(\left|\begin{array}{cc}A_i^1& A_i^3\\ A_j^1& A_j^3\end{array}\right|+\left|\begin{array}{cc}A_i^2& A_i^4\\ A_j^2& A_j^4\end{array}\right|\right).`$
By (10), we have
$`[\overline{X}_i,\overline{X}_j]`$ $`=2\pi i{\displaystyle \frac{R_iR_j}{(detT)^2}}\left\{(1)^{i+1}(1)^{j+1}\left|\begin{array}{cc}B_{1i}& B_{3i}\\ B_{1j}& B_{3j}\end{array}\right|+(1)^i(1)^j\left|\begin{array}{cc}B_{2i}& B_{4i}\\ B_{2j}& B_{4j}\end{array}\right|\right\}`$
$`=2\pi i(1)^{i+j+1}{\displaystyle \frac{R_iR_j}{(detT)^2}}\left\{\left|\begin{array}{cc}B_{1i}& B_{3i}\\ B_{1j}& B_{3j}\end{array}\right|+\left|\begin{array}{cc}B_{2i}& B_{4i}\\ B_{2j}& B_{4j}\end{array}\right|\right\}`$
$`=2\pi i(1)^{i+j+1}{\displaystyle \frac{R_iR_j}{detT}}\gamma _{ij}.`$
Now we should find generators of the set of operators which commute with $`U_i`$’s. To find such operators we need to describe an embedding map which corresponds to the dual lattice of the lattice defined by the embedding map $`T`$ as discussed in Section II. For such a map, let
$`S=\left(\begin{array}{cccc}0& 0& 1& 0\\ 0& 0& 0& 1\\ 1& 0& 0& 0\\ 0& 1& 0& 0\end{array}\right)(T^t)^1={\displaystyle \frac{1}{detT}}\left(\begin{array}{cccc}B_{31}& B_{32}& B_{33}& B_{34}\\ B_{41}& B_{42}& B_{43}& B_{44}\\ B_{11}& B_{12}& B_{13}& B_{14}\\ B_{21}& B_{22}& B_{23}& B_{24}\end{array}\right).`$ (11)
Using the matrix (11), we define operators acting on $`E_T`$ by
$`(Z_1f)(s_1,s_2,k_1,k_2)`$ $`=e^{\frac{2\pi i(s_1B_{11}+s_2B_{21})}{M_1|T|}}e^{\frac{2\pi ib_1k_1}{M_1}}f(s_1+{\displaystyle \frac{B_{31}}{M_1|T|}},s_2{\displaystyle \frac{B_{41}}{M_1|T|}},k_1,k_2)`$
$`(Z_2f)(s_1,s_2,k_1,k_2)`$ $`=e^{\frac{2\pi i(s_1B_{12}s_2B_{22})}{M_1|T|}}f(s_1{\displaystyle \frac{B_{32}}{M_1|T|}},s_2+{\displaystyle \frac{B_{42}}{M_1|T|}},k_11,k_2)`$
$`(Z_3f)(s_1,s_2,k_1,k_2)`$ $`=e^{\frac{2\pi i(s_1B_{13}+s_2B_{23})}{M_2|T|}}e^{\frac{2\pi ib_2k_1}{M_2}}f(s_1+{\displaystyle \frac{B_{33}}{M_2|T|}},s_2{\displaystyle \frac{B_{43}}{M_2|T|}},k_1,k_2)`$
$`(Z_4f)(s_1,s_2,k_1,k_2)`$ $`=e^{\frac{2\pi i(s_1B_{14}s_2B_{24})}{M_2|T|}}f(s_1{\displaystyle \frac{B_{34}}{M_2|T|}},s_2+{\displaystyle \frac{B_{44}}{M_2|T|}},k_1,k_21),`$
where $`|T|=\text{ Pf}(\psi \theta )`$ denotes the determinant of $`T`$ and $`b_1`$, $`b_2`$ are integers such that $`a_iM_i+b_iN_i=1`$, $`a_i`$ are also integers. To check the operators $`Z_i`$ commute with all $`U_j`$’s, let $`Z_iU_j=e^{2\pi i\lambda _{ij}}U_jZ_i`$. Then it is easy to see that
$`\lambda _{ij}={\displaystyle \frac{1}{M_k|T|}}\left\{\left|\begin{array}{cc}x_{1i}& x_{3i}\\ (1)^{3+j}B_{3j}& (1)^{1+j}B_{1j}\end{array}\right|+\left|\begin{array}{cc}x_{2i}& x_{4i}\\ (1)^{4+j}B_{4j}& (1)^{2+j}B_{2j}\end{array}\right|\right\}\delta _{ij}{\displaystyle \frac{b_kN_k}{M_k}},`$ (12)
where $`k=1,2`$ depending on $`ij`$. From the relation (12),
$`\lambda _{ij}`$ $`=0\text{ when }ij`$
$`\lambda _{ii}`$ $`={\displaystyle \frac{1}{M_k}}{\displaystyle \frac{b_kN_k}{M_k}}={\displaystyle \frac{a_kM_k}{M_k}}=a_k.`$
Thus $`Z_i`$ commute with all $`U_j`$’s.
Furthermore the operators satisfy
$`Z_iZ_j=e^{2\pi i\widehat{\theta }}Z_jZ_i.`$ (13)
Now $`\widehat{\theta }`$ can be calculated directly and it is given by
$`\widehat{\theta }_{12}`$ $`={\displaystyle \frac{a_1N_2+b_1N_2\theta _{12}+a_1M_2\theta _{34}b_1M_2\text{Pf}(\theta )}{M_1M_2\text{Pf}(\psi \theta )}}`$
$`\widehat{\theta }_{13}`$ $`={\displaystyle \frac{\theta _{13}}{M_1M_2\text{Pf}(\psi \theta )}}`$
$`\widehat{\theta }_{14}`$ $`={\displaystyle \frac{\theta _{14}}{M_1M_2\text{Pf}(\psi \theta )}}`$
$`\widehat{\theta }_{23}`$ $`={\displaystyle \frac{\theta _{23}}{M_1M_2\text{Pf}(\psi \theta )}}`$
$`\widehat{\theta }_{24}`$ $`={\displaystyle \frac{\theta _{24}}{M_1M_2\text{Pf}(\psi \theta )}}`$
$`\widehat{\theta }_{34}`$ $`={\displaystyle \frac{a_2N_1+b_2N_1\theta _{34}+a_2M_1\theta _{12}b_2M_1\text{Pf}(\theta )}{M_1M_2\text{Pf}(\psi \theta )}}.`$
Also we have
$`\widehat{\theta }=(A\theta +B)(NM\theta )^1`$ (14)
where
$`A=\left(\begin{array}{cccc}0& a_1& 0& 0\\ a_1& 0& 0& 0\\ 0& 0& 0& a_2\\ 0& 0& a_2& 0\end{array}\right),\text{ }B=\left(\begin{array}{cccc}b_1& 0& 0& 0\\ 0& b_1& 0& 0\\ 0& 0& b_2& 0\\ 0& 0& 0& b_2\end{array}\right)`$
and
$`N=\left(\begin{array}{cccc}N_1& 0& 0& 0\\ 0& N_1& 0& 0\\ 0& 0& N_2& 0\\ 0& 0& 0& N_2\end{array}\right)\text{ }M=\left(\begin{array}{cccc}0& M_1& 0& 0\\ M_1& 0& 0& 0\\ 0& 0& 0& M_2\\ 0& 0& M_2& 0\end{array}\right).`$
From the equation (14), we see that $`\theta `$ and $`\widehat{\theta }`$ are related by SO$`(4,4|)`$ transformation.
Note that U$`(n)`$ theory on $`𝒜_\theta `$ is equivalent to U$`(1)`$ theory on $`𝒜_{\widehat{\theta }}`$. For U$`(1)`$ theory the generators $`Z_i`$ can be identified with functions on the dual torus:
$`Z_je^{i\sigma _j}`$
where $`\sigma _j`$ are coordinates of the dual torus such that
$`[\sigma _i,\sigma _j]=2\pi i\widehat{\theta }_{ij}.`$
Now the general solution of the compactification is given by
$`X_i=\overline{X}_i+{\displaystyle \underset{i_1,\mathrm{},i_4}{}}\mathrm{\Psi }_{i_1i_2i_3i_4}Z_1^{i_1}Z_2^{i_2}Z_3^{i_3}Z_4^{i_4},`$
where the coefficients $`\mathrm{\Psi }_{i_1i_2i_3i_4}`$ are $`c`$-numbers.
Recall that a connection in a module $`E_T`$ is determined by a set of operators $`_1,\mathrm{},_4`$ in $`E_T`$ such that
$`[_i,U_j]=2\pi i\delta _{ij}U_j.`$
From the definition of $`\overline{X}_i`$ given in (9) we have
$`[\overline{X}_i,U_j]=2\pi \delta _{ij}R_jU_j.`$
Thus we see that the special solution $`\overline{X}_i`$ is related with connections by $`\overline{X}_i=\frac{R_i}{i}_i`$ and for such connection $``$, the constant curvature $`=(_{ij})`$ is given by
$`=\gamma ^1\text{Id}_N,\text{ where }N=N_1N_2.`$ (15)
Now the general solution should be identified as
$`X_i={\displaystyle \frac{R_i}{i}}_i+A_i(\sigma _1,\sigma _2,\sigma _3,\sigma _4)`$ (16)
where $`A_i`$ are gauge fields defined on a noncommutative torus.
Note that from the curvature form (15), it corresponds to the $`U(N)`$ gauge theory with vanishing $`su(N)`$ curvature. This type of solutions has been studied in for noncommutative $`𝕋^2`$ and in for higher torus case. This was generalized to a nonvanishing $`su(N)`$ curvature case in and it has been noted that the analysis for noncommutative tori is the same as that of for commutative tori. In fact the above solution has been described by (4220) system with trivial $`SU(N)`$ gauge fields in and its moduli space can be identified with $`𝕋^4`$. So we may expect that the moduli space of constant curvature connections in noncommutative torus is of the same form as in the ordinary torus.
The operators
$`\stackrel{~}{}_j={\displaystyle \frac{i}{R_j}}\overline{X}_j+\alpha _j,\text{ }j=1,\mathrm{},4,`$ (17)
where $`\alpha _j`$ is any real number, determine a Hermitian connection with constant curvature in $`E_T`$. Furthermore connections of the form (16) define a representation on $`L^2(^2,^{M_1}^{M_2})`$ of the Heisenberg commutation relations and from this one can follow the same steps in to show that connections of the form (17) can be found in each gauge orbits and two such connections $`\frac{i}{R_j}\overline{X}_j+\alpha _j`$ and $`\frac{i}{R_j}\overline{X}_j+\mu _j`$ are gauge equivalent if and only if $`\alpha _j\mu _j`$. Thus the moduli space of constant curvature connections can be identified with $`(/)^4(S^1)^4𝕋^4`$. In general, if we consider a projective module consisting of $`n`$ copies of such modules, such as $`E_{T_1}\mathrm{}E_{T_n}`$, where $`T_i`$ is an embedding, then there is a constant curvature connection on each summand such that the overall curvature is given by $`=_k`$, where $`_k`$ is given as in (15) with the same $`\gamma `$. Thus for a constant curvature connection on $`E`$ which breaks a projective module $`E`$ into $`_kE_{T_i}`$, block diagonal construction gives the moduli space of the form $`(𝕋^4)^n/S_n`$, where $`S_n`$ is the symmetric group.
## IV. Compactification on noncommutative toroidal orbifold $`𝕋_\theta /_2`$
In this section we find solutions for the quotient conditions (1)-(4) along with the projective module actions (5) and (6) via the compactification solutions on a noncommutative torus $`𝕋_\theta ^4`$ obtained in Section III. From this we find the moduli space of equivariant constant curvature connections on noncommutative toroidal orbifold $`𝕋_\theta ^4/_2`$.
Consider the module $`E_T:=𝒮(^2)C(_{M_1})C(_{M_2})`$ together with $`U_i`$’s as operators acting on it. The general solution for the quotient conditions has been identified as
$`X_j={\displaystyle \frac{R_j}{i}}_j+A_j(\sigma _1,\sigma _2,\sigma _3,\sigma _4),\text{ }1j4.`$ (18)
To find solutions for the quotient conditions on the compactified part we need to solve for $`\mathrm{\Omega }`$ which satisfies $`\mathrm{\Omega }U_i\mathrm{\Omega }=U_i^1`$ and $`\mathrm{\Omega }^2=1`$. Consider an operator $`\mathrm{\Omega }_0`$ on $`E_T`$ defined by
$`(\mathrm{\Omega }_0f)(s_1,s_2,k_1,k_2)=f(s_1,s_2,k_1,k_2).`$
It is easy to see that $`\mathrm{\Omega }_0U_i\mathrm{\Omega }_0U_i=e^{2\pi i(x_{1i}x_{3i}+x_{2i}x_{4i})}`$. By redefining $`U_ie^{\pi i(x_{1i}x_{3i}+x_{2i}x_{4i})}U_i`$, we get $`\mathrm{\Omega }_0U_i\mathrm{\Omega }_0=U_i^1`$ and $`\mathrm{\Omega }_0^2=1`$. Thus we have a solution for (6) i.e., $`\mathrm{\Omega }_0`$ together with $`U_i`$’s define a projective module over $`𝒜_\theta _2`$. As was indicated in , there might be other $`_2`$ actions on the module. To get other actions on the module, consider the operators $`Z_i`$ defined in Section III. As for the $`U_i`$’s, rescale $`Z_i`$ by $`e^{\pi i(B_{1i}B_{3i}+B_{2i}B_{4i})}Z_i`$ and we get the relation
$`\mathrm{\Omega }_0Z_i\mathrm{\Omega }_0=Z_i^1.`$ (19)
Since $`Z_i`$ commute with all $`U_j`$’s, the operators $`\mathrm{\Omega }_{n_1\mathrm{}n_4}=e^{i\varphi }\mathrm{\Omega }_0Z_1^{n_1}Z_2^{n_2}Z_3^{n_3}Z_4^{n_4}`$, ($`n_i`$), satisfy the equation (6), where $`\varphi `$ is a phase which is chosen to get the relation $`\mathrm{\Omega }^2=1`$ and it can be calculated explicitly by using the commutation relations given in (13). Now consider the general solution (18) satisfying (1) and (2). Recall $`_i=\frac{i}{R_i}\overline{X}_i`$. For $`\overline{X}_i`$, which was defined in (9), it is easy to verify that $`\mathrm{\Omega }_0\overline{X}_i\mathrm{\Omega }_0=\overline{X}_i`$. But since $`\overline{X}_i`$ do not commute with $`Z_i`$’s, we see that $`\mathrm{\Omega }_0`$ is the unique solution for the equation $`\mathrm{\Omega }\overline{X}_i\mathrm{\Omega }=\overline{X}_i`$. By definition of the functions $`A_i`$ on the dual torus and by the relation (19), we have $`\mathrm{\Omega }_0A_i(\sigma _1,\sigma _2,\sigma _3,\sigma _4)\mathrm{\Omega }_0=A_i(\sigma _1,\sigma _2,\sigma _3,\sigma _4)`$. Applying $`\mathrm{\Omega }_0`$ to the both sides on the equation (18) we see that
$`A_i(\sigma _1,\sigma _2,\sigma _3,\sigma _4)=A_i(\sigma _1,\sigma _2,\sigma _3,\sigma _4),`$ (20)
which implies that the functions $`A_i`$ are odd functions. If we consider a constant curvature connection $``$ on $`E_T`$, the functions $`A_i`$ in (20) can be represented by a real constant and hence it vanishes. In other words the moduli space has no Higgs branch. Note that this type of solutions has been studied in for the ordinary torodial orbifold $`𝕋^4/_2`$ under the name of Rep. II.
In the above representation, the moduli space of constant curvature connections on $`E_T`$ over $`𝕋_\theta ^4`$ is not preserved by the $`_2`$ action on $`E_T`$. So it may be more natural to consider two copies of $`E_T`$ which respect the $`_2`$ action and this corresponds to Rep. I of . Consider the bundle of the form $`E_T^2=E_TE_T`$ and define operators acting on $`E_T^2`$ by
$`\mathrm{\Omega }=\left(\begin{array}{cc}\mathrm{\Omega }_0& 0\\ 0& \mathrm{\Omega }_0\end{array}\right),\text{ and }𝐔_i=\left(\begin{array}{cc}U_i& 0\\ 0& U_i\end{array}\right),`$
where $`\mathrm{\Omega }_0`$ and $`U_i`$’s are operators on $`E_T`$ given as above and in Section III. Then it is easy to check that
$`𝐔_i𝐔_j`$ $`=e^{2\pi i\theta _{ij}}𝐔_j𝐔_i,`$
$`\mathrm{\Omega }𝐔_i\mathrm{\Omega }`$ $`=𝐔_i^1\text{ and }\mathrm{\Omega }^2=1.`$ (21)
Thus the relations (IV. Compactification on noncommutative toroidal orbifold $`𝕋_\theta /_2`$) defines a projective module over $`𝒜_\theta _2`$. Since $`\overline{X}_i`$ defines a particular solution, we may write the general solution on the torus as follows
$`X_i=\overline{X}_i+\left(\begin{array}{cc}A_i^{11}& A_i^{12}\\ A_i^{21}& A_i^{22}\end{array}\right).`$
Since the matrix $`\left(\begin{array}{cc}A_i^{11}& A_i^{12}\\ A_i^{21}& A_i^{22}\end{array}\right)`$ should commute with all the $`𝐔_i`$’s, each entries $`A_i^{jk}`$ commute with $`U_i`$’s. In other words, the operators $`A_i^{jk}`$ are generated by $`Z_i`$’s. Thus they can be identified with functions on the dual torus. Now the general solutions should be identified as
$`X_i={\displaystyle \frac{R_i}{i}}_i+\left(\begin{array}{cc}A_i^{11}(\sigma _j)& A_i^{12}(\sigma _j)\\ A_i^{21}(\sigma _j)& A_i^{22}(\sigma _j)\end{array}\right).`$ (22)
By applying $`\mathrm{\Omega }`$ we find
$`\left(\begin{array}{cc}A_i^{11}(\sigma _j)& A_i^{12}(\sigma _j)\\ A_i^{21}(\sigma _j)& A_i^{22}(\sigma _j)\end{array}\right)=\left(\begin{array}{cc}A_i^{11}(\sigma _j)& A_i^{12}(\sigma _j)\\ A_i^{21}(\sigma _j)& A_i^{22}(\sigma _j)\end{array}\right).`$
Note that the diagonal entries of the matrix in (22) are odd functions on the dual torus, and this fact will be used in finding the moduli space below. Meanwhile the off-diagonal entries are even fuctions of $`\sigma `$. Here, the gauge transformation should be invariant under $`\mathrm{\Omega }`$ implementing the $`_2`$ quotient condition. This implies that the gauge parameter in general should be given by $`\mathrm{\Lambda }=\left(\begin{array}{cc}\lambda _{ev}^{11}& \lambda _{od}^{12}\\ \lambda _{od}^{21}& \lambda _{ev}^{22}\end{array}\right)`$ where the subscript $`ev`$ or $`od`$ indicates an even or odd function of $`\sigma `$. This indicates us that not all the $`U(2)`$ group acts. We now consider the constant curvature connection $``$ on $`E_T`$ considered in Section III. In this case, as discussed in Rep. II we have constant gauge field in (22). Thus the diagonal entries vanish and the bundle becomes singular at the fixed points. For the ordinary case this has been related to the existence of two-brane charge at the collapsing two-cycle of the blown-up space .
Now the solutions of the constant curvature connection in this case are given by
$`X_i`$ $`=\overline{X}_i+\left(\begin{array}{cc}0& A_i^{12}(\sigma _j)\\ A_{i}^{12}{}_{}{}^{}(\sigma _j)& 0\end{array}\right).`$
One of the $`A_i`$ components can be gauged away by constant gauge transformation of the type $`\left(\begin{array}{cc}\lambda & 0\\ 0& \widehat{\lambda }\end{array}\right)`$ which can be decomposed into two parts, one propotional to the identity and the other proportional to $`\sigma _3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$. Since we are only considering constant gauge transformations in dealing with the moduli space, the noncommutativity does not affect the result as in Section III. The remaining component of $`A_i`$ has translational symmetry of the commutative 4-torus. This fact together with a residual guage symmetry $`\sigma _3`$ now yields a Higgs branch moduli space of constant curvature connections to be an ordinary torodial orbifold $`𝕋^4/_2`$.
For the uncompactified $`X_\nu `$ sector, the solution is the same as in the commutative case ; the moduli becomes $`^5\times ^5`$ when $`A_i=0`$, and when $`A_i0`$ the transverse moduli becomes $`^5`$ for generic points in $`𝕋^4/_2`$, and $`^5\times ^5`$ at the fixed points in $`𝕋^4/_2`$. Thus this can be viewed as a fibration over the Higgs branch of $`𝕋^4/_2`$, with the fiber $`^5`$ at a generic point and with the fiber $`^5\times ^5`$ at the orbifold fixed points as suggested in the commutative case .
For the ordinary $`𝕋^4`$, the discussion above corresponds to the construction of the theory of zero branes on $`𝕋^4/_2`$. We first considered a T-duality on the covering torus $`𝕋^4`$ to a dual torus $`\widehat{𝕋}^4`$ and then project to $`\widehat{𝕋}^4/_2`$. So, for $`N`$ identical D0-branes on $`𝕋^4/_2`$ we need $`2N`$ zero branes on $`𝕋^4`$. This is described by $`U(2N)`$ gauge theory and the gauge group is broken down to $`U(N)\times U(N)`$. In , it has been shown that the moduli space of the flat connections is identified with $`𝕋^4/_2`$. In fact our above analysis on the moduli space of constant curvature connections is exactly the same as the one in .
## V. Conclusion and prospect
In this paper, we construct a bundle on noncommutative toroidal orbifold $`𝕋_\theta ^4/_2`$. We start with the construction of a bundle on noncommutative $`𝕋^4`$ a la Rieffel and find a CDS type solution of M(atrix) theory compactified on the noncommutative 4-torus. There, we also show explicitly that the dual tori are actually related to each other through SO(4,4$`|`$) transformations. Based on our explicit CDS type solution on noncommutative $`𝕋^4`$, we find its $`_2`$ orbifold solutions, Rep. I and Rep. II, by looking into the systems of D0-branes on the covering space projected onto their invariant parts under the discrete symmetry group. From the solutions obtained, we study the moduli space of equivariant constant curvature connections. The Higgs branch moduli space has been identified with the ordinary toroidal orbifold in the Rep. I case where we consider two copies of a bundle over $`𝕋_\theta `$ which are invariant under the $`_2`$ action on $`𝕋_\theta `$. In the Rep. II case, the moduli space has no Higgs branch. In conclusion, in the noncommutative $`𝕋^4/_2`$ case the moduli space has the same form as its commutative counterpart.
In , the moduli space of D0-branes on commutative $`𝕋^4`$ with torons of $`U(N)`$ Yang-Mills theory was given as $`(𝕋^4)^{p_1}/S_{p_1}\times (𝕋^4)^{p_2}/S_{p_2}`$ where $`U(N)`$ gauge group broken down into $`U(k_1)\times U(k_2)`$ satisfying $`k_1+k_2=N`$, and $`p_i=gcd(k_i,m_i)`$, $`i=1,2`$ with fluxes $`m_i`$ of $`U(k_i)`$. Its extension to the noncommutative case has been recently studied in using the ’t Hooft’s $`SU(N)`$ solution of nontrivial twists , and the resulting moduli space of connections turned out to be of the same form, $`(𝕋^4)^{p_1}/S_{p_1}\times (𝕋^4)^{p_2}/S_{p_2}`$. We expect that the same holds for the noncommutative toroidal $`_2`$ orbifold case.
Note added: After completion of our paper, a related paper has appeared, which has some overlap with our paper. Their methodology to get the relevant moduli spaces is to use the theory of representation of Heisenberg algebra defined by the commutation relations of a fixed connection. On the other hand, our approach is the usual one in that we construct a module on $`𝕋_\theta ^4`$ with explicit computation, and then consider the $`_2`$ orbifold condition on this module finding the moduli space in the specific cases.
Acknowledgments
This work was supported by Korea Research Foundation, Interdisciplinary Research Project 1998-D00001. We would like to thank KIAS and APCTP for their kind hospitality where parts of this work were done. E. K. and H. K. were also supported in part by BK 21. C.-Y. L. was also supported in part by NSF PHY-9511632 in Austin.
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# Dynamics of Cosmic Necklaces
## 1 Introduction
Topological defects are a natural consequence of phase transitions in the early universe and are as such predicted by a number of Grand Unified Theories. They have received much attention in the past two decades, having been considered promising candidates for a number of interesting cosmological phenomena including gravity waves, baryon asymmetry, density perturbations, non-gaussianity in the cosmic microwave background and ultra-high energy cosmic rays. For a review see .
Topological defects can occur when some high energy particle physics symmetry group is broken to a smaller symmetry group. They are formed by the Kibble mechanism, in which regions of space that are uncorrelated (at a minimum, those that are outside the causal horizon) must independently choose vacuum states. The type of defect formed depends on the topology of the manifold of equivalent vacuum states.
Cosmic necklaces, in particular, are hybrid topological defects that can be produced in the sequence of phase transitions
$$G\stackrel{\eta _m}{}H\times U(1)\stackrel{\eta _s}{}H\times Z_2,$$
where $`G`$ is a semi-simple group. In the first phase transition, at an energy $`\eta _m`$, monopoles of mass $`m\eta _m/e`$ are produced. The second transition, occuring at energy $`\eta _s`$, traps the magnetic flux into two strings of mass per unit length $`\mu \eta _s^2`$ connecting monopoles to anti-monopoles. The resulting system is a network of long strings and loops with the monopoles playing the role of beads.
Previous work by Berezinsky and Vilenkin was done with the idea of explaining the origin of the highest energy cosmic rays. The rather complicated underlying dynamics were, for good reason, largely ignored. They found that if one started with a low enough density of monopoles, such that one could approximate the evolution of the system using simple string evolution, and more importantly, if one could disregard the effects of monopole-antimonopole annihilation, the density of monopoles on the string would naturally increase to the point where the approximation of simple string evolution would break down. The monopole density could also be increased with a sufficiently long damping era following the formation of the network. Furthermore, they noted that a large enough density of monopoles would make loop motion non-periodic and therefore loop fragmentation very efficient. Thus, in their scenario, the network consits of a long string with monopoles in the high monopole density regime. The string intercommutes to form loops that rapidly fragment and yield cosmic rays by the annihilation of the monopoles. The estimated cosmic ray energies and fluxes produced by this model seem to fit current observations. They did leave, however, the detailed analysis of the evolution of these systems to numerical simulations and, in particular, the verification that monopole-antimonopole annihilation can initially be disregarded.
In this work we make a first attempt at understanding the dynamics of necklaces by numerically evolving single loops.
In the next section we review general string motion for completeness and derive useful results used in the rest of our paper. In the third section we look at monopole motion on cosmic necklaces when the energy of the monopoles is comparable to the energy in the string. In the fourth section we examine dynamics and monopole annihilations in the limit where the monopole energy is low compared to the string energy. We summarise and conclude in the fifth section.
## 2 Review of String Motion
When the typical length scale of a cosmic string is much larger than its thickness, $`\delta _s\eta _s^1`$, and when there are no long-range interactions between different string segments, as is the case for gauge strings, the string can be accurately modeled by a one dimensional object. Such an object sweeps out a two dimensional surface referred to as the string world-sheet. This surface can be described by a function of two parameters, a timelike parameter $`t`$ which can be identified with the time coordinate, and a spacelike parameter $`\sigma `$,
$$x^\mu =x^\mu (t,\sigma ).$$
The infinitesimal line element in Minkowski space-time with metric $`\eta _{\mu \nu }=\mathrm{diag}(1,1,1,1)`$ is
$$ds^2=\eta _{\mu \nu }dx^\mu dx^\nu =\eta _{\mu \nu }x_{,a}^\mu x_{,b}^\nu d\xi ^ad\xi ^b,$$
where $`a=0,1`$ labels the internal parameters, $`t=\xi ^0`$, $`\sigma =\xi ^1`$ and $`x_{,a}^\mu =x^\mu /\xi ^a`$. One can then write the induced metric on the world-sheet of the string as
$$\gamma _{ab}=\eta _{\mu \nu }x_{,a}^\mu x_{,b}^\nu .$$
(1)
For an infinitely thin string we can use the Nambu-Goto action. It is proportional to the invariant area swept by the string,
$$S_s=\mu 𝑑A=\mu d^2\xi \sqrt{\gamma },$$
(2)
where $`\gamma =det(\gamma _{ab})`$ and $`\mu `$ is the mass per unit length of the string. In geometrical terms this is the action for a 1+1 dimensional space-time with a cosmological constant $`\mu `$. Its variation with respect to $`x^\mu `$ gives the equation
$$\delta S_s=\mu _a(\sqrt{\gamma }\gamma ^{ab}x_{,b}^\mu \delta x_\mu )d^2\xi +\mu _a(\sqrt{\gamma }\gamma ^{ab}x_{,b}^\mu )\delta x_\mu d^2\xi =0.$$
(3)
The first term in this equation is an integral of a total divergence that can be turned into boundary terms at the ends of the string. For ordinary strings these terms vanish because of the absence of boundaries. In the case of cosmic necklaces however, monopoles live on the boundaries of strings and so these turn out to be the most interesting terms. The second term gives the usual equations of motion for the Nambu-Goto string. If we work, as usual, in the conformal gauge,
$$𝐱^{}(\sigma ,t)\dot{𝐱}(\sigma ,t)=0$$
(4)
$$𝐱^2(\sigma ,t)+\dot{𝐱}^2(\sigma ,t)=1,$$
(5)
and choose $`t=x^0`$, the equation of motion is
$$𝐱^{\prime \prime }(\sigma ,t)=\ddot{𝐱}(\sigma ,t).$$
(6)
Primes and dots denote partial derivatives with respect to $`\sigma `$ and $`t`$ respectively. This equation can be readily solved using
$$𝐱(\sigma ,t)=\frac{1}{2}[𝐚(\sigma t)+𝐛(\sigma +t)],$$
(7)
then
$`𝐱^{}(\sigma ,t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}[𝐚^{}(\sigma t)+𝐛^{}(\sigma +t)]`$ (8)
$`\dot{𝐱}(\sigma ,t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}[𝐚^{}(\sigma t)+𝐛^{}(\sigma +t)]`$ (9)
and the constraints coming from (4,5) become
$$𝐚^2(\sigma t)=𝐛^2(\sigma +t)=1.$$
(10)
The functions $`𝐚(\sigma t)`$ and $`𝐛(\sigma +t)`$, often referred to as right- and left-movers or halves of the string, are constant along the lines $`\sigma =t`$ and $`\sigma =t`$ respectively on the string world-sheet.
## 3 Monopole Motion: The Massive Case
### 3.1 Equations of Motion
Assuming we can treat the monopoles as point particles living on the string and there are no unconfined magnetic charges, the motion for a monopole attached to two strings can be described by the action
$$S=m𝑑s\mu \underset{i=1}{\overset{2}{}}_{\sigma _i(t)}^{\sigma _{i+}(t)}𝑑A_i.$$
The first term is the standard action for a relativistic particle of mass $`m`$, the monopole, $`\mu `$ is the mass per unit length of the string and $`A_i`$ are the areas swept by each of the strings attached to the monopole. The integral over the string world-sheet is bounded by the parameters $`\sigma _i(t)`$ and $`\sigma _{i+}(t)`$, the monopole parametric positions on the $`i`$th string. In our convention the string sources are at anti-monopoles and the ends at monopoles (see Fig. 1). Here the number of strings attached is just two but our action can clearly be generalized to any number of strings by adding new string terms in the sum over $`i`$. From (3) one can see the variation of the string part of the action yields the usual Nambu-Goto term and a total divergence. This divergence can be turned into two integrals along the world-lines of the monopoles bounding the string by the 1+1 dimensional analog of Gauss’ theorem ,
$$_a(\sqrt{\gamma }\gamma ^{ab}x_{,b}^\mu \delta x_\mu )d^2\xi =\lambda _a^{}\gamma ^{ab}x_{,b}^\mu \delta x_\mu ^{}𝑑s^{}\lambda _a^+\gamma ^{ab}x_{,b}^\mu \delta x_\mu ^+𝑑s^+,$$
where the superscripts $`+`$ and $``$ refer to monopoles and anti-monopoles respectively. The unit vector $`\lambda _a`$ is orthogonal to the world-line of the monopole and points into the string world-sheet. In external coordinates it can be written
$$\lambda ^\mu (t)=\lambda _a\gamma ^{ab}x_{,b}^\mu =\pm \gamma _m(t)[\dot{\sigma }(t)\dot{x}^\mu (t,\sigma (t))+x^\mu (t,\sigma (t))],$$
where $`\gamma _m(t)=(1\dot{𝐱}_m^2(t))^{1/2}`$ is the Lorentz factor of the monopoles, and the upper sign is for anti-monopoles and the lower for monopoles. These boundary terms can be absorbed into the monopole action to give the equations of motion for the monopoles ,
$$m\frac{d^2x^\nu }{ds^2}=\mu \underset{i=1}{\overset{2}{}}\lambda _i^\nu ,$$
(11)
where the $`i=1,2`$ labels the strings to which the monopole is attached. The $`0`$-component of this equation,
$$m\dot{\gamma }_m(t)=\mu \underset{i=1}{\overset{2}{}}\dot{\sigma }_i(t),$$
(12)
where $`\dot{\sigma }_i(t)`$ is the rate of change of the position of the monopole on the $`i`$th string, is the equation for energy conservation. The monopoles can create string and lose kinetic energy in the process, or annihilate string and gain kinetic energy. The spatial part in turn can be put in the form
$$\ddot{𝐱}_m(t)=\pm \frac{\mu }{m}\gamma _m^3(t)\underset{i=1}{\overset{2}{}}\frac{𝐱_i^{}(\sigma _i(t),t)}{\left|𝐱_i^{}(\sigma _i(t),t)\right|^2}=\frac{\mu }{m}\gamma _m^3(t)\underset{i=1}{\overset{2}{}}\gamma _i(t)\widehat{n}$$
(13)
where $`\gamma _i(t)`$ is the gamma factor of the string at the position of the monopole, $`𝐱_i^{}(\sigma _i(t),t)`$ is the tangent vector of the string at the monopole, and $`\widehat{n}`$ is a unit vector tangent to the string pointing inwards. A similar result was obtained in for the case of strings bounded by monopoles.
The rate of change of monopole position along the string is related to the energy balance of the system. A more useful expression than (12) can be obtained taking the time derivative of the constraint that the monopole must lie on the string, $`𝐱_m(t)=𝐱_i(\sigma (t),t)`$,
$$\dot{𝐱}_m(t)=\frac{d𝐱_i(\sigma _i(t),t)}{dt}=\dot{\sigma }_i(t)𝐱_i^{}(\sigma _i(t),t)+\dot{𝐱}_i(\sigma (t),t).$$
(14)
Using the constraints (4) and (5), (14) can be re-written as
$$\dot{\sigma }_i(t)=\frac{\dot{𝐱}_m(t)𝐱_i^{}(\sigma _i(t),t)}{\left|𝐱_i^{}(\sigma _i(t),t)\right|^2},$$
(15)
which expresses the rate of change of the monopole parametric position along each of the two strings individually.
The equations of motion for the monopoles, (13) and (15), and for the string, (4), (5) and (6), could be solved to give the motion of our system. Except for very special initial conditions, however, the task of finding analytic solutions to this system of equations is hopeless, so we resort to numerical simulation.
### 3.2 The Numerical Algorithm
Although we cannot solve the entire system analytically, we do have an explicit solution for the interior of the string, (7). We can use this solution to give the string motion, taking as input the motion of the monopoles at the ends of the string, which we will simulate. The effect of the string will be to take excitations from the monopole (or anti-monopole) at one end and transport them to affect the anti-monopole (or monopole) at the other. We proceed as follows.
We consider, for simplicity, only a monopole attached to an anti-monopole by a string, and ignore the effects of the strings attached on the other side of either defect. Since the monopole is constrained to live on the boundary $`\sigma _i(t)`$ of the string, its position in terms of the right- and left-movers can be written
$$𝐱_m(t)=𝐱_i(\sigma _i(t),t)=\frac{1}{2}\left[𝐚_i(\sigma _i(t)t)+𝐛_i(\sigma _i(t)+t)\right].$$
(16)
Its velocity is then given by
$$\dot{𝐱}_m(t)=\frac{d𝐱_i(\sigma _i(t),t)}{dt}=\dot{\sigma }_i(t)𝐱_i^{}(\sigma _i(t),t)+\dot{𝐱}_i(\sigma _i(t),t).$$
(17)
which can be expressed in terms of the right- and left-movers,
$$\dot{𝐱}_m(t)=\frac{1}{2}\left\{[\dot{\sigma }_i(t)1]𝐚_i^{}(\sigma _i(t)t)+[\dot{\sigma }_i(t)+1]𝐛_i^{}(\sigma _i(t)+t)\right\}.$$
(18)
For the anti-monopole, we can interpret this expression as giving the value of $`𝐚_i^{}`$, the excitations that travel away from the defect, in terms of $`\dot{𝐱}_m`$, the monopole velocity, and $`𝐛_i^{}`$, the excitations traveling toward the defect. Similarly for the monopole, (18) gives $`𝐛_i^{}`$ in terms of $`\dot{𝐱}_m`$ and $`𝐚_i^{}`$. The outgoing excitations $`𝐚_i^{}(\sigma _i(t)t)`$ and $`𝐛_i^{}(\sigma _i(t)+t)`$ then propagate along the string segment into the future, and eventually reach the other end (see Fig. 2).
The equations for the system, after some algebra, can be cast in the form
$`\dot{\sigma }_i(t)`$ $`=`$ $`{\displaystyle \frac{\dot{𝐱}_m^2(t)+\dot{𝐱}_m(t)𝐚_i^{}(\sigma _i(t)t)}{\dot{𝐱}_m(t)𝐚_i^{}(\sigma _i(t)t)+1}}`$ (19)
$`𝐛_i^{}`$ $`=`$ $`{\displaystyle \frac{2\dot{𝐱}_m[\dot{\sigma }_i(t)1]𝐚_i^{}(\sigma _i(t)t)}{\dot{\sigma }_i(t)+1}}`$ (20)
for the monopole, and
$`\dot{\sigma }_i(t)`$ $`=`$ $`{\displaystyle \frac{\dot{𝐱}_m^2(t)\dot{𝐱}_m(t)𝐛_i^{}(\sigma _i(t)+t)}{\dot{𝐱}_m(t)𝐛_i^{}(\sigma _i(t)+t)1}}`$ (21)
$`𝐚_i^{}`$ $`=`$ $`{\displaystyle \frac{2\dot{𝐱}_m[\dot{\sigma }_i(t)+1]𝐛_i^{}(\sigma _i(t)+t)}{\dot{\sigma }_i(t)1}}`$ (22)
for the anti-monopole. These equations, along with (13) and
$$𝐱_i^{}(\sigma _i(t),t)=\frac{1}{2}[𝐚_i^{}(\sigma _i(t)t)+𝐛_i^{}(\sigma _i(t)+t)]$$
can be used to solve for the motion of any system of this type.
In a typical time-step of our code at, say, an anti-monopole, we have its velocity, $`\dot{𝐱}_m(t)`$, and the values of $`\sigma _1(t)`$ and $`\sigma _2(t)`$, its parametric positions on the strings. We also have a sequence of values of $`𝐛_i^{}`$ emitted at past times from the other end of the string (see Fig. 2). By interpolation, we can find $`𝐛_i^{}(\sigma _i(t)+t)`$ and use it to compute $`\dot{\sigma }_i(t)`$ and $`𝐚_i^{}`$ using (21) and (22). If we do this for both strings we can compute the acceleration according to (13). We can then use the acceleration and the $`\dot{\sigma }_i(t)`$ we just calculated to compute the velocity and the $`\sigma _i`$ at the next time-step (with an appropriate finite differencing scheme). This is what we started with and therefore all we need to continue onto the next timestep. It is, of course, necessary to store enough of the $`𝐚_i^{}`$ and $`𝐛_i^{}`$ to ensure we can perform the appropriate interpolations at every time-step. We use energy conservation, (12), as an independent check on the accuracy of our code (for example see Fig. 5).
The reward for our efforts is that in a numerical computation we can ignore the presence of the strings entirely and simply store the outward going string excitations in the past of the monopole and anti-monopole world-lines to be used when needed.
If one is interested in the string, however, its position or velocity can be readily re-constructed: Starting at one end of the string, one can find the values of $`\sigma _st`$ and $`\sigma _s+t`$ that correspond to a certain point $`\sigma _s`$ on the string, look up $`𝐚^{}`$ and $`𝐛^{}`$, and integrate to get $`𝐚`$ and $`𝐛`$, and thus the string position (see Fig. 3).
### 3.3 Results
A fundamental feature of these systems, first pointed out by Berezinsky and Vilenkin , is the existence of the dimensionless parameter $`\mu L/Nm`$, the ratio of the total string energy to the total monopole energy. Here, $`\mu `$ is the mass per unit length of the string, $`L`$ the invariant length, $`N`$ the number of monopoles and $`m`$ the monopole rest mass. This parameter defines an equivalence class of loops in the sense that loops with the same geometrical form (shape) and the same value of that parameter will evolve in an analogous way. One can easily see this equivalence by looking at the equations of motion for the monopole (11). This is, in fact, a straightforward generalization to our case of the usual conformal invariance of cosmic strings. For example, if one changes the lengths of string on a loop, $`L\lambda L`$, and the monopole masses, $`m\lambda m`$, the value of $`\mu L/Nm`$ remains constant and the evolution will be the same in the sense that after a time $`\lambda t`$ the larger loop will look the same as the smaller one at time $`t`$, remaining, of course, $`\lambda `$ times bigger.
We have evolved a considerable number of different loop configurations with values of the dimensionless parameter $`\mu L/Nm1`$. As was anticipated, one of the most important effects is the non-periodicity of the solutions which has rather notable consequences, making the evolution of these systems very different from that of ordinary cosmic strings. In particular, the non-periodic nature of the solutions allows the system to explore the phase space. Since there are many more ways in which the string can be wiggly than straight it ends up being mostly wiggly; this is a vastly different situation from that of regular cosmic strings where the loops can get trapped in smooth low entropy configurations and never explore the phase space because of the periodicity of their motion.
Here we only show the results for one of the loop configurations we have evolved. However, the result is typical of all our other simulations and the difference between the usual cosmic string evolution and cosmic necklace evolution is very well illustrated.
Figs. 4, 5 and 6 show the evolution of a necklace using as initial conditions a Kibble-Turok loop with values for the loop parameters of $`\kappa =0.7`$ and $`\varphi =\pi /7`$, twelve monopoles and a value for the dimensionless parameter $`\mu L/Nm=5`$. One can see from Fig. 5 that the monopole and string energies go through a series of maxima and minima. Bearing in mind that the motion is not truly periodic we shall nevertheless abuse the language and refer to these as oscillations.
In the early stages of evolution, when the strings are still quite straight, the monopoles can easily gain kinetic energy from the strings. However, after a few oscillations, thirty or so, the strings become so wiggly that the monopoles can no longer gain significant amounts of kinetic energy from them. While the details of the mechanism for wiggle formation on the string elude us it is clear that they arise from the non-linear back-reaction of the monopoles on the string 20, 22. After the wiggles are formed the energy of the monopoles is mostly just their rest mass. We have observed the behaviour described here in all loops in our simulations with values of the dimensionless parameter $`\mu L/Nm1`$.
Although we have not included self-intersections in our code, it seems clear that had we included them the loop would have fragmented well before reaching the wiggly state shown in Fig. 6. While the exact timescale is not obvious the general evolution is unambiguous. Because of the non-periodicity of necklace motion, necklace loops fragment into smaller daughter loops; some containing monopoles and others not. The ones not containing monopoles evolve like ordinary Nambu-Goto string loops: They self-intersect until they find themselves in stable trajectories and then decay by gravitational radiation. The ones that do contain monopoles self-intersect again until only monopole-antimonopole pairs are left on each of the loops. These pairs then radiate small loops and gravitational radiation until they decay and the monopole and anti-monopole annihilate, along the lines discussed in .
## 4 Monopole Motion: The Massless Limit
The starting regime for the necklace network proposed by Berezinsky and Vilenkin is the limit where the string energy is much larger than the monopole energy. Unless there is a sufficiently long damping era following network formation we expect the dynamics to be quite different from what we have described in the previous section.
Because of the equivalence of loops with equal values of the dimensionless parameter $`\mu L/Nm`$ the large string energy limit can be thought of as the limit where the monopole mass $`m`$ is negligible compared with the mass per unit length of the string $`\mu `$. Unfortunately, we can no longer use the equations derived in the previous section because they do not have a well defined $`m0`$ limit.
It turns out to be convenient to think of the monopoles as test particles constrained to live on a 1+1 dimensional dynamical space-time, the string. In this case the monopoles do not back-react on the string and therefore also do not produce small-scale structure so we expect them to get accelerated by the curvature of the string and eventually collide and annihilate.
The motion of the string, being unaffected by the presence of the test particles, is given by some solution of the Nambu-Goto equations of motion. The monopole motion can in turn be derived by minimising the 1+1 action
$$S_m=m𝑑s=m\sqrt{\gamma _{ab}d\xi ^ad\xi ^b}$$
which is the world-line of the monopole on the string world-sheet. One can re-write this action by explicitly evaluating $`\gamma _{ab}`$ in the conformal gauge, yielding
$$S_m=m\sqrt{(1\dot{𝐱}_s^2)dt^2𝐱_s^{}d\sigma ^2}=m𝑑t\sqrt{1\dot{𝐱}_m^2}$$
where
$$\dot{𝐱}_m(t)=\frac{d𝐱_s(\sigma (t),t)}{dt}=\dot{\sigma }(t)𝐱_s^{}+\dot{𝐱}_s$$
is the monopole velocity. It should be noted that this action could have also been derived using the constraint $`𝐱_m(t)=𝐱_s(\sigma (t),t)`$ and taking the action to be the world-line of the monopole in 3+1 Minkowski space instead of the 1+1 world-sheet of the string. We feel, however, that the derivation outlined above makes the degrees of freedom of the monopole motion more manifest. The only degrees of freedom of the monopole are its parametric position $`\sigma (t)`$ and velocity $`\dot{\sigma }(t)`$ along the string.
This action yields the equation for the motion of the monopole along the string which, after some algebra, can be put in the form
$$\ddot{\sigma }(t)=(\dot{\sigma }(t)1)^2\frac{𝐱_s^{}(\sigma ,t)\left\{\ddot{𝐱}_s(\sigma ,t)+\dot{\sigma }(t)\dot{𝐱}_s^{}(\sigma ,t)\right\}}{\left|𝐱_s^{}(\sigma ,t)\right|^2}$$
(23)
This equation can be readily solved by numerically integrating the velocity $`\dot{\sigma }(t)`$ and position $`\sigma (t)`$ along the string of monopoles and anti-monopoles for a number of different loop trajectories.
We have again used a wide variety of initial conditions and found that, as we expected, in all cases the monopoles annihilate with the anti-monopoles within a few oscillation periods. The typical behaviour is well illustrated in Fig. 7 which shows the evolution of the parametric position of twelve monopoles on a loop identical to the one in Fig. 4. All twelve monopoles annihilate before the end the second oscillation period. Although the exact timescale does depend on the particular shape of the loop, in no case have we found it to be longer than a few oscillation periods of the loop. The reason for this is that the motion of the string is unaffected by the presence of the monopoles and is therefore periodic, while the motion of the monopoles given by (23), is not in general periodic. Being constrained to live on a string, it is inevitable that monopoles eventually collide and annihilate.
## 5 Concluding Remarks
We have examined the dynamics of cosmic string loops with monopoles in two limits, when the energy in the string is comparable to the monopole energy and when the string energy is much larger than the monopole energy.
The behaviour in both cases seems to point to the annihilation of the monopoles in the system. When the monopole energy is very small compared with the string energy, the monopoles on the string collide and annihilate in a few oscillation periods. When the monopole energy is large, self intersections chop the string into small loops which subsequently decay by gravitational radiation. If the monopole and string scales are not too different, this process inevitably leads to rapid monopole annihilation. On the other hand if the monopoles are very heavy and the strings very light, loops with one monopole and one anti-monopole may be long-lived, as in . Except for this last possibility, though, it would appear that necklace loops are rather short lived.
However, one could envision an intermediate regime in which the back reaction of the monopoles on the string is not large enough to produce such a rapid process of decay by self intersections but not so small as to decay directly by annihilations. Such a regime could significantly increase the life-time of loops. Unfortunately, at present, we are not able to faithfully simulate such a regime.
Furthermore, in the case when the monopole energy is small compared to the string energy, the possibility exists that monopoles may miss each other upon colliding and avoid annihilation. The likelihood of this would depend on the cross section for monopole-antimonopole annihilation and the thickness of the string. If the probability of crossing is large then having crossed each other they would feel a force towards each other (twice the tension on the string) and undergo a bouncing process. Such a process could in principle extend the lifetime of the monopoles on the string until the strings have radiated enough energy to gravitational waves that the monopole energy becomes comparable to the string energy. At this point the system would decay by self-intersections, which could in principle be happening now. Our results therefore indicate that if the string energy is dominant when the system is formed and there is not sufficient damping to overturn this regime, the scenario for cosmic ray production proposed in is viable only if the cross section for monopole-antimonopole annihilation is sufficiently small.
The purpose of this work was to study the dynamics of cosmic necklaces to use this knowledge in the construction of cosmologically interesting models. We have found that we can understand the behaviour of these systems in two different regimes. A more definite statement regarding the viability of these systems as a source of ultra-high energy cosmic rays and other cosmological consequences of interest depends crucially on the cosmological evolution of the dimensionless parameter $`\mu L/Nm`$. Work in this direction is in progress.
## 6 Acknowledgements
We would like to thank Jose Blanco-Pillado, Christopher Borgers, Carlos Serna, Alexander Vilenkin and Serge Winitzki for useful discussions. The work of K.D.O. was supported in part by the National Science Foundation.
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# 1 Introduction
## 1 Introduction
We have recently proposed new sets of Fock states which can be exploited in the contexts of coherence and squeezing -. Due to the inclusion of a (real continuous) parameter $`\lambda `$ (in the bosonic creation Heisenberg operator), the corresponding oscillatorlike ”Hamiltonians” led to stationary Schrödinger problems characterized by $`\lambda `$-independent eigenvalues but $`\lambda `$-dependent eigenfunctions, the latter ones being particularly interesting for the study of new squeezed states . Moreover, this approach was in a certain sense a kind of deformation of the current one but following Wigner’s point of view .
Let us insist strongly on the fact that we were considering ”squeezing” through the $`\lambda `$-dependent eigenfunctions of our Schrödinger problems and, evidently, through the associated meanvalues of position and momentum, their variances and (in)equalites coming from the Heisenberg relations. This differs from Yuen’s approach which is based on the study of ”squeezing” through the famous two-photon coherent states of the radiation field asking for eigenstates of the oscillator operators and not of the Hamiltonian.
Here we want to generalize such developments by including a priori more than one parameter when, simultaneously, we study new creation and annihilation operators as well as the corresponding oscillatorlike ”Hamiltonians” appearing as physically admissible or nonadmissible ones.
The contents are then distributed as follows. In Section 2, we recall a few relations issued from our first approach . Section 3 is devoted to its generalization already suggested. In Section 4, we apply these considerations to the squeezing problem and find real improvements with respect to the one-parameter previous results. Finally some conclusions and comments are included in Section 5.
## 2 A short survey of our recent proposal
Let us define the new (bosonic) creation Heisenberg operator by
$$a_\lambda ^{}a^{}+\lambda I,\lambda R,$$
(1)
where $`\lambda `$ refers to a real continuous parameter and where $`a^{}`$ is the Hermitian conjugate of the annihilation operator $`a`$ satisfying altogether the expected Heisenberg commutation relations, i.e.
$$[a,a_\lambda ^{}]=I,[a,a]=[a_\lambda ^{},a_\lambda ^{}]=0.$$
(2)
These quantum harmonic oscillatorlike considerations lead to an analog of a (non-Hermitian) Hamiltonian of the type
$$H_\lambda =\frac{1}{2}\{a,a_\lambda ^{}\}=\frac{1}{2}\{a,a^{}\}+\lambda a=H_{H.O.}+\lambda a$$
(3)
where the harmonic oscillator Hamiltonian is obviously given by
$$H_{H.O.}=\frac{1}{2}\frac{d^2}{dx^2}+\frac{1}{2}x^2,H_{H.O.}^{}=H_{H.O.}.$$
(4)
Moreover they ensure that
$$[H_\lambda ,a]=a,[H_\lambda ,a_\lambda ^{}]=a_\lambda ^{},$$
(5)
so that (generalized) Wigner’s approach of quantum mechanics is still working. The energy eigenvalues and eigenfunctions have been determined as
$$E_{n,\lambda }=n+\frac{1}{2}(n=0,1,2,\mathrm{})$$
(6)
and
$$\psi _{n,\lambda }=\frac{2^{\frac{n}{2}}\pi ^{\frac{1}{4}}}{\sqrt{n!}\sqrt{L_n^{(0)}(\lambda ^2)}}e^{\frac{x^2}{2}}H_n(x+\frac{\lambda }{\sqrt{2}})$$
(7)
where, as usual, we have chosen units such that $`\omega =1`$, $`\overline{h}=1`$ and where $`H_n`$ and $`L_n^{(0)}`$ refer to Hermite and generalized Laguerre polynomials , respectively. Let us insist on the unchanged spectrum (6) with respect to well known oscillator results but now with $`\lambda `$-modified eigenfunctions. Moreover we have shown that these new eigenfunctions (7) correspond to specific squeezed states (,). Let us recall that squeezed states have already been experimentally detected being seen as “two-photon coherent states”for the electromagnetic field. Our new states lead to the characteristic inequality for squeezing given by
$$(\mathrm{\Delta }x)_\lambda ^2=2n+\frac{1}{2}(2\lambda ^2+1)\frac{L_{n1}^{(1)}(\lambda ^2)}{L_n^{(0)}(\lambda ^2)}2\lambda ^2(\frac{L_n^{(1)}(\lambda ^2)}{L_n^{(0)}(\lambda ^2)})^2<\frac{1}{2}$$
(8)
if
$$n=1,2,3,\mathrm{}and\lambda R\backslash ]r,+r[,r0ifn\mathrm{}.$$
(9)
## 3 A simple way to get generalized developments
The qualities and defects of our above approach suggest the following new position of the problem:
to search for (bosonic) oscillatorlike annihilation ($`b`$) and creation ($`b^+`$) operators ensuring the following conditions
$$[b,b^+]=1,[H,b]=b,[H,b^+]=b^+$$
(10)
and
$$H=\frac{1}{2}\{b,b^+\}=\alpha \frac{d^2}{dx^2}+\beta (x)\frac{d}{dx}+\gamma (x),$$
(11)
where $`b`$ and $`b^+`$ have to be general expressions of the usual operators $`a`$ and $`a^{}`$. Let us note that we have introduced different notations for the usual Hermitian conjugate operator $`a^{}`$ of $`a`$ and the so-called $`b^+`$ associated to $`b`$ with a general meaning discussed in the following.
Such a set of conditions obviously contains the Heisenberg and Wigner requirements through eqs. (10) and, moreover, restricts the Hamiltonian to Schrödingerlike ones through eq. (11) where $`\alpha `$ is a real constant and $`\beta `$, $`\gamma `$ are arbitrary real functions of the space variable.
According to such a programme, let us introduce the generalized operators $`b`$ and $`b^+`$ in terms of $`a`$ and $`a^{}`$ by the following definitions
$$b=(1+c_1)a+c_2a^{}+c_3$$
(12)
and
$$b^+=c_4a+(1+c_5)a^{}+c_6,$$
(13)
where $`c_1,c_2,\mathrm{},c_6`$ are arbitrary (real) parameters and where the current harmonic oscillator context has been included by equating all the parameters to zero while the definition (1) is also incorporated in equating only $`c_6`$ with the $`\lambda `$-parameter, all the other $`c`$’s being identically zero.
At this stage let us point out that the generalized operators (12) and (13) are intimately connected with the construction of the so-called ”two-photon coherent states” due to Yuen but here from inhomogeneous linear canonical transformations which could be summarized by the following form
$$\left(\begin{array}{ccc}b& & \\ b^+& & \end{array}\right)=\left(\begin{array}{ccc}1+c_1& c_2& \\ c_4& 1+c_5& \end{array}\right)\left(\begin{array}{ccc}a& & \\ a^{}& & \end{array}\right)+\left(\begin{array}{ccc}c_3& & \\ c_6& & \end{array}\right)$$
where the matrix has to be invertible so that, for example,
$$(1+c_1)(1+c_5)c_2c_4=1.$$
This leads to
$$c_1+c_5+c_1c_5c_2c_4=0$$
(14)
which is the constraint between the $`c`$’s issued from (10) and leaving, in fact, only five independent parameters in the whole discussion.
By taking care of the definitions (12) and (13) in the Hamiltonian (11) and by remembering that
$$a=\frac{1}{\sqrt{2}}(\frac{d}{dx}+x),a^{}=\frac{1}{\sqrt{2}}(\frac{d}{dx}+x),$$
(15)
the possible Hamiltonians are then found on the following form
$$H=A\frac{d^2}{dx^2}+(Bx+C)\frac{d}{dx}+Dx^2+Ex+F$$
(16)
transferring the parametrization on the six parameters $`A,B,\mathrm{},F`$ given by
$`A={\displaystyle \frac{1}{2}}c_2c_4+{\displaystyle \frac{1}{2}}c_4(1+c_1)+{\displaystyle \frac{1}{2}}c_2(1+c_5),`$
$`B=c_4(1+c_1)c_2(1+c_5),`$
$`C={\displaystyle \frac{1}{\sqrt{2}}}[c_6(c_1c_2+1)+c_3(c_4c_51)],`$
$`D={\displaystyle \frac{1}{2}}+c_2c_4+{\displaystyle \frac{1}{2}}c_4(1+c_1)+{\displaystyle \frac{1}{2}}c_2(1+c_5),`$ (17)
$`E={\displaystyle \frac{1}{\sqrt{2}}}[c_6(c_1+c_2+1)+c_3(c_4+c_5+1)],`$
$`F={\displaystyle \frac{1}{2}}c_4(1+c_1){\displaystyle \frac{1}{2}}c_2(1+c_5)+c_3c_6.`$
These developments compared to the previous ones clearly appear as a generalization; moreover it permits an interesting discussion at the level of physically admissible Hamiltonians as well as at the level of coherence and (or) squeezing, once we have solved the eigenvalue and eigenfunction problems associated with such Hamiltonians.
In terms of the new parameters, let us point out again that the current harmonic oscillator Hamiltonian (4) corresponds to
$$A=D=\frac{1}{2},B=C=E=F=0,$$
(18)
while our previous deformation (1) leading to the Hamiltonian (3) is given by
$$A=D=\frac{1}{2},B=F=0,C=E=\frac{\lambda }{\sqrt{2}}.$$
(19)
With the Hamiltonian (16) and the relations (17), the stationary Schrö-
dinger problem can now be solved by conventional quantum mechanical methods . It leads to the general answer
$$E_n=F\frac{B}{2}\frac{C^2}{4A}\frac{A}{p^2}(2n+1)+q^2(D\frac{B^2}{4A})+q(E\frac{BC}{2A})$$
(20)
while the corresponding eigenfunctions take the form
$$\psi _n(x)=exp[\frac{B}{4A}x^2\frac{C}{2A}x]exp[\frac{x^2}{2p^2}+\frac{qx}{p^2}\frac{q^2}{2p^2}]H_n(\frac{xq}{p}),$$
(21)
where $`p`$ and $`q`$ enter the necessary change of variable
$$x=py+q,(p0).$$
(22)
Let us mention the two constraints
$$\frac{p^4}{A}(D\frac{B^2}{4A})=1$$
(23)
and
$$2q(D\frac{B^2}{4A})+E\frac{BC}{2A}=0,$$
(24)
issued from these calculations.Together with eqs. (17), these relations (23), (24) fix the parameters $`p`$ and $`q`$ of our change of variable (22) to the unique values:
$$p^2=2A,q=2EABC$$
(25)
in order to get in particular a positive spectrum. By requiring to deal with square integrable eigenfunctions, we finally have to ask for
$$A<0andB<1$$
(26)
compatible with the specific cases (18) and (19). We thus get (up to a normalization factor $`N_n`$) the solutions (21) as given by
$`\psi _n(x)=N_nexp[{\displaystyle \frac{1B}{4A}}x^2]exp[({\displaystyle \frac{(B1)C}{2A}}E)x]exp[{\displaystyle \frac{(2EABC)^2}{4A}}]`$
$`H_n({\displaystyle \frac{1}{\sqrt{2A}}}(x2EA+BC)).`$ (27)
Moreover we obtain the remarkable result of an unchanged real spectrum
$$E_n=n+\frac{1}{2},n=0,1,2,\mathrm{},$$
(28)
inside this general context as it was already the case in our first study (see (6)). Let us here insist on this real character without having required the selfadjointness of the Hamiltonian (16), a similar property to the one which has recently been quoted by Bender and Boettcher although we have not required any specific discrete symmetries. Nevertheless, we have also noticed that a necessary and sufficient condition ensuring that $`H^{}=H`$ is simply
$$B=C=0$$
(29)
leading to a large class of physically admissible Hamiltonians.
As a last remark in this Section, let us point out that such eigenfunctions $`\psi _n(x)`$ like (27) are once again associated with Fock states - let us call them $`n>_c`$ referring to the $`c`$-parametrization included in eqs.(12) and (13) - and it is interesting to quote the action of $`b`$ and $`b^{}`$ on such states. We obtain
$$bn>_c=\frac{n}{\sqrt{A}}\frac{N_n}{N_{n1}}(1+c_1c_2)n1>_c$$
(30)
and
$$b^+n>_c=\frac{1}{2\sqrt{A}}\frac{N_n}{N_{n+1}}(1+c_5c_4)n+1>_c$$
(31)
and point out that
$$bb^+n>_c=(n+1)n>_c,b^+bn>_c=nn>_c$$
(32)
so that the conditions (10) and (11) are obviously satisfied, ensuring in particular that
$$\{b,b^+\}n>_c=2Hn>_c=(2n+1)n>_c.$$
(33)
## 4 On implications in squeezing
Let us, first, extract new information by considering the lowest energy eigenvalue $`E_0`$ of the spectrum and its associated eigenfunction $`\psi _0(x)`$. Then, as a second step, let us come back very briefly on known cases corresponding to the Hamiltonian (16) with the conditions (18) or (19). Finally, let us consider thirdly new parametrizations exploiting the results obtained in Section 3 mainly with a view of interesting improvements in squeezing.
### 4.1 From the lowest eigenvalue of the spectrum
Due to the fundamental and specific role played by the lowest energy eigenvalue $`E_0`$, let us study coherence and squeezing through the eigenfunction $`\psi _0(x)(27)`$ which takes the explicit form
$$\psi _0(x)=N_0exp[\frac{1}{2}\alpha x^2\beta x\frac{1}{2}\gamma ],N_0=(\frac{\alpha }{\pi })^{\frac{1}{4}}exp[\frac{\alpha \gamma \beta ^2}{2\alpha }]$$
(34)
in order to ensure that
$$_{\mathrm{}}^+\mathrm{}\psi _0^2(x)𝑑x=1$$
(35)
where, for brevity, we have introduced the notations
$$\alpha =\frac{B1}{2A},\beta =E\frac{(B1)C}{2A},\gamma =\frac{1}{2A}(2EABC)^2.$$
(36)
In such a $`n=0`$-context, the meanvalues and consequences are readily obtained as follows:
$`x_0={\displaystyle \frac{2EA+C}{1B}},x^2_0={\displaystyle \frac{A}{B1}}+({\displaystyle \frac{2EA}{B1}}C)^2,`$
$`p_0=0,p^2_0={\displaystyle \frac{B1}{4A}},`$ (37)
so that we get
$$(\mathrm{\Delta }x)_0^2=\frac{A}{B1},(\mathrm{\Delta }p)_0^2=\frac{B1}{4A}.$$
(38)
These results ensure coherence due to the Heisenberg relation
$$(\mathrm{\Delta }x)_0(\mathrm{\Delta }p)_0=\frac{1}{2}$$
(39)
and squeezing on the x-variable
$$(\mathrm{\Delta }x)_0^2<\frac{1}{2}iffB<2A+1$$
(40)
or on the p-variable
$$(\mathrm{\Delta }p)_0^2<\frac{1}{2}iffB>2A+1.$$
(41)
Such inequalities on $`A`$ and $`B`$ only will suggest our future parametrizations in Sections 4.2 and 4.3. In fact, let us immediately inform the reader that we plan to priviledge the discussion on the x-variable so that eq. (40) will play the main role.
### 4.2 From known cases
(i) The harmonic oscillator context characterized by the condition (18) is obviously well known as far as coherence and squeezing are visited (-). As already mentioned, this case is contained in our study but we learn only that it corresponds to all $`c`$’s equal to zero in eqs. (12) and (13), it generates a selfadjoint Hamiltonian (4) and deals with hermitian conjugated operators $`ba`$ and $`b^{}a^{}`$ verifying the condition
$$(b^{})^{}=b.$$
(42)
(ii) The deformed context characterized by the condition (19) has already been discussed in : it breaks down the condition (42) and the selfadjointness of the Hamiltonian (3) so that physical connections are here questionable although they correspond to real spectra and to new possibilities of squeezing for $`n0`$ . Let us point out that the conditions (29) are obviously in contradiction with eqs. (19) and that, for $`n=0`$, the inequalities (40) cannot be satisfied.
### 4.3 To new contexts
By keeping the conditions (29) in order to maintain the selfadjointness of the Hamiltonians, we can also require the condition (42). The latter leads to very simple demands of the types:
$$c_1=c_5,c_2=c_4,c_3=c_6$$
(43)
so that we then get families of physically admissible Hamiltonians which can be further exploited.
(i) Within such conditions, let us go to a one-parameter $`\lambda `$-deformation with, for example, the values
$$c_1=c_5=\frac{2}{3},c_2=c_4=\frac{4}{3},c_3=c_6=\lambda .$$
(44)
Such a case corresponds to the parametrization (17) given by
$$A=\frac{1}{18},C=\frac{9}{2},E=3\sqrt{2}\lambda ,F=\lambda ^2,B=C=0$$
(45)
and the eigenvalues and eigenfunctions problem can be completely solved. We evidently get the spectrum (28) and the eigenfunctions (27) take the final form
$$\psi _n(x)=N_nexp[\frac{9}{2}x^2\frac{6}{\sqrt{2}}\lambda x\lambda ^2]H_n(3x+\sqrt{2}\lambda )$$
(46)
where the normalization factor is found on the following form
$$N_n=\frac{\sqrt{3}\pi ^{\frac{1}{4}}2^{\frac{n}{2}}}{\sqrt{n!}}.$$
(47)
Meanvalues and Heisenberg constraints can then be evaluated and we get
$$x_\lambda =\frac{\sqrt{2}}{3}\lambda ,x^2_\lambda =\frac{1}{9}(2\lambda ^2+n+\frac{1}{2})$$
(48)
and
$$p_\lambda =0,p^2_\lambda =9(n+\frac{1}{2}),$$
(49)
so that
$$(\mathrm{\Delta }x)_\lambda ^2=\frac{1}{9}(n+\frac{1}{2}),(\mathrm{\Delta }p)_\lambda ^2=9(n+\frac{1}{2})$$
(50)
leading to
$$(\mathrm{\Delta }x)_\lambda (\mathrm{\Delta }p)_\lambda =n+\frac{1}{2}.$$
(51)
This result is analogous to the one of the undeformed case but, here, it permits, moreover, squeezing (but on x only) due to the relations (40) and (50). In fact, such a squeezing can only take place for $`n=0,1,2,3.`$
(ii) A final improvement of this example consists in the possible increase of such $`n`$-values permitting the squeezing and maintaining the nice property (42) and the selfadjointness of $`H`$. This can be realized through the new $`\lambda `$-deformation ($`\lambda >0`$) characterized by the values
$$c_1=c_5=\frac{(\sqrt{\lambda }1)^2}{2\sqrt{\lambda }},c_2=c_4=\frac{\lambda 1}{2\sqrt{\lambda }},c_3=c_6=0,$$
(52)
leading to the relations (17) given now on the form
$$A=\frac{1}{2\lambda },D=\frac{\lambda }{2},B=C=E=F=0.$$
(53)
satisfying once again the inequalities (40) when $`n=0`$.
Here the eigenfunctions are found as
$$\psi _n(x)=N_nexp[\frac{\lambda }{2}x^2]H_n(\sqrt{\lambda }x),N_n=\frac{\lambda ^{\frac{1}{4}}\pi ^{\frac{1}{4}}2^{\frac{n}{2}}}{\sqrt{n!}}$$
(54)
and we get in correspondence with eqs. (50)
$$(\mathrm{\Delta }x)_\lambda ^2=\frac{1}{\lambda }(n+\frac{1}{2}),(\mathrm{\Delta }p)_\lambda ^2=\lambda (n+\frac{1}{2}).$$
(55)
We thus notice once more the validity of eq. (51) ensuring coherence for the particular value $`n=0`$ only but squeezing (in the $`x`$-coordinate) for all the values $`n`$ satisfying the following inequality
$$\lambda >2n+1>0.$$
(56)
A further interesting property of the above eigenfunctions (54) (and evidently (46)) is that, due to the characteristics of Hermite polynomials , these solutions are not only normalized but are also orthogonal as it can be easily established.
If physical applications require a fixed finite set of levels in the energy spectrum, we can always choose, due to the inequality (56), our $`\lambda `$-parameter in order to guarantee the squeezing up to this $`n`$-value.
(iii) As a last context, let us relax the condition (42) and the selfadjointness of the Hamiltonian. This corresponds to an extension of the context discussed in and recalled here in the subsection (4.2.ii). We can choose, for example,
$$b=a+\lambda a^{},b^+=a^{}$$
(57)
corresponding to all the null parameters $`c`$ except $`c_2=\lambda `$ or to
$$A=\frac{1}{2}(\lambda 1),B=\lambda ,C=E=0,D=\frac{1}{2}(\lambda +1),F=\frac{\lambda }{2}$$
(58)
ensuring squeezing on $`x`$ in the $`n=0`$-case if $`1>\lambda >0`$. With the spectrum (28), the associated eigenfunctions here take the form
$$\psi _n(x)=N_nexp[\frac{1}{2}(\frac{1+\lambda }{1\lambda })x^2]H_n(\frac{x}{\sqrt{1\lambda }}).$$
(59)
They are normalizable with
$$N_n=\frac{\pi ^{\frac{1}{4}}}{n!}(\frac{1+\lambda }{1\lambda })^{\frac{1}{4}}(1+\lambda )^{\frac{n}{2}}F_n^{\frac{1}{2}}(\lambda )$$
(60)
but not orthogonal. Depending on the even or odd character of $`n`$ the functions $`F_n(\lambda )`$ are respectively given by
$$F_n(\lambda )=\underset{l=0}{\overset{\frac{n}{2}}{}}\frac{2^{2l}\lambda ^{n2l}}{(2l)![(\frac{n}{2}l)!]^2},(neven),$$
(61)
or
$$F_n(\lambda )=\underset{l=0}{\overset{\frac{n1}{2}}{}}\frac{2^{2l+1}\lambda ^{n12l}}{(2l+1)![(\frac{n1}{2}l)!]^2},(nodd).$$
(62)
These functions enter the evaluation of meanvalues and Heisenberg constraints for each $`n`$-value. Specific values are of interest in order to learn the general behaviour of the corresponding meanvalues and their consequences but these are only exercises. Let us just point out that, for $`n=0`$, we get
$`x_\lambda =0,x^2_\lambda ={\displaystyle \frac{1}{2}}({\displaystyle \frac{1\lambda }{1+\lambda }})=(\mathrm{\Delta }x)_\lambda ^2`$
$`p_\lambda =0,p^2_\lambda ={\displaystyle \frac{1}{2}}({\displaystyle \frac{1+\lambda }{1\lambda }})=(\mathrm{\Delta }p)_\lambda ^2`$ (63)
giving us coherence due to
$$(\mathrm{\Delta }x)_\lambda ^2(\mathrm{\Delta }p)_\lambda ^2=\frac{1}{4},\lambda ,$$
(64)
while squeezing requires parametrizations according to
$$1>\lambda >0or1<\lambda <0$$
(65)
in the $`x`$\- or $`p`$\- context respectively. Coherence is then lost if $`n0`$ but squeezing can be installed when specific refined inequalities of the type (65) are valid. The upper and lower bounds on these $`\lambda `$-values can be determined by entering the results (59)-(62) depending on the $`n`$-values we are considering.
## 5 Some further conclusions and comments
Among the above results, let us point out those obtained more particularly in the subsection (4.3.ii) leading to an attracting class of one-parameter selfadjoint Hamiltonians
$$H_\lambda =\frac{1}{2}\{b_\lambda ,b_\lambda ^{}\},$$
(66)
with
$$b_\lambda =(1+\frac{(\sqrt{\lambda }1)^2}{2\sqrt{\lambda }})a+\frac{\lambda 1}{2\sqrt{\lambda }}a^{},b_\lambda ^{}=b^+,(b_\lambda ^{})^{}=b_\lambda ,$$
(67)
characterized by a deformation parameter $`\lambda >0`$ and corresponding to the current harmonic oscillator case when $`\lambda =1`$. Appearing nearly as a trivial result, this family opens possible new studies of squeezing through energy eigenfunctions (54) which are not only normalizable but also orthogonal among themselves. The possible choice $`\lambda >2n+1`$ with a fixed set of energy eigenvalues given by the usual spectrum (28) is maybe an interesting connection with possible experimental realizations for oscillatorlike systems in order to test and to realize the associated squeezed states.
One further comment is the possible exploitation of our generalized operators $`b`$ and $`b^+`$ by studying more than one (real) parameter in the definitions (12) and (13) as just noticed elsewhere .
Another point which has to be recalled is that the motivations of deforming our (annihilation and creation) oscillatorlike operators (as realized in eqs. (12) and (13)) were intimately connected with a specific mathematical property called “subnormality of operators” , a property already exploited in our previous letter .
As a final comment, let us also recall that our developments could evidently be extended to the fermionic sector as already specified in our first approach . The generalizations (12) and (13) can be realized on fermionic annihilation and creation operators and their consequences can then be deduced. Then the superposition of these bosonic and fermionic contexts could be considered in order to go towards supersymmetric developments by including simultaneously specific physical as well as mathematical properties.
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# The Derived Picard Group is a Locally Algebraic Group
(Date: 23.12.00)
## Abstract.
Let $`A`$ be a finite dimensional algebra over an algebraically closed field $`𝕂`$. The derived Picard group $`\mathrm{DPic}_𝕂(A)`$ is the group of two-sided tilting complexes over $`A`$ modulo isomorphism. We prove that $`\mathrm{DPic}_𝕂(A)`$ is a locally algebraic group, and its identity component is $`\mathrm{Out}_𝕂^0(A)`$. If $`B`$ is a derived Morita equivalent algebra then $`\mathrm{DPic}_𝕂(A)\mathrm{DPic}_𝕂(B)`$ as locally algebraic groups. Our results extend, and are based on, work of Huisgen-Zimmermann, Saorín and Rouquier.
Partially supported by a grant from the US-Israel Binational Science Foundation
Let $`A`$ and $`B`$ be associative algebras with $`1`$ over a field $`𝕂`$. We denote by $`𝖣^\mathrm{b}(\mathrm{𝖬𝗈𝖽}A)`$ the bounded derived category of left $`A`$-modules. Let $`B^{}`$ be the opposite algebra, so an $`A_𝕂B^{}`$-module is a $`𝕂`$-central $`A`$-$`B`$-bimodule. A two-sided tilting complex over $`(A,B)`$ is a complex $`T𝖣^\mathrm{b}(\mathrm{𝖬𝗈𝖽}A_𝕂B^{})`$ such there exists a complex $`T^{}𝖣^\mathrm{b}(\mathrm{𝖬𝗈𝖽}B_𝕂A^{})`$ and isomorphisms of the derived tensor products $`T_B^\mathrm{L}T^{}A`$ and $`T^{}_A^\mathrm{L}TB`$. Two-sided tilting complexes were introduced by Rickard in \[Rd\].
When $`B=A`$ we write $`A^\mathrm{e}:=A_𝕂A^{}`$. The set
$$\mathrm{DPic}_𝕂(A):=\frac{\{\text{two-sided tilting complexes }T𝖣^\mathrm{b}(\mathrm{𝖬𝗈𝖽}A^\mathrm{e})\}}{\text{isomorphism}}$$
is the derived Picard group of $`A`$ (relative to $`𝕂`$). The identity element is the class of $`A`$, the multiplication is $`(T_1,T_2)T_1_A^\mathrm{L}T_2`$, and the inverse is $`TT^{}=\mathrm{RHom}_A(T,A)`$.
Denote by $`\mathrm{Out}_𝕂(A)`$ the group of outer $`𝕂`$-algebra automorphism of $`A`$, and by $`\mathrm{Pic}_𝕂(A)`$ the Picard group of $`A`$ (the group of invertible bimodules modulo isomorphism). Then there are inclusions
$$\mathrm{Out}_𝕂(A)\mathrm{Pic}_𝕂(A)\mathrm{DPic}_𝕂(A).$$
The first inclusion sends the automorphism $`\sigma `$ to the invertible bimodule $`A^\sigma `$ where the right action is twisted by $`\sigma `$. The second inclusion corresponds to the full embedding $`\mathrm{𝖬𝗈𝖽}A^\mathrm{e}𝖣^\mathrm{b}(\mathrm{𝖬𝗈𝖽}A^\mathrm{e})`$. See \[Ye\] for details.
To simplify notation we use the same symbol to denote an automorphism $`\sigma \mathrm{Aut}_𝕂(A)`$ and its class in $`\mathrm{Out}_𝕂(A)`$. Likewise for a two-sided tilting complex $`T`$ and its class in $`\mathrm{DPic}_𝕂(A)`$. The precise meaning is always clear from the context.
Now assume $`𝕂`$ is algebraically closed and $`A`$ is a finite dimensional $`𝕂`$-algebra. Then the group $`\mathrm{Aut}_𝕂(A)=\mathrm{Aut}_{\mathrm{𝖠𝗅𝗀}𝕂}(A)`$ of $`𝕂`$-algebra automorphisms is a linear algebraic group, being a closed subgroup of $`\mathrm{GL}(A)=\mathrm{Aut}_{\mathrm{𝖬𝗈𝖽}𝕂}(A)`$. This induces a structure of linear algebraic group on the quotient $`\mathrm{Out}_𝕂(A)`$. Denote by $`\mathrm{Out}_𝕂^0(A)`$ the identity component.
Examples calculated in \[MY\] indicated that the whole group $`\mathrm{DPic}_𝕂(A)`$ should carry a geometric structure (cf. Example 3 below). This is our first main result Theorem 2.
A result of Brauer says that the group $`\mathrm{Out}_𝕂^0(A)`$ is a Morita invariant of $`A`$: if $`A`$ and $`B`$ are Morita equivalent $`𝕂`$-algebras then $`\mathrm{Out}_𝕂^0(A)\mathrm{Out}_𝕂^0(B)`$. In \[HS\] and \[Ro\] this is extended to derived Morita equivalence. Our Theorem 4 extends these results further.
We shall need the following variant of the result of Huisgen-Zimmermann, Saorín and Rouquier.
###### Theorem 1.
Let $`A`$ and $`B`$ be finite dimensional $`𝕂`$-algebras. Suppose $`T`$ $`𝖣^\mathrm{b}(\mathrm{𝖬𝗈𝖽}A_𝕂B^{})`$ is a two-sided tilting complex over $`(A,B)`$, with inverse $`T^{}𝖣^\mathrm{b}(\mathrm{𝖬𝗈𝖽}B_𝕂A^{})`$. Then for any element $`\sigma \mathrm{Out}_𝕂^0(A)`$ the two-sided tilting complex
$$\varphi _T^0(\sigma ):=T^{}_A^\mathrm{L}A^\sigma _A^\mathrm{L}T\mathrm{DPic}_𝕂(B)$$
is in $`\mathrm{Out}_𝕂^0(B)`$. The group homomorphism
$$\varphi _T^0:\mathrm{Out}_𝕂^0(A)\mathrm{Out}_𝕂^0(B)$$
is an isomorphism of algebraic groups.
###### Proof.
According to \[HS, Theorem 17\] or \[Ro, Théorème 4.2\] there is an isomorphism of algebraic groups $`\varphi ^0:\mathrm{Out}_𝕂^0(A)\mathrm{Out}_𝕂^0(B)`$ induced by $`T`$. Letting $`\tau :=\varphi ^0(\sigma )\mathrm{Out}_𝕂^0(B)`$ one has
$$T_BB^\tau A^\sigma _AT\text{ in }𝖣(\mathrm{𝖬𝗈𝖽}A_𝕂B^{}).$$
Applying $`T^{}_A^\mathrm{L}`$ to this isomorphism we see that $`B^\tau \varphi _T^0(\sigma )`$ in $`𝖣(\mathrm{𝖬𝗈𝖽}B^\mathrm{e})`$, so $`\tau =\varphi _T^0(\sigma )`$ in $`\mathrm{DPic}_𝕂(B)`$. We conclude that $`\varphi _T^0=\varphi ^0`$. ∎
A locally algebraic group over $`𝕂`$ is a group $`G`$, with a normal subgroup $`G^0`$, such that $`G^0`$ is a connected algebraic group over $`𝕂`$, each coset of $`G^0`$ is a variety, and multiplication and inversion are morphisms of varieties. A morphism $`\varphi :GH`$ of locally algebraic groups is a group homomorphism such that $`\varphi (G^0)H^0`$ and the restriction $`\varphi ^0:G^0H^0`$ is a morphism of varieties. We call $`\varphi `$ an open immersion if $`\varphi `$ is injective and $`\varphi ^0`$ is an isomorphism.
In other words $`G`$ is the group of rational points $`𝑮(𝕂)`$ of a reduced group scheme $`𝑮`$ locally of finite type over $`𝕂`$, in the sense of \[SGA3, Exposé VI<sub>A</sub>\]. A morphism $`\varphi :GH`$ corresponds to a morphism $`\mathit{\varphi }:𝑮𝑯`$ of group schemes over $`𝕂`$.
Here is our first main result.
###### Theorem 2.
Let $`A`$ be a finite dimensional $`𝕂`$-algebra. Then the derived Picard group $`\mathrm{DPic}_𝕂(A)`$ is a locally algebraic group over $`𝕂`$. The inclusion $`\mathrm{Out}_𝕂(A)\mathrm{DPic}_𝕂(A)`$ is an open immersion.
In particular the identity components coincide: $`\mathrm{Out}_𝕂^0(A)=\mathrm{DPic}_𝕂^0(A)`$.
###### Proof.
Theorem 1 with $`A=B`$ implies that the subgroup $`\mathrm{Out}_𝕂^0(A)\mathrm{DPic}_𝕂(A)`$ is normal, and for any two-sided tilting complex $`T`$ the conjugation $`\varphi _T^0:\mathrm{Out}_𝕂^0(A)\mathrm{Out}_𝕂^0(A)`$ is an automorphism of algebraic groups.
Let us now switch to the notation $`T_1T_2`$ and $`T^1`$ for the operations in $`\mathrm{DPic}_𝕂(A)`$. Define an algebraic variety structure on each coset $`C=T\mathrm{Out}_𝕂^0(A)\mathrm{DPic}_𝕂(A)`$ using the multiplication map $`PTP`$, $`P\mathrm{Out}_𝕂^0(A)`$. Since $`\varphi _T^0`$ is an automorphism of algebraic groups, the variety structure is independent of the representative $`TC`$.
Let us prove that $`\mathrm{DPic}_𝕂(A)`$ is a locally algebraic group. For $`P_1,P_2\mathrm{Out}_𝕂^0(A)`$ and $`T_1,T_2\mathrm{DPic}_𝕂(A)`$, multiplication is the morphism
$$(T_1P_1)(T_2P_2)=(T_1T_2)(\varphi _{T_2}^0(P_1)P_2).$$
Similarly for the inverse:
$$(TP)^1=T^1\varphi _T^0(P)^1.$$
###### Example 3.
Let $`\stackrel{}{\mathrm{\Omega }}_n`$ be the quiver with two vertices $`x,y`$ and $`n`$ arrows $`x\stackrel{\alpha _i}{}y`$. Let $`A`$ be the path algebra $`𝕂\stackrel{}{\mathrm{\Omega }}_n`$. According to \[MY, Theorem 5.3\], $`\mathrm{Out}_𝕂(A)\mathrm{Pic}_𝕂(A)\mathrm{PGL}_n(𝕂)`$ and
$$\mathrm{DPic}_𝕂(A)\times \left(\mathrm{PGL}_n(𝕂)\right).$$
In the semi-direct product a generator $`T`$ of $``$ acts on a matrix $`\sigma \mathrm{PGL}_n(𝕂)`$ by $`\varphi _T^0(\sigma )=(\sigma ^1)^\mathrm{t}`$. This is clearly a morphism of varieties, so $`\mathrm{DPic}_𝕂(A)`$ is indeed a locally algebraic group.
Our second main result relates two algebras. Recall that the algebras $`A`$ and $`B`$ are derived Morita equivalent over $`𝕂`$ if there is a $`𝕂`$-linear equivalence of triangulated categories $`𝖣^\mathrm{b}(\mathrm{𝖬𝗈𝖽}A)𝖣^\mathrm{b}(\mathrm{𝖬𝗈𝖽}B)`$.
###### Theorem 4.
Suppose $`A`$ and $`B`$ are two finite dimensional $`𝕂`$-algebras, and assume they are derived Morita equivalent over $`𝕂`$. Then $`\mathrm{DPic}_𝕂(A)\mathrm{DPic}_𝕂(B)`$ as locally algebraic groups.
###### Proof.
It is known that there exist two-sided tilting complexes $`T𝖣(\mathrm{𝖬𝗈𝖽}A_𝕂B^{})`$; choose one. We obtain a group isomorphism
$$\varphi _T:\{\begin{array}{cc}\mathrm{DPic}_𝕂(A)\mathrm{DPic}_𝕂(B),\hfill & \\ ST^{}_A^\mathrm{L}S_A^\mathrm{L}T.\hfill & \end{array}$$
By Theorem 1, $`\varphi _T`$ restricts to an isomorphism of algebraic groups $`\varphi _T^0:\mathrm{Out}_𝕂^0(A)\mathrm{Out}_𝕂^0(B)`$. So $`\varphi _T`$ is an isomorphism of locally algebraic groups. ∎
We end the paper with a corollary and some remarks. Suppose $`𝖢`$ is a $`𝕂`$-linear triangulated category that’s equivalent to a small category. Denote by $`\mathrm{Out}_𝕂^{\mathrm{tr}}(𝖢)`$ the group of $`𝕂`$-linear triangle auto-equivalences of $`𝖢`$ modulo natural isomorphism. Let $`\mathrm{𝗆𝗈𝖽}A`$ stand for the category of finitely generated $`A`$-modules.
###### Corollary 5.
Suppose $`𝖢`$ is a $`𝕂`$-linear triangulated category that is equivalent to $`𝖣^\mathrm{b}(\mathrm{𝗆𝗈𝖽}A)`$ for some hereditary finite dimensional $`𝕂`$-algebra $`A`$. Then $`\mathrm{Out}_𝕂^{\mathrm{tr}}(𝖢)`$ is a locally algebraic group.
###### Proof.
Trivially $`\mathrm{Out}_𝕂^{\mathrm{tr}}(𝖢)\mathrm{Out}_𝕂^{\mathrm{tr}}(𝖣^\mathrm{b}(\mathrm{𝗆𝗈𝖽}A))`$, and by \[MY, Corollary 0.11\] we have $`\mathrm{Out}_𝕂^{\mathrm{tr}}(𝖣^\mathrm{b}(\mathrm{𝗆𝗈𝖽}A))\mathrm{DPic}_𝕂(A)`$. ∎
###### Example 6.
Beilinson \[Be\] proved that $`𝖣^\mathrm{b}(\mathrm{𝖢𝗈𝗁}𝐏_𝕂^1)𝖣^\mathrm{b}(\mathrm{𝗆𝗈𝖽}𝕂\stackrel{}{\mathrm{\Omega }}_2)`$, where $`\mathrm{𝖢𝗈𝗁}𝐏_𝕂^1`$ is the category of coherent sheaves on the projective line, and $`\stackrel{}{\mathrm{\Omega }}_2`$ is the quiver from Example 3. Therefore $`\mathrm{Out}_𝕂^{\mathrm{tr}}(𝖣^\mathrm{b}(\mathrm{𝖢𝗈𝗁}𝐏_𝕂^1))`$ is a locally algebraic group. This should be compared to Remark 7 below; see also \[MY, Remark 5.4\].
###### Remark 7.
Suppose $`X`$ is a smooth projective variety over $`𝕂`$ with ample canonical or anti-canonical bundle. Bondal and Orlov \[BO\] prove that
$$\mathrm{Out}_𝕂^{\mathrm{tr}}(𝖣^\mathrm{b}(\mathrm{𝖢𝗈𝗁}X))\left(\mathrm{Aut}_𝕂(X)\mathrm{Pic}(X)\right)\times .$$
Here $`\mathrm{Pic}(X)`$ is the group of line bundles. Thus $`\mathrm{Out}_𝕂^{\mathrm{tr}}(𝖣^\mathrm{b}(\mathrm{𝖢𝗈𝗁}X))G\times D`$ where $`G`$ is an algebraic group and $`D`$ is a discrete group, and in particular this is a locally algebraic group.
###### Remark 8.
In \[Or\], Orlov gives an example of an abelian variety over $`𝕂`$ such that
$$\mathrm{Out}_𝕂^{\mathrm{tr}}(𝖣^\mathrm{b}(\mathrm{𝖢𝗈𝗁}X))D(X\times \widehat{X})(𝕂),$$
where $`D`$ is a discrete group (an extension of $`\mathrm{SL}_2()`$ by $``$) and $`\widehat{X}`$ is the dual abelian variety. The group $`D`$ acts (nontrivially) via $`\mathrm{Aut}_𝕂(X\times \widehat{X})`$ and hence $`\mathrm{Out}_𝕂^{\mathrm{tr}}(𝖣^\mathrm{b}(\mathrm{𝖢𝗈𝗁}X))`$ is a locally algebraic group.
Acknowledgments. I wish to thank Birge Huisgen-Zimmermann for showing me \[HS\], and for very helpful comments on an earlier version of this paper. Thanks also to Raphael Rouquier for calling my attention to \[Ro\] and for illuminating discussions.
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# Spin-accumulation and Andreev-reflection in a mesoscopic ferromagnetic wire
## I Introduction
Much theoretical and experimental work has addressed the effect of a superconductor (S) in proximity to a normal metal (N) on the transport properties during the last years, see Ref. and references therein for a review. Most experimental results can be explained in the framework of the quasiclassical theory of superconductivity accounting for a “long range” proximity effect with a coherence length $`\xi =(\mathrm{}D/2k_\text{B}T)^{1/2}`$, where $`D`$ is the diffusion coefficient of the normal metal and $`T`$ is the temperature. On the other hand, applications of the quasi-classical theory to transport in heterostructures containing ferromagnets (F) are still scarce. In contrast to normal metals the presence of a strong exchange field in the ferromagnet leads to a strong difference in the energy dispersions for the two spin bands. However, long-range coherence in normal metals requires spin degenerate bands close to the Fermi energy, since singlet superconductivity couples quasiparticles of different spins by Andreev reflection. The consequence of the exchange field energy $`h_{\text{xc}}`$ is a strong decoherence of quasiparticles belonging to the different spin bands. Typically the superconducting energy scale $`\mathrm{\Delta }`$ is smaller than $`h_{\text{xc}}`$ by several orders of magnitude for (Al, Nb) vs. (Fe, Ni, and Co), respectively. Thus, the proximity effect in ferromagnetic metals is negligible and a ferromagnet in contact to a superconductor may be considered as an incoherent metal coupled to the superconductor. In this case all changes induced by the contact to a superconductor depend on the properties of the interface itself. This is accomplished by the effect of spin accumulation, which requires no phase coherence in the ferromagnet and can therefore have a much longer range than the proximity effect. The main purpose of this paper is to study the mutual influence of resistance changes by spin accumulation and interface properties.
Recently heterostructures of ferromagnets and superconductors have been experimentally realized and investigated. Several unusual phenomena have been unveiled. The experimental results in point contact geometries can be explained by the reduced, bias-dependent transparency of the interface due to spin-dependent band mismatch between the normal metal and the ferromagnet. The experimental results in diffusive nanostructured samples are more intriguing. The measured conductance changes on the ferromagnetic side can be positive and negative at the superconducting transition with amplitudes much larger than anticipated. The sign and the amplitude of the changes appear to depend strongly on the ferromagnetic - superconductor interface transparency. It has been conjectured that a strong mutual influence of the superconductors and ferromagnetic conductors and a penetration of the superconducting order parameter into the ferromagnet over distances many times longer than expected from the above estimates might explain the observations.
Some effects of the interplay between spin accumulation and Andreev reflection in diffuse systems have been discussed in Refs. . Since the spin-current into a superconductor vanishes at sufficiently low bias and temperature, a non-equilibrium spin-accumulation builds up on the ferromagnetic side in order to conserve the spin-currents. The spin-accumulation causes an additional boundary resistance which is of the order of the resistance of the ferromagnetic wire of a length of the spin-flip diffusion length. Therefore the resistance of the F-S system should be always larger than that of the F-N system, in contradiction with many of the experimental observations. A possible reason for this apparent failure is the neglect of changes in the interface resistance in the transition from F-S to F-N. Previous theories took into account only perfectly ballistic interfaces for which the resistance is determined purely by the matching of the adjacent Fermi surfaces. The interface resistance and its modulation are of the order of the Sharvin resistance, which is negligible compared to the total one. However, in the sputtered samples with relatively large contact areas the interface can contribute significantly, especially when differences of resistances below and above the superconducting transition temperatures are considered. Other transport phenomena in ferromagnet-superconductors systems have been studied in Ref. .
It has been speculated that the triplet component of the order parameter induced by the fluctuations of the spin-orbit scattering potentials is essential in mesoscopic junctions. Neglecting magnetic impurities and spin-orbit coupling the superconducting order parameter is a spin singlet. However, magnetic impurities or spin-orbit coupling, induce a fluctuating spin triplet component with zero average. The triplet component is ‘long-range’ coherent in the ferromagnet since it couples electrons and holes with the same spin and the exchange field in the ferromagnet does not play a role. However, the contribution to the conductivity from the triplet fluctuations is only relevant when the fluctuations are relatively large which is only the case when the conductance is close to the quantum conductance. The experimental samples have a much larger conductance, and we do not expect that such mesoscopic fluctuations play an important role.
None of the above-mentioned theories can explain the recent experimental results. This makes it necessary to study the properties of the contact between the ferromagnet and the normal metal in more detail and to account for a possible spin accumulation and heating effects in the ferromagnet including all possible interfaces between the ferromagnets and the superconductors. In particular we will go beyond the assumptions of a perfect transparent metallic interface and discuss its influence on the observed conductance changes below the superconducting transition temperature and for bias voltages less than the superconducting gap. We will in this work radically disregard the proximity effect. Therefore our results only apply to ferromagnets with $`h_{\text{xc}}\mathrm{\Delta }`$, which is e.g. the case for the magnetic transition metals (Fe, Co, and Ni) in conjunction with superconducting metals like Nb and Al. This assumption is supported by the experimental fact that FS interferometers show no phase-periodic oscillations down to the level of $`0.1e^2/h`$ in strong ferromagnets. In contrast to the calculations presented in this paper, the proximity effect could be important in weak ferromagnets. Below we will show that most of the recent experimental results can be explained in terms of the energy-dependence introduced by interface conductance and the accomplishing change in the spin accumulation. It is important to note that these changes are small in comparison to the total resistance, which is dominated by the long ferromagnetic wire. Nevertheless they play a dominant role when only the resistance changes are measured.
The paper is organized in the following way: Section II gives a description of the diffusive ferromagnetic wire both in the limit of elastic and inelastic scattering between the electrons. Section III treats the boundary condition between the ferromagnet and the superconductor which is crucial to the understanding of the transport properties. The results for the conductance obtained from the description of the ferromagnetic wire with the boundary conditions are discussed in Section IV. Finally we compare our results with experiments in Section V and give our conclusions in Section VI.
## II Description of the ferromagnet
We consider a ferromagnetic diffusive wire connected to an ideal (ferromagnetic or normal metal) reservoir on one side and to a superconducting reservoir on the other side as depicted in Fig. 1. The wire is characterized by length $`L`$, cross-section $`A`$ and spin-dependent conductivities $`\sigma _{}`$ and $`\sigma _{}`$. In this Section we discuss the kinetic equations describing the ferromagnetic wire in the absence of the proximity effect. We consider collision with impurities to be the dominant scattering process and use the diffusion approximation. The electrons in the quasi-one-dimensional wire are described by energy $`ϵ`$ and spatial $`x`$ dependent distribution functions $`f_s(ϵ,x)`$ for the two spin directions $`s=+,`$ for spin $``$ and $``$, respectively. The distribution functions obey two coupled Boltzmann equations in the diffusive limit. Other scattering mechanisms will be specified in the following subsections.
Instead of the spin-dependent conductivities, it is convenient to introduce the total conductivity $`\sigma =\sigma _{}+\sigma _{}`$ and the spin-polarization of the conductivity $`\gamma =(\sigma _{}\sigma _{})/\sigma `$. The spin dependent conductivities are then expressed as $`\sigma _s=(1+s\gamma )\sigma /2`$. We will also make use of the total conductance (resistance) of the ferromagnetic wire $`G_\text{F}=A\sigma /L`$ ($`R_\text{F}=1/G_\text{F}`$).
### A Elastic scattering
In the elastic scattering case the energy is conserved in the scattering processes. This makes it necessary to study the energy dependent distribution functions in the ferromagnetic wire. In addition to elastic impurity scattering we consider only the spin-flip scattering processes, accounted for by the spin-flip length $`l_{\text{sf}}`$. Then, the kinetic equation reads
$$\frac{d^2}{dx^2}f_s(ϵ,x)=\frac{1}{l_{\mathrm{sf}}^2}\left[f_s(ϵ,x)f_s(ϵ,x)\right].$$
(1)
The current for spin $`s`$ is given by $`(e=|e|)`$
$$I_s(x)=\sigma _s\frac{A}{e}𝑑ϵ\frac{df_s(ϵ,x)}{dx}=𝑑ϵI_s(ϵ,x).$$
(2)
This equation defines the spectral current $`I_s(ϵ,x)`$. Electrical and spin currents are $`I_{\text{charge}}=I_{}(x)+I_{}(x)`$ and $`I_{\text{spin}}(x)=I_{}(x)I_{}(x)`$ and similar for the spectral currents. It is convenient to introduce the conductivity-averaged distribution function
$$f_{\text{el}}(ϵ,x)=\frac{\sigma _{}}{\sigma }f_{}(ϵ,x)+\frac{\sigma _{}}{\sigma }f_{}(ϵ,x)$$
(3)
and the nonequilibrium spin distribution function
$$f_{\text{sp}}(ϵ,x)=f_{}(ϵ,x)f_{}(ϵ,x).$$
(4)
The kinetic equations in terms of these functions read
$`{\displaystyle \frac{d^2}{dx^2}}f_{\text{el}}(ϵ,x)`$ $`=`$ $`0`$ (5)
$`{\displaystyle \frac{d^2}{dx^2}}f_{\text{sp}}(ϵ,x)`$ $`=`$ $`{\displaystyle \frac{1}{l_{\text{sf}}^2}}f_{\text{sp}}(ϵ,x).`$ (6)
The first equation is the spectral current conservation and the second describes spatial relaxation of the non-equilibrium spin-distribution.
### B Inelastic scattering
The reason to investigate the role of inelastic scattering is the convenient fact that the ferromagnet is an incoherent metal with rather strong correlations. Both phonon and electron-electron scattering can mediate inelastic scattering. In general it is not obvious which should dominate and both should be treated on equal footing. In order to achieve insight in the physics it is useful to consider limiting cases as well.
In the limit of strong inelastic scattering we assume that the electron-electron interaction is stronger than the electron-phonon relaxation. When a bias voltage is applied, the local electron temperature can therefore be different from the temperature in the reservoirs. This transport regime is relevant when the typical inelastic scattering length is smaller than the spin-diffusion length.
The electrons relax to a local equilibrium
$$f_s(ϵ,x)=f(ϵ;\mu _s(x),T_{\text{el}}(x)),$$
(7)
where $`\mu _s(x)`$ is the spin-dependent chemical potential, $`T_{\text{el}}(x)`$ is the local temperature and
$$f(ϵ;\mu ,T)=\frac{1}{1+\mathrm{exp}(\left(ϵ\mu \right)/k_\text{B}T)}$$
(8)
is the Fermi-Dirac distribution function. This makes it possible to integrate the kinetic equation and the currents over energy and to obtain equations for the local chemical potentials and electron temperature.
The spin-dependent (electric) current from Eq. (2) is
$`I_s(x)`$ $`=`$ $`{\displaystyle \frac{\sigma }{2e}}(1+s\gamma ){\displaystyle \frac{d\mu _s(x)}{dx}}.`$ (9)
Current conservation requires
$$\frac{d^2}{dx^2}\left(\sigma _{}\mu _{}(x)+\sigma _{}\mu _{}(x)\right)=0.$$
(10)
Spin relaxation occurs within the spin-diffusion length $`l_{\text{sf}}`$:
$$\frac{d^2}{dx^2}\left[\mu _{}(x)\mu _{}(x)\right]=\frac{1}{l_{\text{sf}}^2}\left[\mu _{}(x)\mu _{}(x)\right].$$
(11)
The local spin-dependent chemical potentials in the ferromagnet are determined by (10) and (11) and the boundary conditions to be discussed below.
Additionally, we need equations describing energy transport in the system to account for heating of the electrons. The energy current is
$`I_ϵ(x)`$ $`=`$ $`{\displaystyle \frac{A}{e^2}}{\displaystyle \underset{s}{}}\sigma _s{\displaystyle 𝑑ϵϵ\frac{df_s(ϵ,x)}{dx}}`$ (12)
$`=`$ $`\left[\mu _{}(x)I_{}(x)+\mu _{}(x)I_{}(x)\right]/e+I^Q(x),`$ (13)
where the heat current is
$$I_\text{Q}(x)=\kappa _\text{Q}(x)A\frac{dT_{\text{el}}(x)}{dx},$$
(14)
the heat conductivity $`\kappa _\text{Q}(x)=\sigma _0T_{\text{el}}(x)`$ and the Lorentz number is $`_0=\frac{\pi ^2}{3}\left(\frac{k_\text{B}}{e}\right)^2`$ .
The conservation law for the energy current dictates
$$\frac{d}{dx}I_ϵ(x)=A\left(\frac{\rho _ϵ(x)}{t}\right)_{\text{rel.}},$$
(15)
where $`\rho _ϵ(x)`$ is the local energy density. The energy relaxation between the electronic system and the phonons at sufficiently low temperatures is
$$\left(\frac{\rho _ϵ(x)}{t}\right)_{\text{rel.}}=\zeta \left[\left(k_\text{B}T\right)^5\left(k_\text{B}T_{\text{el}}(x)\right)^5\right],$$
(16)
where $`\zeta `$ parameterizes the strength of the electron-phonon interaction $`\zeta =48\pi \zeta (5)N(ϵ_F)\lambda ^{}/(\mathrm{}^3\omega _D^2)`$, $`\zeta (5)1.04`$ is the Riemann zeta function, $`N(ϵ_F)`$ is the density of states of both spins per unit cell, $`\lambda ^{}`$ is of the same order of magnitude as the electron-phonon coupling constant $`\lambda `$, and $`\mathrm{}\omega _D`$ is the Debye energy.
The conservation of energy current (15) together with the expression for the energy relaxation (16) give a differential equation for the local electron temperature which can be solved together with the boundary conditions to be discussed below.
When the electron-phonon interaction is weak there is no exchange of energy between the electron and the phonon systems so that the right hand side of Eq. (15) can be set to zero and we have conservation of the energy current due to the electron transport $`dI^e(x)/dx=0`$. In the opposite limit of strong electron-phonon interaction the electron temperature equals the lattice temperature. The differential equation for the energy conservation with the boundary conditions given above can in these two cases be solved exactly. In the intermediate regime the equations will be solved numerically.
## III Boundary conditions
The condition that the ferromagnet should be completely incoherent leads to simplified boundary conditions for the kinetic equations. These boundary conditions can be derived from the boundary conditions for the quasiclassical Green’s function. A transparent form suitable for diffusive systems has been presented by Nazarov. We will follow the spirit and the notation of this paper. A circuit theory for ferromagnetic-normal metal systems has been presented in Ref. . A contact is described by a set of transmission eigenvalues $`\{T_n\}`$ or, equivalently, by a distribution of the transmission eigenvalues $`\rho (T)`$. The boundary condition at the contact is expressed through a conservation law for the matrix current in the Keldysh formulation. In the framework of superconductivity it is a $`4\times 4`$-matrix comprising $`2\times 2`$ Keldysh space and $`2\times 2`$ particle-hole (Nambu) space. The matrix current through the FS-contact is
$$\stackrel{ˇ}{I}=\frac{2e}{\pi \mathrm{}}\underset{n}{}T_n\frac{\left(\stackrel{ˇ}{G}_\text{F}\stackrel{ˇ}{G}_\text{S}\stackrel{ˇ}{G}_\text{S}\stackrel{ˇ}{G}_\text{F}\right)}{4+T_n(\stackrel{ˇ}{G}_\text{F}\stackrel{ˇ}{G}_\text{S}+\stackrel{ˇ}{G}_\text{S}\stackrel{ˇ}{G}_\text{F}2)}.$$
(17)
This matrix current has to be equated to the diffusive matrix current entering the contact from either side. The two sides of the contact are characterized by the Keldysh matrix Green’s functions $`\stackrel{ˇ}{G}_\text{S}`$ and $`\stackrel{ˇ}{G}_\text{F}`$, which we will specify to be the superconducting reservoir and the ferromagnetic wire, respectively. The Keldysh-Nambu matrix Green‘s function of the superconductor in equilibrium is
$$\stackrel{ˇ}{G}_\text{S}(ϵ)=\left(\begin{array}{cc}\widehat{G}_\text{S}^\text{R}(ϵ)& \widehat{G}_\text{S}^\text{K}(ϵ)\\ 0& \widehat{G}_\text{S}^\text{A}(ϵ)\end{array}\right).$$
(18)
A similar structure holds for any matrix in Keldysh space. In local equilibrium the Keldysh (1,2) component in Nambu space is
$$\widehat{G}_\text{S}^\text{K}(ϵ)=(\widehat{G}_\text{S}^\text{R}(ϵ)\widehat{G}_\text{S}^\text{A}(ϵ))\left(12f^\text{S}(ϵ)\right),$$
(19)
where $`f^\text{S}(ϵ)=[1+\mathrm{exp}(ϵ/k_\text{B}T)]^1`$ is the quasi-particle distribution function in the superconductor and we have set the chemical potential in the superconductor to zero. $`\widehat{G}_\text{S}^\text{R}(ϵ)`$ and $`\widehat{G}_\text{S}^\text{A}(ϵ)`$ are retarded and advanced Nambu Green’s functions determining the spectral properties of the superconductor. In the BCS case with a real order parameter they are
$`\widehat{G}^\text{R}(ϵ)=\left(\widehat{G}^\text{A}(ϵ)\right)^{}={\displaystyle \frac{(ϵ+i0)\widehat{\tau }_3i\mathrm{\Delta }\widehat{\tau }_1}{\sqrt{(ϵ+i0)^2\mathrm{\Delta }^2}}}.`$ (20)
The diagonal component represents the normal retarded Green’s function whereas the off-diagonal component is conventionally called the anomalous Green’s function. On the ferromagnetic side we completely neglect the proximity effect leading to the spectral functions $`\widehat{G}_\text{F}^\text{R}=\widehat{\tau }_3=\widehat{G}_\text{F}^\text{A}`$. The absence of an anomalous component is a result of the absence of the proximity effect. The Keldysh component accounts for the spin-dependent non-equilibrium distribution:
$$\widehat{G}_\text{F}^\text{K}(ϵ)=2\left(\begin{array}{cc}12f_{}^\text{F}(ϵ)& 0\\ 0& 12f_{}^\text{F}(ϵ)\end{array}\right),$$
(21)
where $`f_{}^\text{F}(ϵ)`$ and $`f_{}^\text{F}(ϵ)`$ are the quasi-particle distribution functions close to the interface on the ferromagnetic side. The spectral electrical current is determined by the Keldysh-component of the matrix current according to
$$I_{\text{el}}(ϵ)=\frac{1}{4e}\text{Tr}\left[\widehat{\tau }_3\widehat{I}^\text{K}(ϵ)\right].$$
(22)
Eq. (21) and Eq. (22) suggest a representation of the diagonal components of the Keldysh component of the matrix current in the form
$$\widehat{I}^\text{K}(ϵ)=\left(\begin{array}{cc}I_{}(ϵ)& \mathrm{}\\ \mathrm{}& I_{}(ϵ)\end{array}\right).$$
(23)
Now we are in the position to calculate the spin-resolved currents through the contact. Performing the calculations along the lines of Ref. we find the spectral spin-dependent current
$`I_s(ϵ)`$ $`=`$ $`{\displaystyle \frac{G_{\text{QP}}(ϵ)}{2e}}\left(f^S(ϵ)f_s^F(ϵ)\right)`$ (25)
$`+{\displaystyle \frac{G_\text{A}(ϵ)}{4e}}\left(1f_s^F(ϵ)f_s^F(ϵ)\right).`$
The quasiparticle conductance $`G_{\text{QP}}(ϵ)`$ and the Andreev conductance $`G_A(ϵ)`$ are determined by the properties of the contact and the spectral properties of the two metals connected by the contact. The distribution of transmission eigenvalues can be incorporated in a single characteristic complex function
$`Z(x)={\displaystyle \frac{e^2}{\pi \mathrm{}}}{\displaystyle \underset{n}{}}{\displaystyle \frac{T_n}{2+T_n(x1)}},`$ (26)
where $`x(ϵ)=\text{Tr}\{\widehat{G}_\text{S}^\text{R}(ϵ),\widehat{G}_\text{F}^\text{R}(ϵ))\}/4`$. The conductances are
$`G_{\text{QP}}(ϵ)`$ $`=`$ $`\text{Re}Z(x)\text{Re}x+{\displaystyle \frac{\text{Im}Z(x)}{\text{Im}x}}\text{Im}^2\sqrt{1x^2},`$ (27)
$`G_\text{A}(ϵ)`$ $`=`$ $`{\displaystyle \frac{\text{Im}Z(x)}{\text{Im}x}}\left|1x^2\right|.`$ (28)
The contact is characterized by a transmission distribution, which leads to contact-specific energy dependences of the conductances. The normal state conductance is $`G_{\text{BN}}=(e^2/2\pi \mathrm{})_nT_n`$. For a ballistic model contact all transmission eigenvalues are equal to one for the propagating channels and zero otherwise and $`_nT_n=N`$, where $`N`$ is the number of propagating channels. The distribution function in the case of a dirty interface is
$$\rho (T)=\frac{\mathrm{}}{e^2}G_{\text{BN}}\frac{1}{T^{3/2}\sqrt{1T}}$$
(29)
and in the case of a diffusive contact the distribution is
$$\rho (T)=\frac{\mathrm{}}{2e^2}G_{\text{BN}}\frac{1}{T\sqrt{1T}}.$$
(30)
Finally, for a tunnel conductance a perturbation expansion in terms of the small transmission eigenvalues can be performed. We list the characteristic function $`Z(x)`$ for a number of generic contacts in Table (I): Tunnel junction, ballistic contact, diffusive contact and dirty interface. In the case of an incoherent metal on one side (i.e. $`\widehat{G}_\text{F}^\text{R}=\widehat{\tau }_3`$), the argument of the characteristic function reduces to $`x=\text{Tr}\widehat{\tau }_3\widehat{G}_\text{S}^\text{R}/2`$. The result in this case is demonstrated explicitely in Table I for a contact of a BCS-superconductor with spectral functions (20). The energy dependence of these spectral conductances is depicted in Fig 2. Below the superconducting gap only the Andreev conductance is nonzero, gradually decreasing from the value of $`2G_{\text{BN}}`$ for the metallic junction to zero in the tunnel junction. Above the gap the Andreev conductance vanishes rather quickly $`1/ϵ^2`$. Also quasiparticle transport becomes possible and, thus, spin-transport into the superconductor.
The properties of these contacts are demonstrated by the temperature dependence of the linear conductance following from
$$G_{\text{BS}}(T)=𝑑ϵ(G_{\text{QP}}(ϵ)+G_\text{A}(ϵ))\left(\frac{f(ϵ,0,T)}{ϵ}\right).$$
(31)
This is the conductance that would be measured if the contact would be placed between a normal reservoir and a superconducting reservoir. The temperature dependence of the contact conductance (31) is shown in Fig. 3. The dashed and dotted lines show the conductance of the diffusive contact and the dirty interface, respectively. The resistance of the diffusive contact shows the well known reentrant behavior, i.e. it reaches the normal state conductance at zero temperature. The resistance of the dirty interface after a small drop below the critical temperature is higher than the normal state value and saturates at low temperature at $`\sqrt{2}R_{\text{BN}}`$.
Additionally, we introduce mixed contacts as a model for an inhomogeneous interface with distributed regions with low and high transparency. The relative admixture $`q`$ of a tunnel and $`(1q)`$ of a ballistic contact allows switching continuously from one limit to the other, covering approximately the universal cases of a diffusive contact ($`q0.5`$) and a dirty interface ($`q1/\sqrt{2}`$ with a single parameter $`q`$. A common feature of the temperature dependence of all these contacts except the tunnel junction is that right below $`T_\text{c}`$ the resistance drops. At lower temperatures the resistance increases again except in the case of the purely ballistic contact. The drop of resistance of these contacts close to $`T_\text{c}`$ can be traced back to the temperature dependence of the superconducting order parameter $`\mathrm{\Delta }(T)`$. The resistance drop is caused by the leading order contribution of the change in the superconducting gap $`\mathrm{\Delta }(T)(1T/T_\text{c})^{1/2}`$ to the Andreev contribution and the conductance.
The boundary condition presented so far imply that the transmission ensembles and the number of channels are the same for the two spin species. In reality the transmission matrices for spin-up and spin-down states can be different. A microscopic calculation of the transmission eigenvalues is beyond the scope of the present paper. We will therefore heuristically generalize the boundary conditions to spin-dependent interfaces by taking different transmission ensembles for the two spin directions. These ensembles can differ in the total number of channels and/or in the transmission distribution. Thus, we replace the spin-dependent current through the interface (25) by
$`I_\text{s}(ϵ)`$ $`=`$ $`{\displaystyle \frac{G_\text{s}(ϵ)}{2e}}(f^S(ϵ)f_\text{s}^\text{F}(ϵ,0))`$ (33)
$`{\displaystyle \frac{G_\text{A}(ϵ)}{4e}}\left(1f_\text{s}^\text{F}(ϵ,0)+f_{\text{-s}}^\text{F}(ϵ,0)\right),`$
In general the spin-dependent quasi-particle conductances $`G_{}(ϵ)`$ and $`G_{}(ϵ)`$ entering the first term are of different magnitude and have different energy dependences. Similar as for the ferromagnetic wire we introduce the total conductance of the boundary $`G_B(ϵ)=G_{}(ϵ)+G_{}(ϵ)`$ and a dimensionless factor $`\gamma _B(epsilon)=(G_{}(ϵ)G_{}(ϵ))/G_B(ϵ)`$, which we call polarization of the boundary conductance. Since the definitions (27) of quasi-particle and (28) of Andreev conductance have been derived for a spin-degenerate interface these definition are not valid anymore for spin-dependent interface scattering. It is, however, reasonable to assume that the energy dependence of all conductances is well approximated by the same transmission ensemble, but different numbers of channels. We can motivate this choice by the fact that in the experiments that we have in mind the interfaces are strongly disordered regions, with a possible formation of an alloy layer extending over several monolayers. In such contacts the number of channels is more or less controlled by the differences of the cross sections of the Fermi surface. But, on the other hand, the transmission ensemble and, hence, the energy dependence of the conductance is not expected to vary much in typical disordered contacts on the scale of the superconductor gap.
We will in the following only take into account the differences in magnitude, but not in energy dependence. In the language of transmission distributions this means that the distributions are the same, but the number of channels differ. In this approximation the spin-polarization of the boundary conductance $`\gamma _\text{B}`$ is energy independent. The energy dependence of the Andreev conductance follows from the same transmission ensemble, but its magnitude will be reduced in comparison to the unpolarized case. It is important to notice that the boundary polarization and the polarization of the ferromagnetic wire need not to have the same sign, since they are parametrically independent. The possibility of this is demonstrated by microscopic numerical calculations.
## IV Results and discussions
We solve the kinetic equations presented in Sec. II in the three cases:
* purely elastic scattering
* inelastic scattering in linear response
* inelastic scattering in nonlinear response
The boundary condition on the superconducting side of the wire has been derived in Sec. III. The boundary conditions at the ferromagnetic reservoir are
$$f_{}(ϵ,L)=f_{}(ϵ,L)=f(ϵ;eV,T),$$
(34)
where $`f(ϵ;eV,T)=(\mathrm{exp}((ϵeV)/k_\text{B}T)+1)^1`$ is the Fermi-Dirac equilibrium distribution at a constant voltage $`V`$ and temperature $`T`$.
In the case of inelastic scattering Eq. (34) also implies that the electron temperature equals the lattice temperature in the ferromagnetic reservoir $`T_{\text{el}}(x=L)=T`$. The other boundary condition for the electron temperature comes from the conservation of energy current in the ferromagnet and into the superconductor.
As a reference we calculate the resistance of the system in the normal state
$$R_{\text{FN}}=R_\text{F}+R_{\text{BN}}+R_{\text{sf}}\frac{(\gamma \gamma _\text{B})^2}{1+R_{\text{sf}}/R_{\text{BN}}},$$
(35)
The third term is due to the spin-accumulation in the ferromagnetic wire determined by the ‘spin-flip resistance’ $`R_{\text{sf}}=1/G_{\text{sf}}=R_\text{F}l_{\text{sf}}/L\mathrm{tanh}(l_{\text{sf}}/L)(1\gamma ^2)`$. In the limit of a weak ferromagnet $`\gamma ^21`$ and a short spin-flip relaxation length $`l_{\text{sf}}L`$ the spin-flip resistance reduces to $`R_{\text{sf}}R_\text{F}l_{\text{sf}}/L`$, i.e. the resistance of a piece of the ferromagnetic wire of length $`l_{\text{sf}}`$. We see that the excess resistance due to the spin-accumulation increases with increasing asymmetry between the polarization of the bulk conductivity and the polarization of the interface conductance. The expression (35) will be used in the following to calculate resistance changes below the transition to the superconducting state:
$$\mathrm{\Delta }R_{\text{FS}}(T,V)=R_{\text{FS}}(T,V)R_{\text{FN}}.$$
(36)
The differential resistance is defined by
$$R_{\text{FS}}(T,V)=\left(\frac{I(T,V)}{V}\right)^1.$$
(37)
In the linear response regime we will omit the arguments of the differential resistance $`R_{\text{FS}}R_{\text{FS}}(T,V0)`$.
In the following analysis it will be useful to define the following temperature dependent average
$`\mathrm{}={\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{}\left({\displaystyle \frac{f(ϵ;0,T)}{ϵ}}\right)𝑑ϵ`$ (38)
This average occures, e.g., in the temperature dependent conductances of a contact between an incoherent metal and a superconductor (31).
### A Elastic scattering
When the scattering in the wire is elastic, the general solution of (5) satisfying (34) may be written as
$`f_{\text{el}}(ϵ,x)`$ $`=`$ $`{\displaystyle \frac{G_{\text{FS}}(ϵ)}{G_\text{F}}}\left(f(ϵ;0,T)f(ϵ;eV,T)\right)\left(1+{\displaystyle \frac{x}{L}}\right)`$ (40)
$`+f(ϵ;eV,T).`$
The spatially independent spectral conductance $`G_{\text{FS}}(ϵ)`$ determines the current through the structures and remains to be found. The solution of the second kinetic equation (6) satisfying the boundary condition (34) is
$$f_{\text{sp}}(ϵ,x)=2\alpha \mathrm{sinh}\left(\frac{L+x}{l_s}\right),$$
(41)
where the parameter $`\alpha `$ should be found from the continuity of the spin currents into the superconductor (33) and ferromagnetic wire (2). We find the electrical current
$`I(T,V)`$ $`=`$ $`{\displaystyle \frac{1}{2e}}{\displaystyle }dϵG_{\text{FS}}(ϵ)\times `$ (43)
$`\left[1f(ϵ;eV,T)f(ϵ;eV,T)\right].`$
This expression shows that the spectral conductance determines the transport in each energy slice depending on the difference in occupation of states at this energy in the reservoirs. This form is analogous to the classical definition of a conductance as the proportionality factor between current and voltage difference.
The spectral conductance is given by
$`{\displaystyle \frac{1}{G_{\text{FS}}(ϵ)}}`$ $`=`$ $`{\displaystyle \frac{1}{G_\text{F}}}+{\displaystyle \frac{1}{G_{\text{QP}}(ϵ)+G_\text{A}(ϵ)}}`$ (45)
$`+{\displaystyle \frac{\left(\gamma \gamma _\text{B}\frac{G_{\text{QP}}(ϵ)}{G_{\text{QP}}(ϵ)+G_\text{A}(ϵ)}\right)^2}{G_{\text{sf}}+G_{\text{QP}}(ϵ)\left(1\gamma _\text{B}^2\frac{G_{\text{QP}}(ϵ)}{G_{\text{QP}}(ϵ)+G_\text{A}(ϵ)}\right)}}.`$
In the general case the full expression has to be used to calculate the resistance change in the superconducting state.
When the ferromagnetic wire dominates the resistance of the whole structure a simplified expression for the linear resistance change may be obtained. We first limit the discussion to the case of a weak ferromagnet and vanishing boundary polarization to obtain
$`\mathrm{\Delta }R_{\text{FS}}`$ $`=`$ $`{\displaystyle \frac{1}{G_{\text{QP}}(ϵ)+G_\text{A}(ϵ)}}R_{\text{BN}}`$ (47)
$`+\gamma ^2\left[{\displaystyle \frac{1}{G_{\text{sf}}+G_{\text{QP}}(ϵ)}}{\displaystyle \frac{1}{G_{\text{sf}}+G_{\text{BN}}}}\right].`$
We see that the resistance change consists of two contributions. The first is the resistance change due to the change of the boundary resistance, which would also be present in the absence of spin polarization. Note, however, that this term can be qualitatively different from the case of a normal metal wire in contact to a superconductor since in this case the proximity effect would not be negligible. The second term accounts for the difference in spin accumulation between normal and superconducting state.
First we discuss the influence of spin accumulation on the FS-resistance for a spin-degenerate interface. In Fig. 4 resistance changes for two types of contacts are shown for different polarizations of the ferromagnet. Solid curves are for a relatively good contact ($`q=0.75`$) and dashed curves for a less transparent contact ($`q=0.25`$). In this plot the total resistance of the system is dominated by the resistance of the ferromagnetic wire $`R_\text{F}=100R_{\text{BN}}`$ and the spin-relaxation length is $`l_{\text{sf}}=0.03L`$, resulting in a spin accumulation resistance $`R_{\text{sf}}3R_{\text{BN}}`$. Accordingly, the resistance change is normalized to $`R_{\text{BN}}`$ to show the relevant scale of the effect produced by the superconducting transition. For both contacts spin accumulation (increasing from the bottom to the top curves) leads to an enhancement of the resistance. Specifically the low temperature resistance is well accounted for by Eq. (47) in the limit $`T0`$:
$$\mathrm{\Delta }R_{\text{FS}}(T=0)=\frac{1}{G_\text{A}(0)}R_{\text{BN}}+\gamma ^2\frac{R_{\text{sf}}}{1+R_{\text{BN}}/R_{\text{sf}}}.$$
(48)
The second term of this equation shows that the spin-accumulation always enhances the resistance, maximally by an amount $`\gamma ^2R_{\text{sf}}`$. The enhancement for the $`q=0.25`$-contact has a uniform temperature dependence and does not change qualitatively. This is different for the $`q=0.75`$ contact. Here the resistance decreases monotonically in the unpolarized case as a result of the Andreev enhanced conductance. A small polarization $`\gamma 0.20.4`$ results in a nonmonotonic temperature dependence, i.e. an increase of resistance at lower temperatures. This can lead (for specific parameters) to a reentrant behavior of the resistance change, even overshooting the normal state value for larger spin-accumulation. At even higher spin polarizations $`\gamma ^2R_{\text{BN}}/R_{\text{sf}}`$ the Andreev contribution is completely masked and the resistance increases monotonically. This behavior resembles that of a less transparent contact if the absolute scale is properly chosen.
Let us now discuss the effect of the interface polarization $`\gamma _\text{B}`$ on the resistance change. In Fig. 5 the temperature dependent resistance of a $`q=0.75`$ contact is shown for different interface polarizations. Other parameters are $`R_\text{F}=100R_{\text{BN}}`$, $`l_{\text{sf}}=0.03L`$, and $`\gamma =0.3`$. The interface polarization $`\gamma _\text{B}`$ changes from the symmetric value $`+0.5`$ to the antisymmetric value $`0.5`$, as indicated in the plot. The reduction of the Andreev conductance by the spin-dependent interface resistance is taken into account by a phenomenological renormalization factor $`(1\gamma _\text{B}^2)`$. To gain some insight it is useful to look at the low temperature limit of the resistance change in the limit $`R_{\text{BN}}R_{\text{sf}}`$. From Eq. (45) it follows that
$$\mathrm{\Delta }R_{\text{FS}}(T=0)=\frac{1}{G_\text{A}(0)}R_{\text{BN}}+R_{\text{sf}}(4\gamma \gamma _\text{B}\gamma _\text{B}^2).$$
(49)
The spin-dependent contribution depends on the relative sign of the two polarizations $`\gamma `$ and $`\gamma _\text{B}`$ and can also be negative (if $`4\gamma \gamma _\text{B}<\gamma _\text{B}^2`$). This effect is seen from the lower two curves in Fig. 5 with an antisymmetric interface polarization. An increasing interface polarization leads to a lowering of the resistance change, despite the increase of the resistance due to the renormalization of the Andreev conductance. It is worthwhile noting that for the largest negative interface polarization shown ($`\gamma _\text{B}=0.5`$) the total resistance drop is larger than the resistance drop which would result from the pure Andreev reflection in the absence of spin polarization of the interface and the F-wire. This apparent contradiction to the intuition that any spin-accumulation should decrease the Andreev-caused resistance drop stems from the fact that we plot the resistance change below the superconducting transition. The contradiction is resolved by noting that the total resistance $`R_{\text{FS}}(T)=R_{\text{FN}}+\mathrm{\Delta }R_{\text{FS}}(T)`$ is always higher than for the unpolarized case. However, in a real experiment (with fixed polarizations) the Andreev conductance in the absence of a polarization can not be measured separately. It may therefore appear that the measured resistance drop is larger than one would expect from a simple estimate of the reduction of the interface resistance due to Andreev reflection.
### B Inelastic scattering - linear response
We will now proceed to study the case of inelastic scattering in the ferromagnetic wire. It is assumed that the current in the ferromagnet is weakly polarized, $`\gamma 1`$. In order to simplify the discussions we disregard the possible asymmetry in the interface transparency in the following discussions and set $`\gamma _\text{B}=0`$. An extension is straightforward.
An analytical expression for the total conductance of the system can be found in the linear response regime. In this regime the effects of electron heating vanish since they will only contribute to the current in higher orders of the source-drain bias. The coupled equations for the spin-dependent chemical potential distributions and the electron temperature are simplified by letting $`T_{\text{el}}(x)T`$. By solving (10) and (11) together with the boundary condition (25) and (34) we find the linear response resistance. Assuming a weak ferromagnet, $`\gamma ^21`$, and a small interface resistance compared to the resistance of the ferromagnetic wire $`R_{\text{BN}}R_\text{F}`$ the resistance change can be written as
$`\mathrm{\Delta }R_{\text{FS}}(T)`$ $`=`$ $`{\displaystyle \frac{1}{G_{\text{QP}}(ϵ)+G_\text{A}(ϵ)}}R_{\text{BN}}`$ (51)
$`+\gamma ^2\left[{\displaystyle \frac{1}{G_{\text{sf}}+G_{\text{QP}}(ϵ)}}{\displaystyle \frac{1}{G_{\text{sf}}+G_{\text{BN}}}}\right].`$
The first two terms in (51) are due to the effective interface resistance between the ferromagnet and the superconductor and the third term is due to the spin-accumulation. The latter term vanishes when $`\gamma 0`$ or $`l_{\text{sf}}/L0`$. This equation has to compared with Eq. (36) for the case of purely elastic scattering. Only the quasi-particle conductance enters the spin-accumulation contribution since spins cannot be injected into the superconductor by means of the Andreev process. The temperature averaged conductances directly determine the temperature dependence of the total resistance in the case of dominant inelastic scattering processes. The quasi particle conductance vanishes at zero temperature since then no spin-current can propagate into the superconductor. At zero temperature $`G_{\text{QP}}(ϵ)+G_\text{A}(ϵ)=G_\text{A}(0)`$ and the resistance the FS system in the case of inelastic scattering (51) equals the result in the case of elastic scattering (48).
The results with inelastic scattering in general differ from those with purely elastic scattering when the temperature is non-zero or when the current is measured in the non-linear source-drain response regime. The remarkable difference between Eq. (36) and Eq. (51) is the way the thermal averaging is carried out. E.g. in the first term we have to average the inverse contact conductance in the case of elastic scattering, whereas we first have to average the conductance and than invert the result in the case of inelastic scattering. A similar consideration holds for the spin-accumulation term. The origin of this difference can be understood in the following way: we may visualize our wire (or any system) as mapped onto an electric circuit which contains energy-dependent conductors. In the case of purely elastic scattering we first have to calculate the total conductance of the system for each energy. The current is then found by averaging this spectral conductance with the difference of distribution functions of the adjacent reservoirs. This procedure yields Eq. (36) for the change in the resistance. In contrast, inelastic scattering equilibrates the local distribution of electrons in a way that the chemical potential is equal to the potential found from solving the circuit problem of the corresponding electric circuit. Thus, Eq. (51) follows the Kirchhoff’s laws for our system. As we will demonstrate below this difference can have significant consequences for the temperature dependence of the resistance.
Let us now illustrate the temperature dependence of the linear response conductance in the case of a metallic contact with an interface conductance much larger than the conductance of the ferromagnetic wire. The total conductance at sufficiently low temperatures is then $`G_{\text{QP}}(ϵ)+G_\text{A}(ϵ)G_\text{A}(0)=2G_{\text{BN}}`$ and the quasiparticle conductance is $`G_{\text{QP}}(ϵ)(8\pi k_\text{B}T/\mathrm{\Delta })^{1/2}\mathrm{exp}(\mathrm{\Delta }/k_\text{B}T)`$. The temperature must then be so low that
$$k_\text{B}T\frac{\mathrm{\Delta }}{\mathrm{ln}\left(R_{\text{BN}}/R_{\text{sf}}\right)}$$
(52)
in order to prevent thermally assisted spin-current into the superconductor.
We show in Fig. 6 the ratio of the linear response resistance change $`\mathrm{\Delta }R_{\text{FS}}`$ to the interface resistance $`R_{\text{BN}}`$ as a function of the temperature $`T`$ for a metallic interface with $`R_{\text{BN}}=0.05R_\text{F}`$, polarization $`\gamma =0.3`$ and spin-flip diffusion length $`l_{\text{sf}}/L=0.2`$. For these parameters we have a ‘spin-flip’ resistance corresponding to $`R_{\text{sf}}=0.2R_\text{F}`$. The change in resistance below the superconducting transition temperature is due to a competition between the excess resistance caused by the spin-flip relaxation and the reduced interface resistance caused by Andreev reflection. At $`T=0`$ we find from the approximate result (51) that $`R_{\text{FS}}R_{\text{FN}}=0.5R_{\text{BN}}\gamma ^2R_{\text{sf}}=0.14R_{\text{FN}}`$ rougly corresponding to the numerical value which has been obtained without making the approximation $`\gamma ^21`$ and $`R_\text{F}R_{\text{BN}}`$. Using the condition (52) we find that the spin-accumulation is strongly reduced around $`T/T_\text{c}=0.7`$ and consequently the resistance of the system decreases before increasing again around $`T/T_\text{c}=1`$ where the boundary resistance is increased. This explains the non-monotonic behavior of the linear response resistance as a function of the temperature.
### C Inelastic scattering - nonlinear response
At a finite bias voltage the electron heating effects have to be taken into account and the coupled equations for the electron temperature, the spin-dependent chemical potentials (10), (11), (13), and (15) and the boundary conditions (25) and (34) have to be solved numerically.
First let us discuss the transport properties when the electron-phonon interaction is weak so that we have perfect conservation of energy current and the left hand side of (15) can be set to zero. From the discussions in the previous section we understand that there will be a reduction in the excess resistance due to the spin-accumulation when the electron temperature on the ferromagnetic side reaches condition (52) so that there is a significant spin-current entering the superconductor. Roughly speaking, the electron temperature on the ferromagnetic side is proportional to the applied source drain bias. Thus, as a crude approximation, we expect that the excess resistance due to the spin-accumulation is lowered when
$$eV\frac{\mathrm{\Delta }}{\mathrm{ln}\left(R_{\text{BN}}/R_{\text{sf}}\right)}.$$
(53)
We show in Fig. 7 the resistance change $`R_{\text{FS}}(T=0,V)R_{\text{BN}}`$ (36) normalized by the interface resistance $`R_{\text{BN}}`$ as a function of the bias voltage $`V`$. As before, the interface resistance is $`R_{\text{BN}}=0.05R_\text{F}`$, the polarization $`\gamma =0.3`$ and the spin-flip diffusion length $`l_{\text{sf}}/L=0.2`$. For these parameters we have a ‘spin-flip’ resistance corresponding to $`R_{\text{sf}}=0.2R_\text{F}`$. The change in resistance below the superconducting gap is due to a competition between the excess resistance caused by the spin-flip relaxation and the reduced interface caused by the Andreev reflection. A dip in the resistance is seen around $`V=0.7\mathrm{\Delta }`$ which is correctly described by (53). Below this bias voltage the resistance is caused by the competetion between spin-accumulation which enhances the resistance and the effective interface resistance which reduces the resistance. At higher voltages the resistance is only caused by the effective interface resistance and the reduction of the resistance change $`R_{\text{FS}}(T=0,V)R_{\text{BN}}`$ as a function of the bias voltage is small.
In the limit of strong electron-phonon interaction the electron temperature equals the lattice temperature. The spin-current into the superconductor is then not enhanced due to thermal activation and consequently the spin-accumulation on the ferromagnetic side is only reduced when the potential on the ferromagnetic side of the interface is higher than the superconducting gap. This occurs when
$$eV=\mathrm{\Delta }\left(1+2\frac{R_\text{F}}{R_{\text{BN}}}\right),$$
(54)
and thus at a potential that is much larger than the superconducting gap, in contrast to the case of weak electron-phonon interaction.
In the intermediate regime the electron-phonon interaction should be included. In order to illustrate the main physics we consider the case of a weak polarization and set $`\gamma =0`$ and consequently there are no effects due to spin-accumulation and the resistance change of the wire is only due to the change of the effective boundary resistance. The chemical potential in the ferromagnetic wire is thus spin-independent. Furthermore we consider the case that the lattice temperature is zero, $`T=0`$ so that the electron temperature arises solely due to electron heating. Solving the diffusion equation (10) on the ferromagnetic side of the interface gives $`\mu (x)=eVx/L+\mu (0)\left[1+x/L\right]`$, where $`eV`$ is the applied bias and $`\mu (0)`$ is the potential drop across the ferromagnet-superconductor interface. The superconducting energy gap $`\mathrm{\Delta }`$ presents a natural energy scale for the problem. We will characterize the strength of electron-photon energy exchange by a dimensionless constant $`\kappa =AL\zeta e^2\mathrm{\Delta }^3/G_\text{F}`$, $`\zeta `$ is defined by the relation (16). The energy diffusion equation then simplifies to
$`{\displaystyle \frac{\pi ^2}{6}}\left(L_T_x\right)^2(k_\text{B}T_{\text{el}}/\mathrm{\Delta })^2`$ $`=`$ $`(k_\text{B}T_{\text{el}}/\mathrm{\Delta })^5`$ (56)
$`\left((eV\mu (0))/\mathrm{\Delta }\right)^2/\kappa ,`$
where we introduce a typical length-scale for the energy exchange $`L_T=L/\sqrt{\kappa }`$. If $`\kappa 1`$ $`LL_T`$ and the exchange is not effective. For longer wires, $`\kappa `$ becomes bigger than unity. In this case, the electron temperature develops a constant plateau in the ferromagnetic wire and only changes rapidly within the length-scale $`L_T`$ near the end-points $`x=L`$ and $`x=0`$. It follows from (56) that in this case the temperature in the middle of the ferromagnetic wire becomes
$$\frac{k_\text{B}T_{\text{el}}}{\mathrm{\Delta }}=\kappa ^{1/5}\left(\frac{eV\mu (0)}{\mathrm{\Delta }}\right)^{2/5}.$$
(57)
We will now present numerical results of the temperature profile in the ferromagnetic wire and the resulting resistance change using $`\gamma =0`$, a metallic interface $`R_{\text{BN}}/R_\text{F}=0.05`$ for various values of the electron-phonon coupling interaction. We show in Fig. 8 the spatially dependent electron temperature in the ferromagnetic wire for $`\kappa =10^6`$ at a bias voltage $`eV=40\mathrm{\Delta }`$ (upper curve), $`eV=20\mathrm{\Delta }`$ (mid curve) and $`eV=10\mathrm{\Delta }`$ (lower curve). The electron temperature in the middle of the wire follows from (57). There are rapid changes of the electron temperature close to the ferromagnetic and superconducting reservoirs and the temperature in the middle of the wire is lower than the electron temperature close to the superconductor. The latter temperature is important for the effective interface resistance.
We show in Fig. 9 the resistance change as a function of the bias voltage. The different solid lines show the current for different ratios of the electron-phonon coupling starting from no electron-phonon interaction (a) $`\kappa =0`$ going through intermediate electron-phonon interaction (b) $`\kappa =10^2`$, (c) $`\kappa =10^6`$, (d) $`\kappa =10^8`$ to strong electron-phonon interaction (f) $`\kappa =\mathrm{}`$ when the electron temperature equals the lattice temperature, e.g. when there is no energy transfer between the electron and phonon system. The cross-over bias voltage for the excess resistance is sensitive to the strength of the electron-phonon interaction and occurs from around $`\mathrm{\Delta }`$ (a) to around $`40\mathrm{\Delta }`$ (f) (according to (54)). The dependence on the electron-phonon interaction parameter $`\kappa `$ is rather weak as can be understood from (57). The local electron temperature in the middle of the ferromagnetic wire is proportional to $`\kappa ^{1/5}`$ and thus only has a very weak dependence on $`\kappa `$.
## V Discussion of experiments
In this section we discuss the connection of our results with experiments of Petrashov et al. and Giroud et al. . It will turn out that most of the experimental results can be understood on the basis of our calculations. Both experimental arrangements we will discuss below contain F-S junctions where the superconductor and the ferromagnet overlap in a certain region. The current redistribution in these junctions will play an important role in the following. Let us therefore introduce parameters characterizing these juntions: the resistance of the interface is called $`R_{\text{BN}}`$ in accordance with our previous consideration. Additionally $`R_{\text{SJ}}`$ will be the resistance of the superconducting part of the overlap junction in the normal state and $`R_{\text{FJ}}`$ the resistance of the ferromagnetic part of the overlap junction.
In the experiment by Giroud et al. a non-monotonic behaviour of the resistance below the superconducting critical temperature was observed. The sample consisted of a ferromagnetic wire, the resistance of which was measured in a 4-point arrangement. At some point a superconducting strip was on top of the wire. In a second sample two such strips were present and the resulting resistance change was twice as big as in the case of one strip. Since our formulation is based on a single interface and no coherent coupling between the two superconducting strips was found experimentally, we concentrate here on the sample with one strip. The resistance change in the two-strip sample is then simply twice that for the single strip sample. The experimental arrangement is such that in the region of the strip the current is redistributed among the ferromagnet and the superconductor. In Appendix A we introduce a simple quasi-one-dimensional model to calculate the effective resistance, these results being used for comparison with experiment. The resistance of the superconducting Al-strip is $`0.4\mathrm{\Omega }`$, the resistance of the ferromagnetic part below the strip is $`10\mathrm{\Omega }`$, and the resistance of the interface is estimated to be $`0.1\mathrm{\Omega }`$. Since the measured resistance change of the F-wire shows no signature of the vanishing of the resistance of the superconducting part, we believe that the real interface resistance is higher than estimated in Ref. , in particular higher than $`R_{\text{SJ}}`$. This yields a total resistance of $`R_{\text{eff}}=2(R_{\text{FJ}}R_{\text{BN}})^{1/2}`$, which is approximately of the order of a few $`\mathrm{\Omega }`$. A resistance change of the interface resistance $`\mathrm{\Delta }R_\text{B}`$ will than lead to an change of the effective resistance $`\mathrm{\Delta }R_{\text{eff}}=(R_{\text{FJ}}/R_{\text{BN}})^{1/2}\mathrm{\Delta }R_\text{B}`$, which in the case $`R_{\text{FJ}}>R_{\text{BN}}`$ is large than the resistance change of the interface resistance itself. For the experimental values we have $`(R_{\text{FJ}}/R_{\text{BN}})^{1/2}10`$ and thus a resistance change of $`0.2\mathrm{\Omega }`$, as observed in the experiment may result from a change of the interface resistance $`R_{\text{BN}}0.1\mathrm{\Omega }`$ by $`20\%`$.
The results of Petrashov et al. are more intriguing, since the magnitude of the measured resistance drop in some of the samples seems to be far too large to be explained without a ‘long-range’ proximity effect in the ferromagnet. We will concentrate here on three of the four samples discussed in Ref. . In these samples the transport through a long ferromagnetic wire with one ferromagnetic and one superconducting contact is studied, this is in contrast to experiments of Ref. . The geometry is such that the superconducting contact overlaps the ferromagnetic wire at one end and the current has to pass through a tiny piece of the superconductor. The three samples differ in the interface resistance. Two samples with a low interface resistance show large drops of the resistance of the order of $`8\mathrm{\Omega }`$ respectively $`16\mathrm{\Omega }`$ below the superconducting critical temperature. The third sample has a higher interface resistance ($`R_{\text{BN}}=41\mathrm{\Omega }`$) and shows a small resistance increase of the order of $`1.5\mathrm{\Omega }`$.
This agrees qualitatively with the results of our model. Indeed, the bigger resistance of the boundary usually means a formation of a ticker tunnel barrier such that the transmission eigenvalues are shifted towards zero. Our model does predict a resistance decrease for a fairly transparent interface and changes to an increase for a more tunnel-like interface. This is shown e.g. in Fig. 4. However, quantitatively one would expect that the resistance changes below the superconducting transition temperature are always of the order of the boundary resistance itself.
This is obviously not the case in the experiment with $`R_{\text{BN}}=41\mathrm{\Omega }`$, where the measured resistance change is about 40 times smaller. The first idea is that the resistance drop of the samples with better interface may possibly be accounted for by combining the effect of current redistribution and the apparent enhancement of Andreev reflection discussed in Sec.IV. Again we calculate the effective interface resistance in a quasi-one-dimensional model (see Appendix A). In the limit of a small interface resistance $`R_{\text{eff2}}=(R_{\text{FJ}}R_{\text{BN}})^{1/2}`$. A change of the interface resistance again results in an apparently larger change of the effective resistance $`\mathrm{\Delta }R_{\text{eff2}}=(R_{\text{FJ}}/R_{\text{BN}})^{1/2}\delta R_\text{B}/2`$. We may speculate that the large resistance drop observed in the experiment by Petrashov et al. can possibly be explained by this effect together with the observation made in IV that a spin dependent interface may cause another apparent enhancement of the Andreev-reflection.
There may be a more radical explanation for a small relative resistance change. In fact, the morphology of the metal-ferromagnet interfaces has not been yet sufficiently studied. The actual structure of the interface may be complicated. To illustrate how this can affect the results let us consider a simplistic model of a double interface. We speculate that a thin layer of magnetic alloy separates the ferromagnet and the superconductor. The boundary scattering then occurs in two stages: at the “inner” interface between the ferromagnet and the alloy and at the “outer” interface between the superconductor and the alloy. Since the proximity effect is quenched in magnets, the resistance of the “inner” interface is not affected by the superconducting transition whereas the resistance of the “outer” interface acquires a change described above. This leads to a smaller relative resistance change.
In Fig. 10 we show a comparison between the experimental results of Giroud et al. and our calculation for a contact with a transmission eigenvalue distribution of the model just described (see Appendix B). The maximal relative resistance change of the contact is $`16\%`$ in our calculation. According to above considerations this may be enhanced to a measured resistance change of the order of $`(R_{\text{FJ}}/R_{\text{BN}})^{1/2}16\%\delta R_\text{B}160\%\delta R_\text{B}`$ in agreement with measurements.
## VI Conclusions
To summarize, we have calculated the resistance of a diffusive ferromagnetic wire in contact with a superconducting reservoir in the linear and non-linear regime with purely elastic and inelastic scattering. It has been demonstrated that most of the recent experimental results can be understood in the absence of a superconducting proximity effect in the ferromagnet. Spin accumulation leads to an enhanced resistance below the superconducting transition temperature whereas Andreev reflection can lead to a decreased resistance below the superconducting transition temperature. The competition between these two mechanism determines the sign of the resistance change. The magnitude of the resistance change is of the order of the interface resistance or the spin-relaxation resistance. Electron heating can dramatically modify the nonlinear response resistance and change the bias dependence by orders of a magnitude.
###### Acknowledgements.
This work was partially supported by the “Stichting voor Fundamenteel Onderzoek der Materie” (FOM) and a Feodor Lynen Fellowship of the “Alexander von Humboldt-Stiftung” (W. B.) and the Norwegian Research Council (A. B.).
## A Current redistribution in an overlap junction
Here we introduce a quasi-one-dimensional model to account for the redistribution of the current under an overlap junction as used in the experiments. The two geometries we have in mind are depicted in Fig. 11.
To estimate the measured resistance we use the following quasi-one-dimensional model for the current redistribution in the overlap region of length $`d`$. The currents $`I_F(x)`$ in F and $`I_S(x)`$ in S in direction of the ferromagnetic wire follow from Ohm’s law
$`I_\text{F}(x)={\displaystyle \frac{d}{R_{\text{FJ}}}}{\displaystyle \frac{dU_\text{F}(x)}{dx}}`$ $`,`$ $`I_\text{S}(x)={\displaystyle \frac{d}{R_{\text{SJ}}}}{\displaystyle \frac{dU_\text{S}(x)}{dx}},`$ (A1)
where $`R_{\text{FJ(SJ)}}`$ is the resistance of the ferromagnetic part under (superconducting part above) the contact and $`U_{F(S)}`$ the respective voltage. Current conservation dictates
$$\frac{dI_\text{F}(x)}{dx}=\frac{dI_\text{S}(x)}{dx}=\frac{U_S(x)U_F(x)}{R_BNd},$$
(A2)
is the current per unit length through the contact resistance is $`R_{BN}`$. Boundary conditions are obviously that the total voltage drop is equal to $`V`$ and no current leaves the system through the boundary to vacuum.
Solving these equations for the geometry of Giroud et al. we find the effective resistance of this part to be
$`R_{\text{eff}}`$ $`=`$ $`{\displaystyle \frac{R_{\text{FJ}}}{R_{\text{FJ}}+R_{\text{S}J}}}`$ (A4)
$`\times \left(R_{\text{SJ}}+\sqrt{{\displaystyle \frac{4R_{\text{BN}}}{R_{\text{SJ}}+R_{\text{FJ}}}}}\mathrm{tanh}\left(\sqrt{{\displaystyle \frac{R_{\text{FJ}}+R_{\text{SJ}}}{4R_{\text{BN}}}}}\right)\right).`$
Of specific interest if the case that $`R_{\text{FJ}}R_{\text{BN}}R_{\text{SJ}}`$ in which case we obtain $`R_{\text{eff}}2(R_{\text{FJ}}R_{\text{BN}})^{1/2}`$.
A calculation similar to the previous for the geometry of Petrashov et al. leads to an effective measured resistance of
$$R_{\text{eff2}}=\frac{\sqrt{R_{\text{FJ}}R_{\text{BN}}}}{\mathrm{tanh}\sqrt{R_{\text{FJ}}/R_{\text{BN}}}}$$
(A5)
in the limit of vanishing resistance of the superconductor on top of the ferromagnet. The difference to the previous calculation is that here the current enters the junction through the ferromagnet, but has to leave the junction through the superconductor.
## B Transmission eigenvalues of a double interface
The distribution of transmission eigenvalues of a double interface as described in the text can be found with the technique described in Ref. . For details we refer to these articles. We model the double interface by a ballistic contact and a tunneling barrier in series. The tunneling barrier of conductance $`G_\text{T}`$ models the sharp drop in the potential due to the band structure mismatch, whereas the region close to that interface is treated as a collection of unit transmission channels with a total conductance $`G_{\text{QPC}}`$ The distribution of transmissions can be found from the solution of
$$I(\mathrm{\Phi })=G_\text{T}\mathrm{sin}(\mathrm{\Phi }\theta )=2G_{\text{QPC}}\mathrm{tan}\left(\frac{\theta }{2}\right).$$
(B1)
The distribution of transmission eigenvalues $`\rho (T)`$ is found by analytic continuation into the complex plane
$$\rho (T)=\frac{1}{e^2}\frac{1}{T\sqrt{1T}}\text{Re}\left[I\left(\pi +2i\text{acosh}\frac{1}{\sqrt{T}}\right)\right].$$
(B2)
The dependence on the two separate conductances may be eliminated in favor of the total conductance of the contact $`G_{\text{BN}}=G_\text{T}G_{\text{QPC}}/(G_\text{T}+G_{\text{QPC}})`$ and the ratio of the two $`\alpha =G_\text{T}/2G_{\text{QPC}}`$. The transmission eigenvalue distribution then only depends on $`\alpha `$. It is plotted in Fig. 12 for several values of $`\alpha `$. For small values of $`\alpha `$ the contact is dominated by the tunnel barrier resulting in a shift of the transmission eigenvalues to lower values and a gap above a certain $`T`$. Higher $`\alpha `$ shift the distribution to larger transmission eigenvalues and a gap opens up for low transmission eigenvalues. For a range $`0.1\alpha 0.5`$ the distribution restricted to a finite interval of transmission eigevalues. At even higher values of $`\alpha `$ the upper gap closes and the distribution becomes more and more peaked at $`T=1`$.
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# 1 Introduction
## 1 Introduction
The prospective high-energy $`e^+e^{}`$ linear collider TESLA is being designed to operate on top of the $`Z`$ boson resonance,
$$e^+e^{}Z,$$
(1)
by adding a bypass to the main beam line . Given the high luminosity, $`=7\times 10^{33}\mathrm{cm}^2\mathrm{s}^1`$, and the cross section, $`\sigma _Z30\mathrm{nb}`$ (including radiative corrections), about $`2\times 10^9Z`$ events can be generated in an operational year of $`10^7\mathrm{s}`$. We will therefore refer to this option as the GigaZ mode of the machine. Moreover, by increasing the collider energy to the $`W`$-pair threshold,
$$e^+e^{}W^+W^{},$$
(2)
about $`10^6`$ $`W`$ bosons can be generated at the optimal energy point for measuring the $`W`$ boson mass, $`M_W`$, near threshold and about $`3\times 10^6`$ $`W`$ bosons at the energy of maximal cross section . The large increase in the number of $`Z`$ events by two orders of magnitude as compared to LEP 1 and the increasing precision in the measurements of $`W`$ boson properties, open new opportunities for high precision physics in the electroweak sector .
By adopting the Blondel scheme for running $`e^+e^{}`$ colliders with polarized beams, the left-right asymmetry, $`A_{LR}2(14\mathrm{sin}^2\theta _{\mathrm{eff}})/(1+(14\mathrm{sin}^2\theta _{\mathrm{eff}})^2)`$, can be measured with very high precision, $`\delta A_{LR}\pm 10^4`$ , when both electrons and positrons are polarized longitudinally. This accuracy can be achieved since the total cross section, the left-right asymmetry and the polarization factor, $`𝒫=(P_++P_{})/(1+P_+P_{})`$, can be measured by individually flipping the electron and positron helicities, generating all $`2\times 2`$ spin combinations in $`\sigma _{ij}(i,j=L,R)`$; only the difference between the moduli $`|P_+|`$ and $`|P_{}|`$ before and after flipping the polarizations of both the positron and electron beams need to be monitored by laser Compton scattering. From $`A_{LR}`$ the mixing angle in the effective leptonic vector coupling of the on-shell $`Z`$ boson, $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$, can be determined to an accuracy
$$\delta \mathrm{sin}^2\theta _{\mathrm{eff}}\pm \mathrm{\hspace{0.17em}1}\times 10^5,$$
(3)
while the $`W`$ boson mass is expected to be measurable to a precision of
$$\delta M_W\pm \mathrm{\hspace{0.17em}6}\mathrm{MeV},$$
(4)
by scanning the $`W^+W^{}`$ threshold .
Besides the improvements in $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$ and $`M_W`$, GigaZ has the potential to determine the total $`Z`$ width within $`\delta \mathrm{\Gamma }_Z=\pm 1`$ MeV; the ratio of hadronic to leptonic partial $`Z`$ widths with a relative uncertainty of $`\delta R_l/R_l=\pm 0.05\%`$; the ratio of the $`b\overline{b}`$ to the hadronic partial widths with a precision of $`\delta R_b=\pm 1.4\times 10^4`$; and to improve the $`b`$ quark asymmetry parameter $`A_b`$ to a precision of $`\pm 1\times 10^3`$ . These additional measurements offer complementary information on the Higgs boson mass, $`M_H`$, but also on the strong coupling constant, $`\alpha _s`$, which enters the radiative corrections in many places. This is desirable in its own right, and in the present context it is important to control $`\alpha _s`$ effects from higher order loop contributions to avoid confusion with Higgs effects. Indirectly, a well known $`\alpha _s`$ would also help to control $`m_t`$ effects, since $`m_t`$ from a threshold scan at a linear collider will be strongly correlated with $`\alpha _s`$. We find that via a precise measurement of $`R_l`$, GigaZ would provide a clean determination of $`\alpha _s`$ with small error
$$\delta \alpha _s\pm \mathrm{\hspace{0.17em}0.001},$$
(5)
allowing to reduce the error of the top-quark mass from the threshold scan. The anticipated precisions for the most relevant electroweak observables at the Tevatron (Run IIA and IIB), the LHC, a future linear collider, LC, and GigaZ are summarized in Tab. 1.
In this note, we study the potential impact of such measurements on the parameters of the Standard Model (SM) and its minimal supersymmetric extension (MSSM) . Higgs boson masses and SUSY particle masses affect the high precision observables through loop corrections. These loop corrections are evaluated in this note at the presently available level of theoretical accuracy, still leaving many refinements to be worked out in the coming years . Even though a complete set of calculations is lacking at the present time, the essential features of the GigaZ physics potential can nevertheless be studied in first exploratory steps. In Ref. it has been demonstrated that very stringent consistency tests of the SM and the MSSM will become feasible with the GigaZ precision, and the prospects for $`b`$ physics at GigaZ have been discussed. The latter topic has been studied in more detail in Ref. . In the present note, we will focus in a systematic way on the Higgs sectors of the SM and the MSSM, and also on the scalar top sector of the MSSM.
## 2 Higgs Sector of the SM
In the canonical form of the SM, the precision observables measured at the $`Z`$ peak are affected by two high mass scales in the model: the top quark mass, $`m_t`$, and the Higgs boson mass, $`M_H`$. They enter as virtual states in loop corrections to various relations between electroweak observables. For example, the radiative corrections entering the relation between $`M_W`$ and $`M_Z`$, and between $`M_Z`$ and $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$, have a strong quadratic dependence on $`m_t`$ and a logarithmic dependence on $`M_H`$. We mainly focus on the two electroweak observables that are expected to be measurable with the highest accuracy at GigaZ, $`M_W`$ and $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$. Our analysis is based on the results for the electroweak precision observables including higher order electroweak and QCD corrections. The current theoretical uncertainties are dominated by the parametric uncertainties from the errors in the input parameters $`m_t`$ (see Tab. 1) and $`\mathrm{\Delta }\alpha `$. The latter denotes the QED-induced shift in the fine structure constant, $`\alpha \alpha (M_Z)`$, originating from charged-lepton and light-quark photon vacuum polarization diagrams. The hadronic contribution to $`\mathrm{\Delta }\alpha `$ currently introduces an uncertainty of $`\delta \mathrm{\Delta }\alpha =\pm 2\times 10^4`$ . Forthcoming low-energy $`e^+e^{}`$annihilation experiments may reduce this uncertainty to about $`\pm 5\times 10^5`$ . Combining this value with future (indistinguishable) errors from unknown higher order corrections, we assign the total uncertainty of $`\delta \mathrm{\Delta }\alpha =\pm 7\times 10^5`$ to $`\mathrm{\Delta }\alpha `$, which is used throughout the paper unless otherwise stated. For the future theoretical uncertainties from unknown higher-order corrections (including the uncertainties from $`\delta \mathrm{\Delta }\alpha `$) we assume,
$$\delta M_W(\text{theory})=\pm 3\mathrm{MeV},\delta \mathrm{sin}^2\theta _{\mathrm{eff}}(\text{theory})=\pm 3\times 10^5\text{(future)}.$$
(6)
Given the high precision of GigaZ, also the experimental error in $`M_Z`$, $`\delta M_Z=\pm 2.1\mathrm{MeV}`$ , results in non-negligible uncertainties of $`\delta M_W=\pm 2.5\mathrm{MeV}`$ and $`\delta \mathrm{sin}^2\theta _{\mathrm{eff}}=\pm 1.4\times 10^5`$. The experimental error in the top-quark mass, $`\delta m_t=\pm 130\mathrm{MeV}`$, induces further uncertainties of $`\delta M_W=\pm 0.8\mathrm{MeV}`$ and $`\delta \mathrm{sin}^2\theta _{\mathrm{eff}}=\pm 0.4\times 10^5`$. Thus, while currently the experimental error in $`M_Z`$ can safely be neglected, for the GigaZ precision it will actually induce an uncertainty in the prediction of $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$ that is larger than its experimental error.
(a) The relation between $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$ and $`M_Z`$ can be written as
$$\mathrm{sin}^2\theta _{\mathrm{eff}}\mathrm{cos}^2\theta _{\mathrm{eff}}=\frac{A^2}{M_Z^2(1\mathrm{\Delta }r_Z)},$$
(7)
where $`A=[(\pi \alpha )/(\sqrt{2}G_F)]^{1/2}=37.2805(2)`$ GeV is a combination of two precisely known low-energy coupling constants, the Fermi constant, $`G_F`$, and the electromagnetic fine structure constant, $`\alpha `$. The quantity $`\mathrm{\Delta }r_Z`$ summarizes the loop corrections, which at the one-loop level can be decomposed as
$$\mathrm{\Delta }r_Z=\mathrm{\Delta }\alpha \mathrm{\Delta }\rho ^\mathrm{t}+\mathrm{\Delta }r_Z^\mathrm{H}+\mathrm{}.$$
(8)
The leading top contribution to the $`\rho `$ parameter , quadratic in $`m_t`$, reads
$$\mathrm{\Delta }\rho ^\mathrm{t}=\frac{3G_Fm_t^2}{8\pi ^2\sqrt{2}}.$$
(9)
The Higgs boson contribution is screened, being logarithmic for large Higgs boson masses
$$\mathrm{\Delta }r_Z^\mathrm{H}=\frac{G_FM_W^2}{8\pi ^2\sqrt{2}}\frac{1+9s_W^2}{3c_W^2}\mathrm{log}\frac{M_H^2}{M_W^2}+\mathrm{}.$$
(10)
(b) An independent analysis can be based on the precise measurement of $`M_W`$ near threshold. The $`M_W`$$`M_Z`$ interdependence is given by,
$`{\displaystyle \frac{M_W^2}{M_Z^2}}\left(1{\displaystyle \frac{M_W^2}{M_Z^2}}\right)`$ $`=`$ $`{\displaystyle \frac{A^2}{M_Z^2(1\mathrm{\Delta }r)}},`$ (11)
where the quantum correction $`\mathrm{\Delta }r`$ has the one-loop decomposition,
$`\mathrm{\Delta }r`$ $`=`$ $`\mathrm{\Delta }\alpha {\displaystyle \frac{c_W^2}{s_W^2}}\mathrm{\Delta }\rho ^\mathrm{t}+\mathrm{\Delta }r^\mathrm{H}+\mathrm{},`$ (12)
$`\mathrm{\Delta }r^\mathrm{H}`$ $`=`$ $`{\displaystyle \frac{G_FM_W^2}{8\pi ^2\sqrt{2}}}{\displaystyle \frac{11}{3}}\mathrm{log}{\displaystyle \frac{M_H^2}{M_W^2}}+\mathrm{},`$ (13)
with $`\mathrm{\Delta }\alpha `$ and $`\mathrm{\Delta }\rho ^\mathrm{t}`$ as introduced above.
Owing to the different dependences of $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$ and $`M_W`$ on $`m_t`$ and $`M_H`$, the high precision measurements of these quantities at GigaZ (combined with the other supplementary electroweak observables) can determine the mass scales $`m_t`$ and $`M_H`$. The expected accuracy in the indirect determination of $`M_H`$ from the radiative corrections within the SM is displayed in Fig. 1. To obtain these contours, the error projections in Tab. 1 are supplemented by central values equal to the current SM best fit values for the entire set of current high precision observables . For the theoretical uncertainties, Eq. (6) is used, while the parametric uncertainties, such as from $`\alpha _s`$ and $`M_Z`$, are automatically accounted for in the fits. The allowed bands in the $`m_t`$$`M_H`$ plane for the GigaZ accuracy are shown separately for $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$ and $`M_W`$. By adding the information on the top-quark mass, with $`\delta m_t\stackrel{<}{}\mathrm{\hspace{0.33em}130}`$ MeV obtained from measurements of the $`t\overline{t}`$ production cross section near threshold, an accurate determination of the Higgs boson mass becomes feasible from both, $`M_W`$ and $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$. If the two values are found to be consistent, they can be combined and compared to the Higgs boson mass measured in direct production through Higgs-strahlung (see the last row in Tab. 1). In Fig. 1 this is shown by the shaded area labeled as “GigaZ (1$`\sigma `$ errors)”, where the measurements of other $`Z`$ boson properties as anticipated for GigaZ are also included (the best fit value for $`m_t`$ is assumed to coincide with the central $`m_t`$ value in Fig. 1). For comparison, the area in the $`m_t`$$`M_H`$ plane corresponding to the current experimental accuracies, labeled as “now (1$`\sigma `$ errors)”, is also shown.
The results can be summarized by calculating the accuracy with which $`M_H`$ can be determined indirectly. The expectations for $`\delta M_H/M_H`$ in each step until GigaZ are collected in Tab. 2. Extending an earlier analysis , where $`\delta M_H/M_H`$ was separately determined from $`M_W`$ and $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$, we additionally employ here the full set of precision observables for our analysis. Concerning the experimental errors, the values from Tab. 1 are taken. For the uncertainty in $`\mathrm{\Delta }\alpha `$, we use $`\delta \mathrm{\Delta }\alpha =\pm 7\times 10^5`$ (yielding Eq. (6) upon combination with our estimate of unknown higher order corrections), except for the first row displaying the present situation, where $`\delta \mathrm{\Delta }\alpha =\pm 2\times 10^4`$ is employed.
It is apparent that GigaZ, reaching $`\delta M_H/M_H=\pm 7\%`$, triples the precision in $`M_H`$ relative to the anticipated LHC status. On the other hand, a linear collider without the high luminosity option would provide only a modest improvement.
(c) A direct formal relation between $`M_W`$ and $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$ can be established by combining the two relations Eqs. (7) and (11) as
$$M_W^2=\frac{A^2}{\mathrm{sin}^2\theta _{\mathrm{eff}}(1\mathrm{\Delta }r_W)}.$$
(14)
The quantum correction $`\mathrm{\Delta }r_W`$ is independent of $`\mathrm{\Delta }\rho ^\mathrm{t}`$ in leading order and has the one-loop decomposition
$`\mathrm{\Delta }r_W`$ $`=`$ $`\mathrm{\Delta }\alpha \mathrm{\Delta }r_W^\mathrm{H}+\mathrm{},`$ (15)
$`\mathrm{\Delta }r_W^\mathrm{H}`$ $`=`$ $`{\displaystyle \frac{G_FM_Z^2}{24\pi ^2\sqrt{2}}}\mathrm{log}{\displaystyle \frac{M_H^2}{M_W^2}}+\mathrm{}.`$ (16)
Relation (14) can be evaluated by inserting the measured value of the Higgs boson mass as predetermined at the LHC and the LC. This is visualized in Fig. 2, where the present and future theoretical predictions for $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$ and $`M_W`$ (for different values of $`M_H`$) are compared with the experimental accuracies at various colliders. Besides the independent predictions of $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$ and $`M_W`$ within the SM, the $`M_W\mathrm{sin}^2\theta _{\mathrm{eff}}`$ contour plot in Fig. 2 can be interpreted as an additional indirect determination of $`M_W`$ from the measurement of $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$. Given the expected negligible error in $`M_H`$, this results in an uncertainty of
$$\delta M_W(\mathrm{indirect})\pm 2\mathrm{MeV}\pm 3\mathrm{MeV}.$$
(17)
The first uncertainty reflects the experimental error in $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$, while the second is the theoretical uncertainty discussed above (see Eq. (6)). The combined uncertainty of this indirect prediction is about the same as the one of the SM prediction according to Eq. (11) and is close to the experimental error expected from the $`W^+W^{}`$ threshold given in Eq. (4).
Consistency of all the theoretical relations with the experimental data would be the ultimate precision test of the SM based on quantum fluctuations. The comparison between theory and experiment can also be exploited to constrain possible physics scales beyond the SM. These additional contributions can conveniently be described in terms of the S,T,U or $`ϵ`$ parameters . Adopting the notation of Ref. , the errors with which they can be measured at GigaZ are given as follows:
$$\begin{array}{cc}\mathrm{\Delta }S=\pm 0.05,\hfill & \hfill \mathrm{\Delta }\widehat{ϵ}_3=\pm 0.0004,\\ \mathrm{\Delta }T=\pm 0.06,\hfill & \hfill \mathrm{\Delta }\widehat{ϵ}_1=\pm 0.0005,\\ \mathrm{\Delta }U=\pm 0.04,\hfill & \hfill \mathrm{\Delta }\widehat{ϵ}_2=\pm 0.0004.\end{array}$$
(18)
The oblique parameters in Eq. (18) are strongly correlated. On the other hand, many types of new physics predict $`U=\widehat{ϵ}_2=0`$ or very small (see Ref. and references therein). With the $`U`$ ($`\widehat{ϵ}_2`$) parameter known, the anticipated errors in $`S`$ and $`T`$ would decrease to about $`\pm 0.02`$, while the errors in $`\widehat{ϵ}_1`$ and $`\widehat{ϵ}_3`$ would be smaller than $`\pm 0.0002`$.
In the context of a spontaneously broken gauge theory, the above mentioned comparisons shed light on the basic theoretical components for generating the masses of the fundamental particles. On the other hand, an observed inconsistency would be a clear indication for the existence of a new physics scale.
## 3 Supersymmetry
The second step in this GigaZ analysis is based on the assumption that supersymmetry would be discovered at LEP 2, the Tevatron, or the LHC, and further explored at an $`e^+e^{}`$ linear collider. The high luminosity expected at TESLA can be exploited to determine supersymmetric particle masses and mixing angles with errors from $`𝒪(1\%)`$ down to one per mille , provided they reside in the kinematical reach of the collider, which we assume to be about 1 TeV. In this context we will address two problems arising in the Higgs sector and the scalar top sector within the MSSM.
For the SUSY contributions to $`M_W`$ and $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$ we use the complete one-loop results in the MSSM as well as the leading higher order QCD corrections . The recent electroweak two-loop results of the SM part in the MSSM have not been taken into account, since no genuine MSSM counterpart is available so far. As above, concerning the future theoretical uncertainties of $`M_W`$ and $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$ we use Eq. (6).
In contrast to the Higgs boson mass in the SM, the lightest $`𝒞𝒫`$-even MSSM Higgs boson mass, $`M_h`$, is not a free parameter but can be calculated from the other SUSY parameters. In the present analysis, the currently most precise result based on Feynman-diagrammatic methods is used, relating $`M_h`$ to the pseudoscalar Higgs boson mass, $`M_A`$. The numerical evaluation has been performed with the Fortran code FeynHiggs . In our analysis we assume a future uncertainty in the theoretical prediction of $`M_h`$ of $`\pm 0.5\mathrm{GeV}`$.
(a) The relation between $`M_W`$ and $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$ is affected by the parameters of the supersymmetric sector, especially the $`\stackrel{~}{t}`$-sector. At the LHC and especially at a prospective LC, the mass of the light $`\stackrel{~}{t}`$, $`m_{\stackrel{~}{t}_1}`$, and the $`\stackrel{~}{t}`$-mixing angle, $`\theta _{\stackrel{~}{t}}`$, may be measurable very well, particularly in the process $`e^+e^{}\stackrel{~}{t}_1\overline{\stackrel{~}{t}_1}`$ . On the other hand, background problems at the LHC and insufficient energy at the LC may preclude the analysis of the heavy $`\stackrel{~}{t}`$-particle, $`\stackrel{~}{t}_2`$.
In Fig. 3 it is demonstrated how in this situation limits on $`m_{\stackrel{~}{t}_2}`$ can be derived from measurements of $`M_h`$, $`M_W`$ and $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$. As experimental values we assumed $`M_h=115\mathrm{GeV}`$, $`M_W=80.40\mathrm{GeV}`$ and $`\mathrm{sin}^2\theta _{\mathrm{eff}}=0.23140`$, with the experimental errors given in the last column of Tab. 1. We consider two cases for $`\mathrm{tan}\beta `$, the ratio of the vacuum expectation values of the two Higgs doublets in the MSSM: the low $`\mathrm{tan}\beta `$ region, where we assume a band, $`2.5<\mathrm{tan}\beta <3.5`$, and the high $`\mathrm{tan}\beta `$ region where we assume a lower bound, $`\mathrm{tan}\beta 10`$, as can be expected from measurements in the gaugino sector (see e.g. Ref. ). As for the other parameters, the following values are assumed, with uncertainties as expected from LHC and TESLA : $`m_{\stackrel{~}{t}_1}=500\pm 2\mathrm{GeV}`$, $`\mathrm{sin}\theta _{\stackrel{~}{t}}=0.69\pm 0.014`$, $`A_b=A_t\pm 10\%`$, $`\mu =200\pm 1\mathrm{GeV}`$, $`M_2=400\pm 2\mathrm{GeV}`$ and $`m_{\stackrel{~}{g}}=500\pm 10\mathrm{GeV}`$. ($`A_{b,t}`$ are trilinear soft SUSY-breaking parameters, $`\mu `$ is the Higgs mixing parameter, $`M_2`$ is one of the soft SUSY-breaking parameter in the gaugino sector, and $`m_{\stackrel{~}{g}}`$ denotes the gluino mass.)
For low $`\mathrm{tan}\beta `$ the heavier $`\stackrel{~}{t}`$-mass, $`m_{\stackrel{~}{t}_2}`$, can be restricted to $`760\mathrm{GeV}\stackrel{<}{}m_{\stackrel{~}{t}_2}\stackrel{<}{}\mathrm{\hspace{0.33em}930}\mathrm{GeV}`$ from the $`M_h`$, $`M_W`$ and $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$ precision measurements. The mass $`M_A`$ varies between $`200\mathrm{GeV}`$ and $`1600\mathrm{GeV}`$. A reduction of this interval to $`M_A500\mathrm{GeV}`$ by its non-observation at the LHC and the LC does not improve the bounds on $`m_{\stackrel{~}{t}_2}`$. If $`\mathrm{tan}\beta 10`$, the second theoretically preferred range , the allowed region turns out to be much smaller ($`660\mathrm{GeV}\stackrel{<}{}m_{\stackrel{~}{t}_2}\stackrel{<}{}\mathrm{\hspace{0.33em}680}\mathrm{GeV}`$), and the mass $`M_A`$ is restricted to $`M_A\stackrel{<}{}\mathrm{\hspace{0.33em}800}\mathrm{GeV}`$. In deriving the bounds on $`m_{\stackrel{~}{t}_2}`$, both the constraints from $`M_h`$ (see Ref. ) and $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$ play an important role. For the bounds on $`M_A`$, the main effect comes from $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$. We have assumed a value for $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$ slightly different from the corresponding value obtained in the SM limit. For this value the (logarithmic) dependence on $`M_A`$ is still large enough so that in combination with the high precision in $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$ at GigaZ an upper limit on $`M_A`$ can be set. For an error as obtained at an LC without the GigaZ mode (see Tab. 1) no bound on $`M_A`$ could be inferred.
(b) A similar problem of high interest occurs in the sector of the MSSM Higgs particles. It is well known, that the heavy Higgs bosons $`A`$, $`H`$ and $`H^\pm `$, are increasingly difficult to observe at the LHC with rising mass . At $`e^+e^{}`$ linear colliders heavy Higgs particles are produced primarily in pairs $`(HA)`$ and $`(H^+H^{})`$ so that they cannot be analyzed for mass values beyond the beam energy of $`500\mathrm{GeV}`$ in the first phase of such a machine. It has been demonstrated though that the ratio of the decay branching ratios of the light Higgs boson $`h`$ is sensitive to $`M_A`$ up to values of 700 GeV to 1 TeV . Since any such analysis is difficult, it is suggestive to search for complementary channels in which new limits may be derived from other high precision measurements.
The result of such a study is presented in Fig. 4, based on the expected errors for $`M_h`$, $`m_{\stackrel{~}{t}_1}`$, and $`\theta _{\stackrel{~}{t}}`$ from LC measurements, and assuming either a rough measurement of the heavy $`\stackrel{~}{t}`$-mass, $`m_{\stackrel{~}{t}_2}`$, at the LHC, or a precise determination of $`m_{\stackrel{~}{t}_2}`$ at an LC. Fig. 4 shows the exclusion contours in the $`M_A\mathrm{tan}\beta `$ plane based on the following scenario (inspired by the mSUGRA(1) reference scenario studied e.g. in Ref. ): $`M_h=110\pm 0.05\mathrm{GeV}`$ from LC measurements, $`M_W=80.400\pm 0.006\mathrm{GeV}`$ and $`\mathrm{sin}^2\theta _{\mathrm{eff}}=0.23138\pm 1\times \mathrm{\hspace{0.17em}10}^5`$ from GigaZ measurements, $`m_{\stackrel{~}{t}_1}=340\pm 1\mathrm{GeV}`$, $`\mathrm{sin}\theta _{\stackrel{~}{t}}=0.69\pm 0.014`$ from the LC, $`m_{\stackrel{~}{t}_2}=640\pm 10\mathrm{GeV}`$ from the LHC or alternatively $`m_{\stackrel{~}{t}_2}=520\pm 1\mathrm{GeV}`$ from LC measurements; furthermore $`A_b=640\pm 60\mathrm{GeV}`$, $`\mu =316\pm 1\mathrm{GeV}`$, $`M_2=152\pm 2\mathrm{GeV}`$, $`m_{\stackrel{~}{g}}=496\pm 10\mathrm{GeV}`$ based on LHC or LC runs.
If the scenario with lower $`\stackrel{~}{t}`$ masses is realized, this would, due to the dependence of the Higgs boson mass on the MSSM parameters, correspond to higher values for both $`M_A`$ and $`\mathrm{tan}\beta `$. Despite the precise measurement of $`m_{\stackrel{~}{t}_2}`$ up to $`1\mathrm{GeV}`$ at the LC, in the example considered here the restrictions placed on the Higgs sector would be relatively weak. The constraints would be $`450\mathrm{GeV}\stackrel{<}{}M_A\stackrel{<}{}\mathrm{\hspace{0.33em}1950}\mathrm{GeV}`$ and $`6\stackrel{<}{}\mathrm{tan}\beta \stackrel{<}{}\mathrm{\hspace{0.33em}10}`$. However, if the higher $`\stackrel{~}{t}_2`$ mass is realized in nature, corresponding to larger mixing in the $`\stackrel{~}{t}`$ sector, in spite of the relatively rough measurement of $`m_{\stackrel{~}{t}_2}`$, in our example the allowed parameter range is reduced to $`250\mathrm{GeV}\stackrel{<}{}M_A\stackrel{<}{}\mathrm{\hspace{0.33em}1200}\mathrm{GeV}`$ and $`2.5\stackrel{<}{}\mathrm{tan}\beta \stackrel{<}{}\mathrm{\hspace{0.33em}3.5}`$. Again, $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$ plays an important role (cf. discussion of Fig. 3); without the high precision measurement of $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$ no upper limit on $`M_A`$ could be set.
Thus the high precision measurements of $`M_W`$, $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$ and $`M_h`$ do not improve on the direct lower bound on the mass of the pseudoscalar Higgs boson $`A`$. Instead they enable us to set an upper bound on this basic parameter of the supersymmetric Higgs sector, derived from the requirement of consistency of the electroweak precision data with the MSSM.
## 4 Conclusions
The opportunity to measure electroweak observables very precisely in the GigaZ mode of the prospective $`e^+e^{}`$ linear collider TESLA, in particular the electroweak mixing angle $`\mathrm{sin}^2\theta _{\mathrm{eff}}`$ and the $`W`$ boson mass, opens new areas for high precision tests of electroweak theories. We have analyzed in detail two generic examples: (i) The Higgs mass of the Standard Model can be extracted to a precision of a few percent from loop corrections. By comparison with the direct measurements of the Higgs mass, bounds on new physics scales can be inferred that may not be accessible directly. (ii) The masses of particles in supersymmetric theories, which for various reasons may not be accessible directly neither at the LHC nor at the LC, can be constrained. Typical examples are the heavy Higgs bosons and the heavy scalar top quark. In the scenarios studied here, a sensitivity of up to order 2 TeV for the mass of the pseudoscalar Higgs boson and an upper bound of 1 TeV for the heavy scalar top quark could be expected from the virtual loop analyses of the high precision data.
Opening windows to unexplored energy scales renders these analyses of virtual effects an important tool for experiments in the GigaZ mode of a future $`e^+e^{}`$ linear collider.
## Acknowledgements
We thank P. Langacker for a critical reading of the manuscript and A. Hoang and T. Teubner for helpful discussions. Parts of the calculation have been performed on the QCM cluster at the University of Karlsruhe, supported by the Deutsche Forschungsgemeinschaft (Forschergruppe “Quantenfeldtheorie und Computeralgebra”).
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# Non-linear evolution of step meander during growth of a vicinal surface with no desorption.
## I Introduction
The production of solids by Molecular Beam Epitaxy (MBE) having a surface which is abrupt on the atomic scale is often hampered either by a stochastic roughness or due to the presence of morphological instabilities. The stochastic roughness is often attributed to shot noise from the incoming deposition flux. As for determinitic instabilities, there are three general types of surface instabilities leading to kinetic roughnesses: step-bunching, step meandering, and islanding (see Fig.1). The two first categories are met on vicinal surfaces, while the last one can either be present on high symmetry surfaces, a typical mechanism being the Ehrlich-Schwoebel effect, or even on a vicinal surface as a secondary instability of the step meander.
Kinetic roughening has long been as a mystery. With regard to MBE growth on a high symmetry surface, a prominent example is the Kardar-Parisi-Zhang equation introduced in an attempt to describe surface noise-induced-roughening, its one dimensional version is:
$`_ty=a\mathrm{\Omega }F+_{xx}y+\left(_xy\right)^2+\eta ,`$ (1)
where $`y`$ is the surface height, and $`x`$ the coordinate along the surface(1). Derivatives are subscripted, that is $`_ty=y/t`$ and so on. $`F`$ is the incoming flux, $`a`$ is the atomic height and $`\mathrm{\Omega }`$ is the atomic area. We have set the coefficients to unity, since only the form of the equation matters in this discussion. This equation has given rise to a variety of investigations both analytically and numerically (this part has included analytical treatment of the partial differential equation together with numerical Monte-Carlo simulations which mimics KPZ dynamics). Equation (1) is phenomenological in the sense that it is derived on the basis of symmetries.
The KPZ nonlinearity is a natural candidate in the long wavelength limit, if desorption is present, or if allowance is made for defects (such as vacancies–often named overhangs– in the growing solid). No derivation of that equation has been given so far, however. The reason is, in our opinion, the lack of a continuum description of island nucleation. In the absence of both defects and desorption (which are two usual requirements for production of solids of interest!), the nonlinear KPZ term is not permissible. In that case the equation must have a form of a conservation law, that is:
$`_ty=a\mathrm{\Omega }F.𝐉,`$ (2)
oo that upon averaging, the mean velocity is simply be given by $`a\mathrm{\Omega }F`$, as it should be; because the KPZ nonlinearity can not be written as a flux (i.e. as a divergence of a current) it gives an additional contribution to the grwoth velocity by a quantity $`<(_xy)^2>0`$ (the symbol $`<..>`$ stands for to the average), which obviously makes no sense. There has thus been a variety of attempts with the aim of deriving the appropriate surface evolution equation in that limit. Here again no derivation from first principles is available.
As said above, in addition to surface roughness caused by shot noise, nominal high symmetry as well well as vicinal surfaces, may become inherently unstable when brought away from equilibrium. Nominal surfaces may develop mounds due to the ES effect. However, derivation of the appropriate surface evolution equation in that case is still a matter of debate, though a significant progress has been achieved.
In contrast to nominal surfaces, vicinal surfaces in the step flow regime have allowed to derive evolution equations from first principles. In a series of papers, we have shown that vicinal surfaces offer a relatively tractable situation, though often nontrivial, where evolution equations can be extracted from basic transport and kinetic laws. The strategy is to first focus on derivation of step evolution equations. Once this task is achieved, it becomes then possible to derive the surface evolution equation. In its general form, the evolution equation is nonlocal and highly nonlinear. A rather simple information is extracted, however, if we focus on the long-wavelength limit: that is we assume that the wavelength of the step meander, and thus surface modulation, is small in comparison to the natural physical length (diffusion length if desorption is important, otherwise the interstep distance, which is the most frequent situation). More precisely the full growth equations, which are highly nonlinear and nonlocal, can be reduced to nonlinear partial differential equations, which are more tractable and often allow a signficant analytical progress as wil be shown here.
As shown by Bales and Zangwill, a straight step during MBE growth may become morphologically unstable in the presence of an attachment asymmetry (the Ehrlich-Schwoebel (ES) effect) at the step. Close to the instability threshold, starting from the Burton-Cabrera-Frank (BCF) model, we have shown that the step profile in the presence of desorption obeys the Kuramoto-Sivahinsky equation (written in a canonical form):
$`_t\zeta =_{xx}\zeta _{xxxx}\zeta +\left(_x\zeta \right)^2,`$ (3)
where $`x`$ is the coordinate along the step (Fig.1), and $`\zeta `$ designates the step position. In a similar fashion we have shown later that steps on a vicinal surface obey a set of coupled anisotropic Kuramoto-Sivahinsky equations. The ultimate stage of surface dynamics is found to be spatiotemporal chaos. Two remarks are in order: (i) the KPZ nonlinearity is of the KS type –due to desorption– (ii) the first term in the KS equation has a negative sign, signalling an instability; there is a necessity for taking higher order derivatives into account in order to prevent arbitrary short wavelength modes to develop.
A question of major importance arose recently: if desorption is negligible, what kind of nonlinearity should one expect? because of the conserved character of dynamics, only terms which can be written as derivatives of a current are allowed. We could naively have thought that a natural candidate would be the conserved KS equation, namely $`_t\zeta =_{xx}\zeta _{xxxx}\zeta +_{xx}[\left(_x\zeta \right)^2]`$. A close inspection of the BCF equations, as shown here in details, reveals that this is not the case, though symmetry and conservation would dictate that form as the first plausible candidate. We have recently shown that, for an in-phase train of steps, each step obeys the following nontrivial evolution equation:
$`_t\zeta =_x\left[{\displaystyle \frac{1}{1+(_x\zeta )^2}}\left(_x\zeta +_x\left({\displaystyle \frac{_{xx}\zeta }{(1+(_x\zeta )^2)^{3/2}}}\right)\right)\right].`$ (4)
This highly nonlinear equation could not be inferred from scaling and symmetry arguments. It is related to a singular behavior of the amplitude of the meander that behaves as $`1/F^{1/2}`$ when $`F`$ is small. Instead of chaos, a regular pattern is revealed, the modulation wavelength is fixed at the very initial stages while the amplitude of the step deformation follows a scaling law $`wt^{1/2}`$.
The objective of this paper is many fold. We first give an extensive derivation of the above evolution equation starting from the BCF model. We shall also present a general argument on why that singular behavior is present in the absence of desorption. A second line of the investigation concerns higher order contributions. It is clear that the above equation enjoys the up-down symmetry, $`\zeta \zeta `$. We show here that the effect of higher order contributions is to destroy this up-down symmetry. An important fact to be presented here is that the step profile exhibits a plateau-like morphology. This contradicts the preliminary simulation given in Ref. . That simulation had suffered from numerical inaccuracy causing spurious spikes to develop. Finally we show that though the full equation is highly nonlinear it has been possible to provide a quantitative analytical treatment for the step morphology, and evaluate the plateau width along with the meander amplitude. The results are found to be in good agreement with numerical results.
This paper is organized as follow. We write down the basic equations in section II. In section III a linear stability analysis is performed, which allows to evaluate the most unstable wavelength and the typical time for the appearance of the instability. In section IV we shall provide a general argument on the extraction of the scaling of the step position with the incoming flux. Section V presents the detailed derivation of the principal evolution equation (4) in the one-sided limit. We shall then present the situation where there is a finite ES barrier. Section VI deals with the higher order terms and their impact on the up-down symmetry. In section VII we generalize the derivation of the step evolution equations to the two-sided case. Discussion and outlook are presented in section VIII.
## II Basic equations
We present the model based on that of BCF, supplemneted with asymmetric attachment kinetics as introduced by Schwoebel , and line diffusion following Ref. . A vicinal surface, whose mean interstep distance is $`\mathrm{}`$, is considered. On the terraces, the adatom concentration $`c_m`$ between steps $`m`$ and $`m+1`$ evolves according to:
$`_tc_m=D^2c_m+F,`$ (5)
where $`D`$ is the adatom diffusion constant, $`F`$ is an incoming flux of adatoms from a beam, and $`_t`$ denotes the time derivative. Once an atom is attached to the surface, it cannot detach from it (no desorption). We consider the widely used quasistatic limit where the concentration reaches a steady state regime on time scales much faster than that of step motion. We then have to solve Eq.(5) with the l.h.s. equal to zero. For implications due to non-quasisteady effects see Ref..
The excursion of the $`m^{th}`$ step about its straight configuration is denoted $`\zeta _m(x,t)`$, so that its position is $`m\mathrm{}+\zeta _m(x,t)+Vt`$, where $`V`$ is the mean step velocity. We consider the case where no step overhang is present, so that the function $`\zeta (x)`$ is univocal. On both sides ($`+`$ and $``$ designate the lower terrace and the upper one respectively) of step $`m`$, the normal diffusion flux is linearly related to departure from equilibrium with kinetic coefficients $`\nu _\pm `$:
$`D_nc_m|_+`$ $`=`$ $`\nu _+(c_mc_{eq})|_+`$ (6)
$`D_nc_{m1}|_{}`$ $`=`$ $`\nu _{}(c_{m1}c_{eq})|_{},`$ (7)
where $`c_{eq}`$ is the local equilibrium concentration, and $`_n`$ denotes the derivative in the direction which is normal to the step. More precisely $`_n𝐧.`$ where $`𝐧=(\zeta _x,1)/\sqrt{1+\zeta _x^2}`$ is the unit vector normal to the step, and $``$ is the two dimensional gradient operator: $`=(_x,_z)`$ where $`x`$ is the coordinate along an originally straight step, and $`z`$ the one orthogonal to it. The attachment lengths on both sides of the steps will be used later: $`d_+=D/\nu _+`$ and $`d_{}=D/\nu _{}`$. If $`c_{eq}^0`$ is the adatom concentration close to a straight step, the concentration for a curved step is given by:
$`c_{eq}`$ $`=`$ $`c_{eq}^0(1+\mathrm{\Gamma }\kappa _m),`$ (8)
where $`\mathrm{\Gamma }=\mathrm{\Omega }\stackrel{~}{\gamma }/k_BT`$ (the definition of $`\mathrm{\Gamma }`$ is slightly different from that of Ref. ) with $`\stackrel{~}{\gamma }`$ the step stiffness, and $`\kappa _m`$, the step curvature is given by:
$$\kappa _m=\frac{_{xx}\zeta _m}{[1+(_x\zeta _m)^2]^{3/2}}.$$
(9)
Here for simplicity we disregard step-step elastic interaction. We shall come back to this point in the discussion.
At the steps, mass conservation, in the limit where the adatom concentration is much smaller than that of the solid $`1/\mathrm{\Omega }`$, imposes:
$$V_n=\mathrm{\Omega }\left(D_nc_m|_+D_nc_{m1}|_{}\right)+a_s[D_L_s(\mathrm{\Gamma }\kappa _m)],$$
(10)
where $`a`$ is an atomic distance. Using Einstein’s relation, the macroscopic diffusion constant along steps is defined as $`D_L=D_{st}ac_{st}`$, where $`D_{st}`$ and $`c_{st}`$ are the diffusion constant and the concentration of mobile atoms along the step, respectively. This expression is in agreement with that of Mullins . One may object that $`c_{st}`$ is in fact not well defined along a step. We shall therefore use a more general expression derived from the Kubo formula : $`D_L=a^2/\tau _L`$ where $`\tau _L`$ is the characteristic time for detachment of an atom from a kink. Non-equilibrium effects related to line diffusion are not considered in this expression.
Two sources of nonlinearities can be identified. The first one is apparent in the boundary conditions (8, 10) because both the normal to the step and the curvature (see Eq.(9)) are nonlinear functions of the step profile. The second one originates the free boundary character (Stephan problem) is a hidden source of nonlinearity: the concentration field on a terrace –which is a nonlinear function of the position– depends on the step profile, leading thus to a nonlinear concentration field as a function of the step position.
In addition to elastic interactions (not included here), steps are coupled via adatom diffusion. Dynamics are nonlocal in space and time. With the help of an integral formulation of the model equations, we have made explicit this nonlocality in a previous work . The use of the quasistatic approximation suppresses delay effects, whereas spatial nonlocality persists.
## III Linear stability analysis
The linear stability analysis is the first step in any stability problem. Moreover it will allow us to prepare some preliminaries for the nonlinear analysis. Let us define the Fourier transform of the meander as:
$`\zeta _{\omega k\varphi }={\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle _{\mathrm{}}^+\mathrm{}}{\displaystyle _{\mathrm{}}^+\mathrm{}}\zeta _m(x,t)e^{i\omega tikxi\varphi m}𝑑x𝑑t,`$ (11)
where $`i\omega `$ is the pulsation of the perturbation of wavevector $`k`$ and phase shift between two neighboring steps, $`\varphi `$. The phase varies between $`0`$ and $`2\pi `$. Let us quote two special cases. The in-phase mode $`\varphi =0`$, corresponds to the case where all step meanders are identical, i.e. $`\zeta _m(x,t)=\zeta _m^{}(x,t)`$ for any $`m`$, $`m^{}`$. The out of phase mode $`\varphi =\pi `$ corresponds to the situation $`\zeta _m(x,t)=\zeta _{m+1}(x,t)`$.
The derivation of the full dispersion relation can be performed in this case along the same lines as in Ref.. We shall not repeat here the calculation, but give directly the result. The quantity $`i\omega `$ is complex, and let us discuss separately the real and imaginary parts. The real part of $`i\omega `$ takes the form
$`\mathrm{}e(i\omega )`$ $`=`$ $`\mathrm{\Omega }F{\displaystyle \frac{q}{𝒟}}\left({\displaystyle \frac{d_{}d_+}{\mathrm{}+d_{}+d_+}}\right)\left[(d_{}+d_+)\left(q\mathrm{}\mathrm{sinh}(q\mathrm{})\mathrm{cosh}(q\mathrm{})+\mathrm{cos}(\varphi )\right)+{\displaystyle \frac{\mathrm{}}{2}}q\mathrm{}\mathrm{sinh}(q\mathrm{})\right]`$ (13)
$`\mathrm{\Gamma }q^2\left[D_S{\displaystyle \frac{q}{𝒟}}\left(2\left(\mathrm{cosh}(q\mathrm{})\mathrm{cos}(\varphi )\right)+q(d_++d_{})\mathrm{sinh}(q\mathrm{})\right)+aD_Lq^2\right],`$
with $`q=|k|`$, and
$`𝒟`$ $`=`$ $`(d_++d_{})q\mathrm{cosh}(q\mathrm{})+(d_+d_{}q^2+1)\mathrm{sinh}(q\mathrm{}).`$ (14)
Both macroscopic diffusion constants (adatom tracer diffusion constant $`D`$ times coverage of mobile atoms) on the terraces $`D_S=D\mathrm{\Omega }c_{eq}^0`$ and $`D_L`$ along the steps enter this relation. The ”bare” (tracer) diffusion constant of adatoms on terraces does not appear alone.
A positive $`\mathrm{}e(i\omega )`$ is a signature of an instability. The straight step is unstable during growth provided that a normal ES effect is present ($`d_{}>d_+`$). Moreover, the most unstable mode is the in-phase mode $`\varphi =0`$. This remark will be exploited later.
The imaginary part of $`i\omega `$ describes propagative effects:
$$\mathrm{}m(i\omega )=\mathrm{\Omega }F\mathrm{sin}(\varphi )\frac{q}{𝒟}(\mathrm{}+d_++d_{}).$$
(15)
The origin of this term is quite transparent. In the limit of a straight step ($`q=0`$), we have $`\mathrm{}m(i\omega )=\mathrm{\Omega }F\mathrm{sin}(\varphi )`$, so that the perturbed solution takes the form (ignoring the real part of $`i\omega `$), $`\zeta _me^{in\varphi +it\mathrm{\Omega }F\mathrm{sin}(\varphi )}=e^{i\varphi [n+t(V_0/\mathrm{})\mathrm{sin}(\varphi )/\varphi ]}`$. Here we have introduced the step velocity of the uniform train, $`V_0=\mathrm{\Omega }F\mathrm{}`$. This means that in order to travel a distance $`n\mathrm{}`$, it takes for a perturbation a time given by $`(n\mathrm{}/V_0)(\varphi /\mathrm{sin}(\varphi ))`$. Since $`\varphi /\mathrm{sin}(\varphi )>1`$, that time is always longer than that needed for a uniform train ($`\varphi =0`$) to travel the same distance. In other words, all perturbations (except the in phase one) travel forward slower that the train velocity $`V_0`$. This means that perturbations are advected backwards in the reference frame with velocity $`V_0`$.
## IV Scaling analysis
Once the instability threshold is reached, any perturbation will amplify exponentially in the course of time, so that nonlinear effects can no longer be disregarded. As discussed in section II, the set of growth equations is highly nonlinear and nonlocal, so that only a ”brute force” numerical analysis would give a general answer. Our idea is to inspect the original equations and try to reduce legitimately the complexity. The key ingredient in our analysis is the identification of a small parameter.
### A Scaling of space and time variables
When inspecting the dispersion relation (13) for an in phase train, one realizes that the band of unstable wavenumbers extends from $`q=0`$ (actually this result is traced back to translational invariance; it corresponds to a global motion of the train) to a critical finite value $`q_c`$ (to be defined below). We shall assume that $`q\mathrm{}`$ remains small in comparison to one, and we come back in the discussion below to the validity of this assumption. In that case Eq.(13) takes a simpler form
$`Re[i\omega (q1,\varphi =0)]={\displaystyle \frac{\mathrm{\Omega }F\mathrm{}^2}{2}}{\displaystyle \frac{d_{}d_+}{\mathrm{}+d_++d_{}}}q^2(D_S\mathrm{}+D_La)\mathrm{\Gamma }q^4.`$ (16)
We have set here $`\mathrm{\Phi }=0`$, which is the exploitation of the fact that the in-phase mode is the most dangerous one. We consider later the situation where small deviations from the in-phase mode are taken into account. It is seen that the range of wavenumbers with positive $`i\omega `$ is given by
$`q_c=\left({\displaystyle \frac{\mathrm{\Omega }F\mathrm{}^2f_s}{\mathrm{\Gamma }(D_S\mathrm{}+D_La)}}\right)^{1/2},`$ (17)
where $`f_s=(d_{}d_+)/(\mathrm{}+d_++d_{})`$ is a parameter describing the Ehrlich-Schwoebel effect. The demand that the small wavenumber expansion make a sense is satisfied by requiring $`ϵ(q_c\mathrm{})^21`$, and this is precisely the definition of our small parameter
$`ϵ={\displaystyle \frac{\mathrm{\Omega }Ff_s\mathrm{}^4}{\mathrm{\Gamma }(D_S\mathrm{}+D_La)}}.`$ (18)
This guarantees the long wavelength regime. It is important to see from the very beginning whether this limit is realistic, or is it rather academic. Experimental data are available on vicinal surfaces of $`Cu(1,1,17)`$ which have recently revealed a meandering instability during step-flow growth. Their data entering the expression of $`ϵ`$ which are best known are $`\mathrm{\Omega }F=\mathrm{3\hspace{0.33em}10}^3s^1`$, $`\mathrm{}=21.7\AA `$. The step stiffness can be written as $`\stackrel{~}{\gamma }(k_BT/2a)\mathrm{exp}(E_k/k_BT)`$, $`E_k`$ being the kink energy. From a simple ”bond counting” argument, one can evaluate the adatom equilibrium concentration on a vicinal surface: $`\mathrm{\Omega }c_{eq}^0\mathrm{exp}(E_a/k_BT)`$ with $`E_a=3E_k`$. Using the result of Ref. from step fluctuations at equilibrium, we have $`E_k=0.13eV`$. The diffusion constant on terraces takes the form $`D=a^2\nu _0\mathrm{exp}(E_D/k_BT)`$, $`\nu _010^{13}s^1`$ is an intrinsic frequency, and $`E_D0.45eV`$ . With a lattice constant of $`2.55\mathrm{\AA }`$, we find: $`D_S\mathrm{}=1.4\times 10^{15}\mathrm{exp}(0.84\mathrm{e}\mathrm{V}/k_BT)`$.
Using the Kubo formula , one can evaluate the line diffusion constant: $`D_L=a^2/\tau _L`$. From experimental data : $`D_La=aD_{L0}\mathrm{exp}(E_L/k_BT)`$, with $`aD_{L0}=6.5\times 10^{18}\mathrm{\AA }^2s^1`$ and $`E_L=0.89eV`$. With these values, we find that $`D_S\mathrm{}/(D_La)10^2`$ in the experimental temperature range. This indicates that the stabilization of steps essentially occurs via line diffusion in this situation. In the one-sided limit ($`d_+=0`$ and $`d_{}\mathrm{}`$), and around $`300K`$ we find $`ϵ10^3`$, and $`\lambda _c=2\pi /q_c10^3`$ atomic distances. This result implies that a priori the longwavelenggth limit is appropriate.
The active modes in the instability are those for which $`q\mathrm{}q_c\mathrm{}ϵ^{1/2}`$, and therefore, lengthscales of interest are those for which $`x\mathrm{}ϵ^{1/2}`$. The characteristic time of the instability development is given by the growth rate of the most unstable mode:
$`t_m{\displaystyle \frac{2\pi \mathrm{}^4}{\mathrm{\Gamma }(D_S\mathrm{}+D_La)}}ϵ^2.`$ (19)
This is obtained as $`t_m2\pi /\mathrm{}e[i\omega (q=q_m,\varphi =0)]`$, $`q_m`$ being the wavevector of the most unstable mode, related to $`q_c`$ by $`q_m=q_c/\sqrt{2}`$. This relation provides the scaling of the time variable $`tϵ^2`$. Using the above data, we find that the instability typically develops after a growth of a thickness of the order of $`100`$ monolayers.
Before proceeding further, it is instructive to analyze briefly the asynchronized train. As pointed out in an earlier work , the ES effect not only induces a morphological instability of steps, but also leads to a “diffusive repulsion” between steps on a vicinal surface. This dynamical repulsion will force steps to evolve in-phase. The time needed for steps to organize in-phase in the unstable train is $`t_\varphi `$, defined as the synchronization time of the most unstable mode with wavevector $`q=q_m=q_c/\sqrt{2}`$ (i.e. the on e having the maximum growth rate):
$`{\displaystyle \frac{1}{t_\varphi }}_{\varphi \varphi }\mathrm{}e\left[\omega (q,\varphi )|_{q=q_m,\varphi =0}\right].`$ (20)
This time corresponds to the decay of a perturbation having phase shifts of order one. From linear dispersion relation (Eq.13):
$`{\displaystyle \frac{t_\varphi }{t_m}}ϵ\left({\displaystyle \frac{(\mathrm{}+d_++d_{})(D_S\mathrm{}+D_La)}{D_S\mathrm{}(d_{}+d_++\mathrm{}/2)+D_La(d_{}+d_+)}}\right),`$ (21)
where $`t_m`$ is the typical time for the instability to develop (Eq.(19)). In the one-sided limit, $`t_\varphi /t_mϵ`$. Thus, we expect the synchronization time to be shorter than the instability time, i.e. $`t_\varphi /t_m1`$ since $`ϵ1`$. This means that steps will be synchronized in the early stage of the instability. This justifies the fact to focus on small phases, or possibly a vanishing phase, as we shall assume later.
For small but finite $`\varphi `$ we have
$`\mathrm{}e[i\omega (q1,\varphi 1)]=`$ $``$ $`{\displaystyle \frac{\mathrm{\Omega }F}{2}}{\displaystyle \frac{d_{}d_+}{\mathrm{}+d_++d_{}}}\left((q\mathrm{})^2{\displaystyle \frac{d_++d_{}}{\mathrm{}+d_++d_{}}}\varphi ^2\right)`$ (22)
$``$ $`\left((D_S\mathrm{}+D_La)q^2+{\displaystyle \frac{D_S}{\mathrm{}+d_++d_{}}}\varphi ^2\right)\mathrm{\Gamma }q^2.`$ (23)
It is seen from the first term that for the phase shift to be relevant we must have $`\varphi q\mathrm{}ϵ^{1/2}`$. This implies that the conjugate variable $`m`$ (the step position along the vicinality) has the following scaling $`mϵ^{1/2}`$ (meaning that one needs to travel a distance of that order to detect phase modulations). In summary we have the following scaling in Fourier space
$`qϵ^{1/2},\omega ϵ^2,\varphi ϵ^{1/2},`$ (24)
and their corresponding conjugate variables in real space
$`xϵ^{1/2},tϵ^2,mϵ^{1/2}.`$ (25)
Beside the instability character, the problem involves progative effects which are related to the imaginary part of $`i\omega `$. Inspection of the imaginary part of the dispersion relation (15) in the long wavelength and small $`\varphi `$ limit, shows that $`\mathrm{}m(i\omega )ϵ^{3/2}`$. This defines a fast timescale $`\tau ϵ^{3/2}`$ related to propagative effects –we mean faster than the time scale associated with the instability $`ϵ^2`$. Since we shall mainly be interested by a synchronized train (in which case the imaginary part vanishes), we shall leave out this additional complication for the moment, and postpone this question to a forthcoming work.
### B Scaling of the meander
In order to determine the nonlinear evolution equation, following our previous work in the presence of desorption, we could expand all physical quantities (concentration, step position) in power series of the small parameter, the leading contribution would be of order $`ϵ^0`$, followed (in principle, and in a regular expansion) by $`ϵ^{1/2}`$, since this is the smallest power encountered above. This strategy worked out when desorption is present, but not in the present context. The reason is that the first contribution in the step profile turned out to be $`\zeta ϵ^{1/2}`$. This was viewed as an ansatz in our previous work . In this paper we provide an explanation of that fact on the basis of general considerations. Without having resort to an explicit derivation we shall show why this scaling is inherently linked with the non-desorption case.
For that purpose it is useful to identify two ’classes’ of adatoms (of course just in terms of a picture): Thermal adatoms of concentration $`c_T`$ detach from a step, diffuse on terraces and re-attach to a step. Mass transport associated to their motion induces relaxation towards equilibrium. Freshly landed adatoms of concentration $`c_F`$ have not yet been incorporated into a step. Their attachment result in the non-equilibrium driving of the steps. We can thus split the full set of equations (5-10) into two pieces by writing the model equations in the following equivalent form:
$`D^2c_T`$ $`=`$ $`0,`$ (26)
$`D^2c_F`$ $`=`$ $`F.`$ (27)
These fields obey the following boundary conditions at the steps:
$`D_nc_T`$ $`=`$ $`\pm \nu _\pm (c_Tc_{eq}),`$ (28)
$`D_nc_F`$ $`=`$ $`\pm \nu _\pm c_F,`$ (29)
where the index $`+`$ and $``$ refer to both sides of the steps. They are coupled only through mass conservation at the steps
$`V_n`$ $`=`$ $`v_F+v_T,`$ (30)
where the driving contribution $`v_F`$ is proportional to the incoming flux $`F`$:
$`v_F=D[_nc_{F_+}_nc_F_{}].`$ (31)
Indeed from the equations obeyed by $`c_F`$ (Eqs.(27), (29)) by making the transformation $`c_Fc_F/F`$ one sees that $`F`$ scales out from the equations, implying thus that $`c_F`$ must directly be proportional to $`F`$.
We can extract from $`c_F`$ the contribution of the uniform train, which leads to a velocity given by $`\mathrm{\Omega }F\mathrm{}`$, plus another contribution due to step modulations which must be compatible with conservation. $`v_F`$ is thus the sum of the mean step velocity and the divergence of a flux $`j`$ that describes how mass is unequally distributed between different steps, and different parts of each step:
$`v_F=\mathrm{\Omega }F(\mathrm{}j),`$ (32)
with $`j`$, according to what is stated above, independent of $`F`$.
The relaxational contribution $`v_T`$ is a thermal part and is obviously independent of $`F`$:
$`v_T`$ $`=`$ $`[D_nc_T]_{}^++a_s[D_L_s\mathrm{\Gamma }\kappa ].`$ (33)
Gradients of chemical potential $`\mu `$ are the driving force of the relaxational contribution. Without loss of generality, and as long as we deal with smooth and large scale perturbations, the thermal part of the normal velocity can be written with help of the Cahn-Hilliard equation:
$`v_T=[𝐌\mu ],`$ (34)
where $`𝐌`$ is the macroscopic mobility of the surface, and $`\mu =\mathrm{\Omega }\stackrel{~}{\gamma }\kappa `$ is the chemical potential. The step index $`m`$ is omitted in this section to simplify notations. Thus we shall from now on use the scalar mobility $`M`$ along $`x`$. The chemical potential is expressed as $`\mu =\mathrm{\Omega }\delta /\delta \zeta `$, where $``$ is the step free energy. Thus, if $`f(_x\zeta )`$ is the free energy density, we have:
$`\mu ={\displaystyle \frac{d}{dx}}[f^{}(_x\zeta )].`$ (35)
The evolution equation of the step meander (i.e. when the step mean velocity is subtracted) now reads:
$`_t\zeta =_x\left[\mathrm{\Omega }Fj+M_{xx}f^{}\right].`$ (36)
Recall that $`F`$ is proportional to $`ϵ`$ (Eq.(18)), so that we can set $`F=ϵ\overline{F}`$, where $`\overline{F}`$ is of order one. On the other hand $`x=Xϵ^{1/2}`$ (Eq.(25)), so that $`_{xx}=ϵ_{XX}`$. $`j`$, $`M`$ or $`f^{}`$ only depend on derivatives of $`\zeta `$, due to translational invariance. We write their argument symbolically as $`\{_x\}=\{ϵ^{1/2}_X\}`$ (which is taken to mean any derivative and any power). Equation (36) can be rewritten as:
$`_t\zeta =ϵ^{3/2}_X\left[\mathrm{\Omega }\overline{F}j\{ϵ^{1/2}_X\}+M\{ϵ^{1/2}_X\}_{XX}f^{}\{ϵ^{1/2}_X\}\right].`$ (37)
The central point lies in the fact that the small parameter $`ϵ`$ appears as a common factor, in the first term it stems from $`F`$ while in the second one it originates from the second spatial derivative.
In a regular expansion, close to the instability point we expect that the amplitude of modulation is vanishingly small when $`ϵ1`$. In reality, and this is the heart of the proof, due to the structure of the above equation, it will follow that no nonlinear term can enter the evolution equation, even if the amplitude were allowed to be of order one. There is even a stronger statement. Indeed even if $`\zeta =ϵ^\vartheta H`$ ($`H`$ is of order one), with $`\vartheta >1/2`$, we show below that any nonlinear term has a vanishing contribution. For that purpose we expand any function noted $`h`$ (which represents $`j`$….) in a Taylor series
$`h=h_0+h_1ϵ^{1/2+\vartheta }(_XH)+h_2ϵ^{1+2\vartheta }(_XH)^2+\mathrm{h}.\mathrm{o}.\mathrm{t}.,`$ (38)
where we have kept the smallest linear and nonlinear terms. For example a term like $`ϵ^{\vartheta +1}_{XX}Hϵ^{\vartheta +1/2}_XH`$. Although our conclusion can be made at this stage, let us be more explicit. Setting $`T=ϵ^2t`$, Eq.(36) now reads:
$`ϵ^{\vartheta +2}_TH=ϵ^{\vartheta +2}_X[j_1_XH+ϵ^{\vartheta +1/2}j_2(_XH)^2`$ (39)
$`+(M_0+ϵ^{\vartheta +1/2}M_1_XH)_{XX}(f_1^{}_XH+ϵ^{\vartheta +1/2}f_2^{}(_XH)^2)]+\mathrm{h}.\mathrm{o}.\mathrm{t}.,`$ (40)
Since $`\vartheta >1/2`$, we have $`ϵ^{\vartheta +1/2}0`$ as $`ϵ0`$. Therefore, nonlinear terms are irrelevant in Eq.(40), and the full equation reduces to a linear evolution equation:
$`_TH=\left(j_1_{XX}H+M_0f_1^{}_{XXXX}H\right).`$ (41)
For nonlinearities to be relevant in Eq.(40), we need $`ϵ^{\vartheta +1/2}O(1)`$, which is obtained when $`\vartheta =1/2`$. But then, the expansion performed in Eq.(38) is a priori not legitimate. Indeed, higher order terms become relevant: $`(_x\zeta )^nϵ^{n(\vartheta +1/2)}O(1)`$ when $`\vartheta =1/2`$ for any integer $`n`$. We therefore expect a highly nonlinear evolution equation, as will be shown explicitly in the next section.
How concentration scales with $`ϵ`$ can also be found using the decomposition of the concentration. From equations (27) and (29), we have $`c_FFϵ`$. From Eq.(26), (28), and (8), we find that
$`c_Tc_{eq}^0c_{eq}c_{eq}^0c_{eq}^0\mathrm{\Gamma }\kappa ϵ^{1/2}.`$ (42)
Thus, $`u=\mathrm{\Omega }c_F+c_Tc_{eq}^0ϵ^{1/2}`$, and the concentration will be written in the following form:
$`u(x,t)`$ $`=`$ $`ϵ^{1/2}U(x,t),`$ (43)
with $`U(x,t)O(1)`$. Similarly, the meander will be written as:
$`\zeta (x,t)`$ $`=`$ $`ϵ^{1/2}H(x,t),`$ (44)
where $`H(x,t)O(1)`$.
It is important to show why in the presence of desorption the expansion is regular, leading to the KS equation (). With desorption the evolution equation can no longer be written in the form of a conservation law (2). Nevertheless, the above mentioned decomposition still holds in a slightly different form: instead of being proportional to $`F`$, the driving part is proportional to $`FF_{eq}`$, where $`F_{eq}`$ is the incoming flux at equilibrium that counterbalances ambiant desorption ($`F_{eq}=c_{eq}^0/\tau `$, $`\tau `$ being the characteristic residence time before desorption on a terrace). Hence, instead of Eq.(36), we have:
$`_t\zeta =(F_{eq}F)g+N_{xx}f^{},`$ (45)
where $`g`$ and $`N`$ are functions of the derivatives of $`\zeta `$. The second term of the r.h.s. is now the one expected for non-conserved relaxation to equilibrium: it is directly proportional to the chemical potential variations (See model A in Ref. , or ). Linearizing this equation, we have to take $`g\stackrel{~}{g}_{xx}\zeta `$ since the first linear term proportional to $`_x\zeta `$ is a propagative term, not contributing to stabilization or destabilization (moreover, this term, not invariant under the $`xx`$ symmetry, is not allowed). We then find:
$`_t\zeta =[(F_{eq}F)\stackrel{~}{g}+N_0f_1^{}]_{xx}\zeta .`$ (46)
The prefactor of $`_{xx}\zeta `$ is the effective stiffness of the step. An instability is signaled by a negative sign of that prefactor. This happens when $`F>F_c=F_{eq}+N_0f_1^{}/\stackrel{~}{g}`$. The small parameter (that is the distance from the instability threshold) is now $`ϵ^{}FF_c`$. Moreover it was found in Ref. that $`xϵ^{1/2}`$ and $`tϵ^2`$. Defining as in the conserved case $`X=ϵ^{1/2}x`$ and $`T=ϵ^2t`$, and $`\zeta =ϵ^\vartheta ^{}H`$, with $`X,T,HO(1)`$, Eq.(45) is now expanded for $`FF_c`$ as:
$`ϵ^{\vartheta ^{}+2}_TH=ϵ^{\vartheta ^{}+2}\stackrel{~}{g}_{XX}H+ϵ^{2\vartheta ^{}+1}(FF_{eq})(_XH)^2g_2`$ (47)
$`+ϵ^{2\vartheta ^{}+2}\left[N_0f_2^{}_{XX}(_XH)^2+N_1f_1^{}_XH_{XXX}H\right]+\mathrm{h}.\mathrm{o}.\mathrm{t}..`$ (48)
As before, we use $`\vartheta ^{}>1/2`$ so that expansion (38) makes a sense. It is seen from this equation that the leading nonlinear term is $`(FF_{eq})ϵ^{2\vartheta ^{}+1}(_XH)^2`$. It counterbalances the linear term $`ϵ^{\vartheta ^{}+2}\stackrel{~}{g}_{XX}H`$ provided that $`\vartheta ^{}=1`$. The nonlinear term here, is that of the Kardar-Parisi-Zhang and Kuramoto-Sivashinsky equations (see Eq.(1) and (3)). It could not be present in the conserved case because it cannot be written as the divergence of a flux. Moreover, it is non-variational, and thus it must vanish at equilibrium, as can bee seen from its prefactor $`FF_{eq}`$.
It must be emphasized that the decomposition into an equilibrium an a nonequilibrium part holds in the present problem, but is not a general property. This does not have to be the case out of equilibrium in general (an example is that of step flow or sublimation in the presence of electromigration, such a decomposition between relaxation and driving parts has not been made possible).
## V One-sided synchronized steps
### A Multiscale analysis
In addition to synchronization, we first assume for simplicity a one-sided limit (steps advance only thanks to atoms from the terrace which is ahead), formally defined as $`d_+=0`$, and $`d_{}+\mathrm{}`$. In this limit Eq.(7) reduces to $`c_+=c_{eq}`$ (which is the Gibbs-Thomson condition) and $`c_{}/n=0`$ (atoms do not descend the steps).
As we have shown in the last section the meander $`\zeta ϵ^{1/2}`$, while the concentration field $`uϵ^{1/2}`$, we find it convenient to set $`\zeta =ϵ^{1/2}H`$ and $`u=ϵ^{1/2}U`$, with $`H`$ and $`U`$ being quantities of order one. Under the assumption that these quantities are analytic functions of $`ϵ^{1/2}`$, we seek solutions of the form
$`U`$ $`=`$ $`U^{(0)}+ϵ^{1/2}U^{(1/2)}+ϵU^{(1)}+ϵ^{3/2}U^{(3/2)}+\mathrm{},`$ (49)
$`H`$ $`=`$ $`H^{(0)}+ϵ^{1/2}H^{(1/2)}+ϵH^{(1)}+ϵ^{3/2}H^{(3/2)}+\mathrm{}.`$ (50)
In order to make explicit the $`ϵ`$ dependence, and to deal with quantities of order 1, and according to (25), we set:
$`x=ϵ^{1/2}X,t=ϵ^2T`$ (51)
It is convenient to rescale space by $`\mathrm{}`$ and time by $`D/\mathrm{}^2`$. Performing the variable change: $`𝒵=z\zeta _m(x,t)`$, mass conservation (5) on terraces reads:
$`0=\rho ^2_{𝒵𝒵}U+ϵ^{1/2}\left(\eta 2_XH_{X𝒵}U_{XX}H_𝒵U\right)+ϵ_{XX}U,`$ (52)
where $`\rho =[1+(_XH)^2]^{1/2}`$, an $`\eta =ϵD/\mathrm{\Omega }F\mathrm{}^2`$. At the steps, the Gibbs-Thomson relation at $`𝒵=0`$, and a zero-flux condition at $`𝒵=1`$, takes the form:
$`U|_{𝒵=0}`$ $`=`$ $`𝒦,`$ (53)
$`\rho ^2_𝒵U|_{𝒵=1}`$ $`=`$ $`ϵ^{1/2}_XH_XU|_{𝒵=1},`$ (54)
where $`𝒦=\mathrm{\Omega }c_{eq}^0\mathrm{\Gamma }_{XX}H/\rho ^3`$.
Mass conservation at the step (Eq.(10)) yields
$`V+ϵ^{3/2}_TH=\rho ^2_𝒵U|_{𝒵=0}ϵ^{1/2}_XH_XU|_{𝒵=0}ϵ^{3/2}_X\left({\displaystyle \frac{\beta }{\rho }}_X𝒦\right),`$ (55)
where $`\beta =D_La/D_S\mathrm{}`$. The strategy is now to solve equations (52-55) in successively higher orders in $`ϵ`$.
order 0
To this order, Eq.(52) reads:
$$_{𝒵𝒵}U^{(0)}=0,$$
(56)
which is solved by $`U^{(0)}=A^{(0)}𝒵+B^{(0)}`$. Equations (53) and (54) provide two conditions from which we get $`A^{(0)}=0`$ and $`B^{(0)}=𝒦^{(0)}`$. No contribution to step velocity is found to $`0^{th}`$ order. That is to say this order corresponds to the equilibrium case.
order 1/2
From (52) we find that $`U^{(1/2)}`$ obeys an inhomogeneous equation on terraces:
$$\rho ^{(0)2}_{𝒵𝒵}U^{(1/2)}=\eta ,$$
(57)
whose general solution takes the form:
$$U^{(1/2)}=\frac{𝒵^2}{2\rho ^{(0)2}}\eta +A^{(1/2)}𝒵+B^{(1/2)}.$$
(58)
From boundary conditions at the steps (53) and (54):
$`U^{(1/2)}`$ $`=`$ $`B^{(1/2)}=𝒦^{(1/2)}`$ (59)
$`\rho ^{(1/2)2}_𝒵U^{(1/2)}|_{𝒵=1}`$ $`=`$ $`_XH^{(0)}_XB^{(1/2)}.`$ (60)
Integration constants are found to be:
$`A^{(1/2)}`$ $`=`$ $`(\eta _XH^{(0)}_X𝒦^{(0)})/\rho ^{(0)^2},`$ (61)
$`B^{(1/2)}`$ $`=`$ $`𝒦^{(1/2)}.`$ (62)
Mass conservation at the step (55) determines the mean step velocity. Going back to physical variables, we find the expected result: $`V=\mathrm{\Omega }F\mathrm{}`$.
order 1
To this order, $`U^{(1)}`$ obeys:
$`_{𝒵𝒵}U^{(1)}`$ $`=`$ $`{\displaystyle \frac{1}{\rho ^{(0)^2}}}[_{XX}H^{(0)}_𝒵U^{(1/2)}+2_XH^{(0)}_{X𝒵}^2U^{(1/2)}`$ (64)
$`2_XH^{(0)}_XH^{(1/2)}_{𝒵𝒵}U^{(1/2)}_{XX}U^{(0)}]`$
$`=`$ $`a+b𝒵,`$ (65)
whose general solution takes the form:
$$U^{(1)}=\frac{b}{6}𝒵^3+\frac{a}{2}𝒵^2+A^{(1)}𝒵+B^{(1)}.$$
(66)
Once again, integration factors $`A^{(1)}`$ and $`B^{(1)}`$ are found from boundary conditions (53) and (54):
$`U^{(1)}|_{𝒵=0}`$ $`=`$ $`B^{(1)}=𝒦^{(1)},`$ (67)
$`_𝒵U^{(1)}|_{𝒵=1}`$ $`=`$ $`{\displaystyle \frac{b}{2}}+a+A^{(1)}`$ (68)
$`=`$ $`{\displaystyle \frac{1}{\rho ^{(0)^2}}}[_XH^{(0)}_XU^{(1/2)}`$ (69)
$`+`$ $`_XH^{(1/2)}_XU^{(0)}`$ (70)
$``$ $`2_XH^{(0)}_XH^{(1/2)}_𝒵U^{(1/2)}]|_{𝒵=1}.`$ (71)
Finally, mass conservation (55) leads to the sought after evolution equation for $`H^{(0)}`$:
$`_TH^{(0)}`$ $`=`$ $`_𝒵U^{(1)}\rho ^{(0)2}_XU^{(1/2)}_XH^{(0)}_X\left({\displaystyle \frac{\beta }{\rho ^{(0)}}}_X𝒦^{(0)}\right).`$ (72)
Upon substitution of the expressions of $`U^{(1)}`$ and $`U^{(1/2)}`$, one realizes that terms containing $`H^{(1/2)}`$ cancel exactly in this expression, leading to a closed form for the evolution equation for $`H^{(0)}`$:
$`_TH^{(0)}`$ $`=_X\left[\eta {\displaystyle \frac{_XH^{(0)}}{2\rho ^{(0)^2}}}+\left(1+\beta \rho ^{(0)}\right){\displaystyle \frac{_X𝒦^{(0)}}{\rho ^{(0)^2}}}\right].`$ (73)
Going back to physical variables, we obtain:
$$_t\zeta =_x\left[\frac{\mathrm{\Omega }F\mathrm{}^2}{2}\frac{_x\zeta }{(1+(_x\zeta )^2)}\left(D_S\mathrm{}+D_La(1+(_x\zeta )^2)^{1/2}\right)\frac{_x(\mathrm{\Gamma }\kappa )}{(1+(_x\zeta )^2)}\right].$$
(74)
Besides the term proportional to $`D_L`$ (line diffusion constant), this is the equation derived in Ref. on which we have given a brief account.
Introducing the step macroscopic mobility $``$, and the chemical potential $`\mu =k_BT\mathrm{\Gamma }\kappa `$, the evolution equation can be rewritten in a more compact and enlightening form:
$$_t\zeta =_x\left[\frac{\mathrm{\Omega }F\mathrm{}_{}^2}{2}_x\zeta _s\mu \right],$$
(75)
where $`s`$ is the arclength along the steps, and $`\mathrm{}_{}=\mathrm{}/[1+(_x\zeta )^2]^{1/2}`$ is the distance between two neighboring steps measured along their normal (see Fig.2 for geometrical definitions.). The effective step mobility reads:
$`={\displaystyle \frac{D_S\mathrm{}_{}+D_La}{k_BT}}.`$ (76)
The expected decomposition of step velocity (see section IV B) is clearly seen here. The first term on the r.h.s. of Eq.(75) is the driving part. To this term a simple geometrical meaning can be assigned (see Appendix A). The second term is the relaxation part with a mobility depending on the local step orientation. Note that the present mobility $``$ and the one introduced in section IV B, noted $`M`$, differ by the scale factor $`[1+(_x\zeta )^2]^{1/2}`$ which relates the arc-length $`s`$ to the Cartesian coordinate $`x`$.
### B Numerical solution
In Ref , numerical solution of Eq.(74) (without line diffusion) was performed using a simple Euler scheme, and it was found that: (i) A cellular structure takes place, the wavelength of which (the most unstable one) is fixed at the initial stage of the instability and no coarsening is seen. (ii) The amplitude grows like $`t^{1/2}`$. (iii) The shape of the cells is similar to the inverse error function, that is to say it develops a spike-like morphology. (iv) The meander is symmetric with respect to the transformation $`\zeta \zeta `$. It has been realized meanwhile that, though all these qualitative features were correct, the spikes are the result of a numerical deficiency in the original code. In order to cure this problem we have to resort to a special numerical treatment. Different successful attempts have been made but we shall here describe only the most robust numerical treatment.
We use a powerful geometrical representation of the meander , in terms of the arclength $`s`$ and the angle $`\theta `$, oriented counterclockwise, between the normal and a given fixed direction (the z-axis direction) (see Fig.2). $`\theta `$ is related to $`\zeta `$ via: $`\mathrm{tan}(\theta )=_x\zeta `$ and the curvature simply reads: $`\kappa =_s\theta `$.
Simple differential geometry provides us with the evolution equation for $`\theta `$, as a function of tangential and normal velocities, $`v_t`$ and $`v_n`$:
$$\frac{\theta }{t}=v_t\kappa \frac{v_n}{s}.$$
(77)
Physics is invariant under a change of definition of the arclength $`s`$. This allows an arbitrary time-dependent re-parameterization of the curve. This ”gauge” can be seen as an additional tangential velocity, with no physical relevance. A particular choice that is convenient here is the one that keeps the relative arclength $`s/L`$ constant in the course of time, where $`L`$ is the total length of the curve. This will ensure that the discretization points remain equally spaced along the curve. The tangential velocity reads:
$$v_t=\frac{s}{L}_0^L\kappa v_n𝑑s^{}_0^s\kappa v_n𝑑s^{}.$$
(78)
The evolution equation of the meander (74) allows one to write the step normal velocity:
$$v_n=_s\left[\mathrm{cos}(\theta )\mathrm{sin}(\theta )+\left(\frac{\beta +\mathrm{cos}(\theta )}{\beta +1}\right)_s\kappa \right],$$
(79)
where time $`t`$ is rescaled by $`4\mathrm{}^4/ϵ^2\mathrm{\Gamma }(lD_S+aD_L)`$ and spatial variables $`x,\zeta ,`$ by $`\sqrt{2}\mathrm{}/\sqrt{ϵ}`$, so that only one parameter survives: $`\beta =D_La/D_S\mathrm{}`$.
Derivatives along the arclength $`s`$ are evaluated using a centered finite difference method. We use a backward differentiation scheme with variable step for time integration. This ”solver” enjoys rather good precision and L-stability (that is to say it is unconditionally stable and optimally attenuates high-frequency (i.e. noise) components of the solution) that makes it well fitted to our specific problem.
Our present simulations show qualitatively similar behavior as that found in Ref. (see Fig.6). A major difference is revealed however: instead of spikes, the cellular structure exhibits a plateau in the extrema regions , as shown on Fig.3. (Here we define a plateau as a region of finite slope, as opposed to regions where the slope diverges with time. Hence, to some scale, plateaus are curved, and this curvature does not tend to zero). The width of a ”plateau” reaches a constant value after a transient regime.
### C Analytical study
We show here that the above numerical results can be accounted for using simple analytical arguments. We give here only the main results, whereas details are relegated into Appendix B. The central assumption is a decomposition in two types of regions, where two different ansatz are used. In the large slope regions, a multiplicative variable separation is used:
$`\zeta _s(x,t)=A(t)g(x),`$ (80)
while and additive variable separation is performed for the plateau regions:
$`\zeta _p(x,t)=B(t)+h(x).`$ (81)
An additional constraint coming from mass conservation allows to determine quantitatively the asymptotic behavior by matching these two solutions. The amplitude of the meander is found to behave as:
$`\zeta _{max}\zeta _{min}=2a_0t^{1/2},`$ (82)
where $`a_0`$ is calculated in Appendix B. The rescaled meander $`\zeta (x,t)/t^{1/2}`$ converges to a well defined profile, and looks as if plateaus were formed in the extrema regions. The width of a plateau is defined as $`\lambda _0/2`$. We find:
$`\lambda _0=\lambda _cI(\stackrel{~}{\beta }),`$ (83)
where $`\stackrel{~}{\beta }=D_S\mathrm{}/(D_S\mathrm{}+D_La)`$, $`\lambda _c=2\pi /q_c`$ is the largest stable wavelength from linear analysis, and $`I`$ is a function given in Appendix B. Since $`I(\stackrel{~}{\beta })`$ decreases monotonously from $`I(0)=1`$ to $`I(1)0.54`$, the plateau size increases as line diffusion is increased, and is always smaller than $`\lambda _c/2`$. Good quantitative agreement between the numerical solution of Eq.(74) and these analytical predictions is found (see Fig.4).
## VI Front-Back symmetry breaking
The expansion performed in section V A can be pushed to next order following the same strategy. We shall here merely give the result and details can be found in Ref. . Instead of a closed equation for $`H^{(0)}`$, here two coupled dynamical equations for $`H^{(0)}`$ and $`H^{(1/2)}`$ are obtained. Going back to the physical quantity $`\zeta =ϵ^{1/2}(H^{(0)}+ϵ^{1/2}H^{(1/2)})`$, the coupled equations can be recast into a single equation for $`\zeta `$:
$`_t\zeta `$ $`=`$ $`{\displaystyle \frac{}{x}}\left[{\displaystyle \frac{\mathrm{\Omega }F\mathrm{}_{}^2}{2}}_x\zeta \left(1{\displaystyle \frac{\kappa \mathrm{}}{3}}\left({\displaystyle \frac{\mathrm{}}{\mathrm{}_{}}}+{\displaystyle \frac{2\mathrm{}_{}}{\mathrm{}}}\right)\right)^{(1/2)}{\displaystyle \frac{\mu }{s}}\right],`$ (84)
where the macroscopic mobility of the step reads:
$$^{(1/2)}=\frac{D_S\mathrm{}_{}+D_La}{k_BT}\frac{D_S\mathrm{}^2\kappa }{2k_BT}.$$
(85)
Hence, to this order, correction to Eq.(74) are proportional to step curvature.
As before, a geometrical formulation with rescaled time and space variables, is used. We now have two parameters $`\beta `$ and $`ϵ`$ in the normal velocity:
$`v_n`$ $`=`$ $`_s\left[\mathrm{cos}(\theta )\mathrm{sin}(\theta )+\left({\displaystyle \frac{\beta +\mathrm{cos}(\theta )}{\beta +1}}\right)_s\kappa \sqrt{ϵ}\kappa \left({\displaystyle \frac{2}{3}}\left(\mathrm{cos}(2\theta )+2\right)\mathrm{sin}(\theta )+_s\kappa \right)\right].`$ (86)
The same qualitative features as for Eq.(74) are observed: the wavelength is fixed at early stages by the one corresponding to the fastest growing mode, and the step roughness increases with time as $`t^{1/2}`$ (see Fig.6). The interesting fact is the symmetry of the shape: the cells do not enjoy the up-down symmetry, $`\zeta \zeta `$. Clearly, Eq.(74) is invariant under the up-down symmetry $`\zeta \zeta `$. The symmetry breaking originates from the new terms as can be seen by performing the transformation $`\zeta \zeta `$ on Eq.(84). As shown in Appendix B, these terms affect the relative sizes of the back and front plateaus, but the $`t^{1/2}`$ scaling law for the roughness amplitude seems to persist as a robust feature.
The results obtained in this section are in complete agreement with full simulations based on a solid-on-solid model in Ref. . Hence we have succeeded in extracting the relevant dynamics of step meander by means of the multiscale analysis.
## VII Two-sided steps in phase
In the two-sided regime (i.e. $`d_+`$ and $`d_{}`$ finite), a similar multi-scale analysis can be performed for in-phase steps. We find:
$`_t\zeta =_x\left[{\displaystyle \frac{\mathrm{\Omega }F}{2}}_x\zeta {\displaystyle \frac{\mathrm{}_{}^2(d_{}d_+)}{d_++d_{}+\mathrm{}_{}}}+{\displaystyle \frac{1}{(1+(_x\zeta )^2)^{1/2}}}\left(D_La+D_S{\displaystyle \frac{\mathrm{}^2+\mathrm{}_{}(d_++d_{})}{d_++d_{}+\mathrm{}_{}}}\right)_x(\mathrm{\Gamma }\kappa )\right].`$ (87)
Although this equation looks more complicated, the meander evolution is qualitatively similar to that found in the one-sided case (which is recovered by taking the limit $`d_{}\mathrm{}`$ in Eq.(87)). Indeed, plateau formation and power law behavior of the roughness (with the same exponent $`t^{1/2}`$) are also found in the two-sided case.
More details on step behavior, such as the plateau size, can be gained from the analytical investigation of Eq.(87), as shown in Appendix B. In the pure line diffusion regime $`D_LaD_S\mathrm{}`$, we have:
$`\lambda _0=\lambda _cI(\stackrel{~}{\delta }),`$ (88)
where $`\lambda _0/2`$ is the plateau size, $`\stackrel{~}{\delta }=\mathrm{}/(\mathrm{}+d_++d_{})`$, and I is the same function as in Eq.(83). In the pure terrace diffusion case $`D_LaD_S\mathrm{}`$, we find:
$`\lambda _0=\lambda _cI(1\stackrel{~}{\delta }).`$ (89)
These result are in good agreement with numerical solution of Eq.(87) (see Fig.7).
## VIII Discussion and summary
Starting from the BCF model, we have extracted a nonlinear evolution equation for the step meander. This equation is highly nonlinear, and thus, could not be expected from traditional phenomenological approaches, where linear terms are simply supplemented with an additive nonlinear term, as in the case of KS or KPZ equations.
A central result of the present study is that the late time power law behavior of the amplitude of the meander $`wt^{1/2}`$ is a robust feature, regardless of the details of the evolution equation (one-sided, two-sided, line diffusion….). It is an important task for future investigations to see whether this property could be derived directly from the basic BCF equations. A second interesting point which is worth of mention is that higher order terms destroy the up-down symmetry, but the power law $`wt^{1/2}`$ remains unaffected. The same conclusion follows from full lattice gas simulations as briefly reported in Ref.. A natural question arises: could the amplitude temporal increase continue to evolve without bound in all circumstances until the surface breaks up into a lamellar-like pattern, or is there a physical mechanism, not accounted for here, leading to saturation of the amplitude? Experimental observation of this instability seems to show such a saturation for the case of Cu(1,1,17), while experiments on Si(001) does not reveal a hint towards a saturation. Possible candidates for amplitude saturations are (i) strong anisotropy, (ii) elastic step interactions. We hope to report along these lines in the near future.
It is worth pointing out that the use of equilibrium formula to evaluate the stabilizing line diffusion effect could be criticized, since densities of kinks and of mobile atoms along steps depend on growth conditions . With regards to line mobility, our analysis allows extraction of the geometry dependence of the mobility for large meander amplitude. This treatment should serve as a basis for the nonlinear study of relaxation towards equilibrium (i.e. thermal smoothening) of large perturbations on vicinal surfaces .
Perhaps one of the most striking result is the manifestation of rather ’stringent’ plateaus, which are likely linked to the non-standard character of the evolution equation. It is interesting to note that the plateaus are a feature of a continuum theory, a finding which is to be compared to a long standing problem in the context of ES-induced mound formation . In all previous studies, mound plateaus were indeed considered as a signature of the breakdown of continuum theories. We have shown here, in contrast, that a single equation in the continuum limit can produce such plateaus, without having resort to specific ingredients in the angular region. It is not yet clear what kind of equation in the continuum limit would describe these dynamics for mound formation. Is it similar or not to the one encountered here? These questions constitute an important line for future inquiries.
## A Geometrical origin of the destabilizing term
In Eq.(74), the relaxation term is interpreted as a Cahn-Hilliard contribution. We present here a derivation of the destabilizing term from geometrical considerations in the one-sided limit.
Let us consider a curved part of the step as shown in Fig.8. In the one-sided model, step motion results from incorporation of adatoms from the terrace ahead of it. Mass conservation for an element of terrace surface $`\mathrm{\Delta }S`$, hatched in Fig.8, reads:
$`v\mathrm{\Delta }x=\mathrm{\Omega }F\mathrm{\Delta }S+j_{}(x)j_{}(x+\mathrm{\Delta }x)`$ (A1)
where $`v`$ is the step velocity along the $`z`$ axis, and $`\mathrm{\Delta }x`$ the extent of the step element $`CC^{}`$ along the $`x`$ axis. The number of atoms entering the step is $`v_n\mathrm{\Delta }s=v\mathrm{\Delta }x`$. $`j_{}(x)`$ is the total flux accross the $`BC`$ segment in Fig.8. $`\mathrm{\Delta }S`$ is written as:
$`\mathrm{\Delta }S\mathrm{}\mathrm{\Delta }x𝒜(x)+𝒜(x+\mathrm{\Delta }x),`$ (A2)
where $`v`$ is the step velocity along the $`z`$ axis, and $`𝒜(x)`$, the area of the triangle $`ABC`$ on Fig.8, is a function of $`_x\zeta `$:
$`𝒜(x)={\displaystyle \frac{\mathrm{}^2}{2}}\mathrm{cos}(\theta )\mathrm{sin}(\theta )={\displaystyle \frac{\mathrm{}^2}{2}}{\displaystyle \frac{_x\zeta }{1+(_x\zeta )^2}},`$ (A3)
where $`\theta `$ is the angle between the $`z`$ axis and the normal to the step. In the long wavelength limit, the local geometry of the terrace is completely described by $`\mathrm{}_{}`$, the length of $`BC`$ on Fig.8, $`\kappa `$, the step curvature, and their derivatives with repect to the arclength $`s`$ along the steps. Since the flux $`j_{}`$ can only come from a variation of the local geometry along $`s`$, we have, at most
$`j_{}_s\mathrm{}_{}_{xx}\zeta 𝒜_x\zeta ,`$ (A4)
which shows that terms containing $`j_{}`$ can be neglected to leading order in Eq.(A1) . Combining Eq.(A1), Eq.(A2), and (A3), and letting $`\mathrm{\Delta }x`$ going to zero, we find:
$`v=\mathrm{\Omega }F\mathrm{}_x\left({\displaystyle \frac{\mathrm{\Omega }F\mathrm{}^2}{2}}{\displaystyle \frac{_x\zeta }{1+(_x\zeta )^2}}\right).`$ (A5)
Once the mean step velocity $`V=\mathrm{\Omega }F\mathrm{}`$ is subtracted, we recover the first term of Eq.(74).
## B Late time behavior
In this appendix we derive analytically the main results obtained numerically. Despite the highly nonlinear character of the evolution equation, we show here that some simple ansatz allows to describe the asymptotic regime with good accuracy.
### 1 Large slope regions and extrema regions
We found the conserved evolution equation of the step meander:
$`_t\zeta =_xj[\zeta ],`$ (B1)
with the mass flux (see Eq.(87)):
$`j[\zeta ]`$ $`=`$ $`{\displaystyle \frac{1}{(1+(_x\zeta )^2)^{1/2}}}[{\displaystyle \frac{\alpha _x\zeta }{\delta +(1+(_x\zeta )^2)^{1/2}}}`$ (B2)
$`+`$ $`(D_La+D_S\mathrm{}{\displaystyle \frac{1+\delta (1+(_x\zeta )^2)^{1/2}}{\delta +(1+(_x\zeta )^2)^{1/2}}})_{xx}\left({\displaystyle \frac{\mathrm{\Gamma }_x\zeta }{(1+(_x\zeta )^2)^{1/2}}}\right)],`$ (B3)
where $`\delta =l/(d_++d_{})`$, and $`\alpha =\mathrm{\Omega }F\mathrm{}^2(d_{}d_+)/2(d_++d_{})`$. For the large slope region we make use of the variable separation:
$`\zeta _s(x,t)=A(t)g(x),`$ (B4)
where $`A1`$, and $`_xg0`$ for any value of $`x`$. Substituting in Eq.(B3), one finds that the destabilizing term (proportional to $`\alpha `$) dominates and:
$`AA^{}=\alpha {\displaystyle \frac{g^{\prime \prime }}{gg^2}}=C,`$ (B5)
where $`C`$ is a constant, and the prime stands for the derivative. The late time solution of these equations reads:
$`A`$ $`=`$ $`(2Ct)^{1/2}`$ (B6)
$`g`$ $`=`$ $`(2\alpha /C)^{1/2}\mathrm{erf}^1(4x/\lambda _s).`$ (B7)
$`\lambda _s`$ being a constant, and erf$`(x)`$ the error function. Inserting these expressions in Eq.(B4), we find that the meander does not depend on $`C`$:
$`\zeta _s(x,t)=2(\alpha t)^{1/2}\mathrm{erf}^1(4x/\lambda _s),`$ (B8)
This solution describes regions of large slopes, but is not expected to accurately describe the shape around the extrema of $`\zeta `$ where the slope $`_x\zeta `$ approaches zero. In those regions of width $`\lambda _0/2`$ and of meander amplitude of the order of $`h_0`$ (see Fig.9), global mass conservation implies:
$`2j_0={\displaystyle \frac{\lambda _0}{2}}_th_0,`$ (B9)
where $`j_0`$ is the mass flux coming from the large slope regions. From Eq.(B3), we have:
$`j_0{\displaystyle \frac{\alpha }{_x\zeta _s(x_0,t)}},`$ (B10)
where $`x_0=(\lambda _m\lambda _0)/4`$ is the abscissa of the crossover point between high slope and extrema regions, and the period of the meander $`\lambda _m`$ is that of the most unstable mode obtained from linear analysis. Using $`h_0=\zeta _s(x_0,t)`$ and Eq.(B8):
$`{\displaystyle \frac{\lambda _0}{4}}={\displaystyle \frac{\lambda _m}{4}}{\displaystyle \frac{\lambda _s}{4}}\mathrm{erf}\left[{\displaystyle \frac{h_0}{2(\alpha t)^{1/2}}}\right],`$ (B11)
so that Eq.(B9) now reads:
$`{\displaystyle \frac{1}{t^{1/2}}}{\displaystyle \frac{\alpha ^{1/2}\lambda _s}{4\sqrt{\pi }}}\mathrm{exp}\left[\left({\displaystyle \frac{h_0}{2(\alpha t)^{1/2}}}\right)^2\right]`$ (B12)
$`=_th_0\left[{\displaystyle \frac{\lambda _m}{4}}{\displaystyle \frac{\lambda _s}{4}}\mathrm{erf}\left[{\displaystyle \frac{h_0}{2(\alpha t)^{1/2}}}\right]\right],`$ (B13)
This equation has the trivial solution: $`h_0=a_0t^{1/2}`$. Using this solution and Eq.(B11), $`\lambda _0`$ is seen not to depend on time. In the extrema region, we therefore look for solutions of the form:
$`\zeta _p(x,t)=B_\pm (t)+h(x).`$ (B14)
with $`B_\pm (t)=\pm a_0t^{1/2}`$, where the plus and minus signs refer to the maxima and the minima regions respectively. Upon substitution in the evolution equation (B3), we find that the problem amounts to finding the stationary solutions $`_xj[h(x)]=0`$. Looking for solutions with left-right symmetry $`xx`$ we finally have to solve
$`j[h(x)]=0.`$ (B15)
This will be exploited in the next section.
The parameters $`\lambda _s`$, $`\lambda _0`$ and $`a_0`$ are not independent. Using Eq.(B9) and Eq.(B11), we have two relations, so that $`\lambda _s`$ and $`a_0`$ can be determined as a function of $`\lambda _0`$. From Eq.(B13), we get an implicit equation for $`a_0`$:
$`{\displaystyle \frac{\lambda _m}{\lambda _0}}1=\sqrt{\pi }\left({\displaystyle \frac{a_0}{2\alpha ^{1/2}}}\right)\mathrm{exp}\left[\left({\displaystyle \frac{a_0}{2\alpha ^{1/2}}}\right)^2\right]\mathrm{erf}\left[{\displaystyle \frac{a_0}{2\alpha ^{1/2}}}\right].`$ (B16)
The expression for $`\lambda _s`$ is:
$`\lambda _s=\lambda _0\sqrt{\pi }\left({\displaystyle \frac{a_0}{2\alpha ^{1/2}}}\right)\mathrm{exp}\left[\left({\displaystyle \frac{a_0}{2\alpha ^{1/2}}}\right)^2\right].`$ (B17)
Hence, the asymptotic behavior of the meander (defined by Eq.(B8) for high slopes and Eq.(B14) for small ones) only depends on the size of the extrema region $`\lambda _0/2`$ (since the two parameters $`a_0`$ and $`\lambda _s`$ are linked to $`\lambda _0`$). How the plateau size is related to the model parameters will be considered in the next section.
### 2 Plateau size
In the one-sided limit we have $`\delta =0`$, and Eq.(B15) yields (in view of Eq.(B3)):
$`0=\alpha {\displaystyle \frac{h^{}}{1+h^2}}+\left({\displaystyle \frac{D_S\mathrm{}\mathrm{\Gamma }}{1+h^2}}+{\displaystyle \frac{D_La\mathrm{\Gamma }}{(1+h^2)^{1/2}}}\right)_{xx}\left({\displaystyle \frac{h^{}}{(1+h^2)^{1/2}}}\right).`$ (B18)
Introducing the abbreviation:
$$m=\frac{h^{}}{(1+h^2)^{1/2}},$$
(B19)
we can rewrite it in a familiar form:
$`\left({\displaystyle \frac{(D_S\mathrm{}+D_La)\mathrm{\Gamma }}{\alpha }}\right)m\mathrm{"}={\displaystyle \frac{m}{\stackrel{~}{\beta }(1m^2)^{1/2}+(1\stackrel{~}{\beta })}}={\displaystyle \frac{dU}{dm}},`$ (B20)
analogous to that describing the motion of a particle of position $`m`$ as a function of time $`x`$ in a potential $`U`$. We have defined $`\stackrel{~}{\beta }=D_S\mathrm{}/(D_S\mathrm{}+D_La)`$, and:
$`U(m)={\displaystyle _0^m}{\displaystyle \frac{mdm}{\stackrel{~}{\beta }(1m^2)^{1/2}+(1\stackrel{~}{\beta })}}.`$ (B21)
Multiplying Eq.(B20) by $`m^{}`$ and integrating with respect to $`x`$, we get the analogue of the ”energy conservation” condition:
$`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{(D_S\mathrm{}+D_La)\mathrm{\Gamma }}{\alpha }}\right)m^2+U(m)=U(m_0),`$ (B22)
where $`m_0`$ is the turning point (i.e. $`m^{}=0`$ when $`m=m_0`$). We look for solutions having two ”vertical” tangents where the step slope diverges: $`_x\zeta \pm \mathrm{}`$, (i.e. $`m_0\pm 1`$) in order to match the extrema solution with the large slope region. The size of the extrema region reads:
$`{\displaystyle \frac{\lambda _0}{2}}={\displaystyle _0^{\lambda _0/2}}𝑑x={\displaystyle _{m_0}^{m_0}}{\displaystyle \frac{dm}{m^{}}},`$ (B23)
where $`m_01`$. Using Eq.(B22), we find:
$`{\displaystyle \frac{\lambda _0}{\lambda _m}}={\displaystyle \frac{I(\stackrel{~}{\beta })}{\sqrt{2}}},`$ (B24)
where
$`\lambda _m=2\pi \left(2{\displaystyle \frac{(D_S\mathrm{}+D_La)\mathrm{\Gamma }}{\alpha }}\right)^{1/2}`$ (B25)
is the the wavelength of the most unstable mode obtained from linear analysis, and
$`I(\stackrel{~}{\beta })={\displaystyle \frac{1}{\pi \sqrt{2}}}{\displaystyle _1^1}{\displaystyle \frac{dm}{[U(1)U(m)]^{1/2}}}`$ (B26)
is plotted in Fig.4. $`I(\stackrel{~}{\beta })`$ is a decreasing function with $`I(0)=1`$ and $`I(1)0.54`$. Hence the extrema region size is finite and always smaller than $`\lambda _c=\lambda _m/\sqrt{2}`$.
The meander variation in this region, as compared to the total amplitude of the meander, decreases as $`t^{1/2}`$, and the step looks as if plateaus were present.
The reader is invited to repeat the calculation in the two-sided case –where $`\delta `$ is finite, in presence of pure line ($`\stackrel{~}{\beta }=0`$) or terrace diffusion ($`\stackrel{~}{\beta }=1`$). Surprisingly, the same integral $`I`$ appears. Let us define $`\stackrel{~}{\delta }=\mathrm{}/(\mathrm{}+d_++d_{})`$. We find in the pure line diffusion case:
$`{\displaystyle \frac{\lambda _0}{\lambda _m}}={\displaystyle \frac{I(\stackrel{~}{\delta })}{\sqrt{2}}},`$ (B27)
and in the pure terrace diffusion case:
$`{\displaystyle \frac{\lambda _0}{\lambda _m}}={\displaystyle \frac{I(1\stackrel{~}{\delta })}{\sqrt{2}}}.`$ (B28)
Hence, $`\lambda `$ never exceeds $`\lambda _m/\sqrt{2}=\lambda _c`$ the largest wavelength for which the meander is linearly stable.
We can use a similar treatment to analyze the case with higher order terms (eq.84). The main point is that terms proportional to $`ϵ^{1/2}`$ do not affect the long time behaviour obtained from the ansatz B4). Consequently, in large slope regions, we expect once again $`\zeta t^{1/2}`$. Terms breaking the front-back symmetry will affect differently maxima and minima regions, as seen in Fig.5, because the effective potential $`U(m)`$ is not invariant under the $`mm`$ transformation anymore, which corresponds to the up-down $`zz`$ for step meander. We shall not develop here further this point; for more details see.
The numerical solution of Eq.(B1) is performed in order to check the validity of the analytical results. First, the qualitative profile of the meander is in good agreement with the above description. We found the predicted scaling of the amplitude of the meander $`t^{1/2}`$ in all simulations performed so far, expect in the case of very small kinetic lengths $`(d_++d_{})/l<10^2`$, where we were not able to explore the late time behavior, due to bad numerical convergence.
As shown in Fig. 4 and 7, the observed plateau size, is in very good agreement with the prediction of Eq.(B24, B28,B27).
The value of $`a_0`$ is extracted from the evolution of the meander amplitude via the relation:
$`\zeta _{p,max}\zeta _{p,min}2a_0t^{1/2}.`$ (B29)
$`\lambda _s`$ is calculated from a fit of $`_x\zeta `$ at $`\zeta =0`$:
$`_x\zeta _s|_{\zeta =0}={\displaystyle \frac{4\sqrt{\pi }}{\lambda _s}}(\alpha t)^{1/2}.`$ (B30)
In Fig.10, both numerical values are compared to the prediction of Eq.(B16) in the one-sided limit, where $`\lambda _0`$ is calculated from Eq.(B24). Once again, good agreement is found.
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# Unquenched Charmonium with NRQCD
## 1 Charmonium on the lattice
One of the most rapidly expanding sectors of lattice QCD in the last decade has been the study of heavy-quark systems. Lattice simulations have successfully reproduced the broad structure of the heavy hadron spectrum, providing a solid piece of evidence for the correctness of QCD. Discrepancies at the level of the hyperfine structure still persist however, and in some cases these are uncomfortably large.
This paper describes a series of highly-improved non-relativistic simulations of the charmonium system, with the aim of estimating the sizes of various systematic uncertainties influencing the spectrum. An understanding of the relative influence of these uncertainties on the heavy-quark spectrum is vital to the interpretation of the current state of lattice simulations.
One very successful approach to simulating heavy quark systems utilises the NRQCD formalism , where the quark dynamics are governed by an effective non-relativistic Hamiltonian, expanded in powers of the heavy-quark velocity. For the bottom and charm quarks, $`v^20.1`$ and $`v^20.3`$ respectively, and so we expect to achieve some success with a non-relativistic theory. Simulations of heavy-light and heavy-heavy charm and bottom systems have shown that NRQCD captures much of the correct physics of the heavy quarks. Understanding the remaining systematic errors in heavy-quark simulations has become a major focus of the lattice NRQCD community.
The first report of a high-statistics NRQCD simulation of charmonium appeared in 1995 by Davies *et al.* . The authors used a NRQCD Hamiltonian with relativistic and discretisation errors corrected to $`𝒪(v^4)`$ to measure ground and excited $`S`$, $`P`$ and $`D`$ state energies in the quenched approximation. Agreement with experiment was very promising, with discrepancies at the order of $`10`$$`30\%`$ in $`S`$\- and $`P`$-state hyperfine spin-splittings, in agreement with the expected size of the next-order corrections.
Disturbingly, charmonium simulations incorporating $`𝒪(v^6)`$ corrections showed a large *decrease* in hyperfine spin splittings, taking theoretical predictions as much as $`50\%`$ further away from experimental values. These simulations also demonstrated a large dependence on the definition of the tadpole correction factor. Given the size of $`v^2`$ for charmonium, sizeable $`𝒪(v^6)`$ corrections are not surprising; however, the disappointingly large discrepancies in the spectrum with such a highly improved theory give pause to the future of charmonium simulations. Evidently, the NRQCD expansion converges slowly for the charm quark.
Even in the less-relativistic $`\mathrm{{\rm Y}}`$ system, the same highly-improved NRQCD action has not provided conclusive agreement with experiment . Certainly, NRQCD to $`𝒪(v^6)`$ is not a closed problem.
The difficulties with the hyperfine spectrum are not limited to the NRQCD approach. A report on the status of charmonium simulations with the relativistic Fermilab approach in 1993 cited a $`20`$$`30\%`$ shortfall for the $`S`$-state hyperfine splitting using an SW-improved Wilson action. In 1999, the UKQCD collaboration reported on a tadpole- and SW-improved simulation of charmonium ; their results for the $`S`$-hyperfine splitting were roughly $`40\%`$ below the experimental value. Both of these simulations used the quenched approximation, and the inclusion of dynamical quark loops would increase the hyperfine splittings. In the 1993 report, quenching effects were estimated to be as large as $`40\%`$, however this seems optimistic—corrections at the $`5`$$`15\%`$ level seem more typical in full QCD simulations of both the $`\mathrm{{\rm Y}}`$ system and of light hadrons .
A very recent report from the CP-PACS collaboration describes unquenched simulations of charmonium and bottomonium using NRQCD over a range of lattice spacings and sea-quark masses, with $`n_f=2`$ SW-improved light sea-quarks. In that report, the authors concentrate mostly on simulations of the $`b\overline{b}`$-system, though some charmonium results are presented. Their results indicate a significant increase in the $`S`$-state hyperfine splittings as the sea-quark mass decreases towards the chiral limit, though no effect is seen on the $`P`$-states.
We have performed a series of highly-improved NRQCD simulations to examine the various systematic uncertainties influencing the charmonium spectrum. We first concentrate on the effects of dynamical quark loops. If these account for the majority of the hyperfine splitting discrepancy in charmonium then we expect to find a large increase in the splittings when dynamical quarks are included, even in the NRQCD formalism. We examine this effect using an ensemble of unquenched configurations provided by the MILC collaboration, seeking to establish whether the effects of dynamical quarks are sufficient to reconcile the hyperfine discrepancy. This work does *not* aim to provide the definitive unquenched charmonium spectrum.
The remainder of the paper is devoted to an examination of other systematic effects. Simulations with two common definitions of the tadpole correction factor result in significantly different spectra, and we find a rough estimate of the effect of $`𝒪(\alpha _s)`$ radiative corrections to the NRQCD expansion coefficients. Finally, we note a sizeable shift in the hyperfine splittings due to an instability in the standard form for the heavy-quark propagator’s evolution equation. Each of these effects is contrasted with the estimated magnitude of the unquenching error, which leads us to several conclusions about NRQCD simulations of charmonium, and heavy-quark simulations in general.
## 2 The standard lattice NRQCD formalism
The NRQCD Hamiltonian is typically presented as an expansion in powers of the heavy-quark velocity. A highly-improved NRQCD Hamiltonian, with corrections to $`𝒪(v^6)`$ in the velocity expansion , is
$$H=H_0+\delta H_{v^4}+\delta H_{v^6},$$
(1)
where
$$H_0=\frac{𝚫^{(2)}}{2M_0},$$
(2)
is the leading kinetic Schrödinger operator, and the $`𝒪(v^4)`$ and $`𝒪(v^6)`$ corrections are
$`\delta H_{v^4}`$ $`=`$ $`c_1{\displaystyle \frac{1}{8M_0^3}}\left(𝚫^{(2)}\right)^2+c_2{\displaystyle \frac{ig}{8M_0^2}}\left(\stackrel{~}{𝚫}\stackrel{~}{𝐄}\stackrel{~}{𝐄}\stackrel{~}{𝚫}\right)`$ (3)
$`+c_3{\displaystyle \frac{g}{8M_0^2}}\sigma \left(\stackrel{~}{𝚫}\times \stackrel{~}{𝐄}\stackrel{~}{𝐄}\times \stackrel{~}{𝚫}\right)c_4{\displaystyle \frac{g}{2M_0}}\sigma \stackrel{~}{𝐁}`$
$`+c_5{\displaystyle \frac{a^2}{24M_0}}𝚫^{(4)}c_6{\displaystyle \frac{a}{16sM_0^2}}\left(𝚫^{(2)}\right)^2,`$
$`\delta H_{v^6}`$ $`=`$ $`c_7{\displaystyle \frac{g}{8M_0^3}}\{\stackrel{~}{𝚫}^{(2)},\sigma \stackrel{~}{𝐁}\}`$ (4)
$`c_8{\displaystyle \frac{3g}{64M_0^4}}\{\stackrel{~}{𝚫}^{(2)},\sigma \left(\stackrel{~}{𝚫}\times \stackrel{~}{𝐄}\stackrel{~}{𝐄}\times \stackrel{~}{𝚫}\right)\}`$
$`c_9{\displaystyle \frac{ig^2}{8M_0^3}}\sigma \stackrel{~}{𝐄}\times \stackrel{~}{𝐄}.`$
A *tilde* signifies the use of improved versions of the lattice operators that remove the leading discretisation errors: the improved lattice derivatives $`\stackrel{~}{\mathrm{\Delta }}`$ and $`\stackrel{~}{\mathrm{\Delta }}^{(2)}`$ are given by
$`\stackrel{~}{\mathrm{\Delta }}_\mu \psi (n)`$ $`=`$ $`\mathrm{\Delta }_\mu \psi (n){\displaystyle \frac{a^2}{6}}\mathrm{\Delta }_\mu ^3\psi (n),`$
$`\stackrel{~}{\mathrm{\Delta }}_\mu ^2\psi (n)`$ $`=`$ $`\mathrm{\Delta }_\mu ^2\psi (n)+{\displaystyle \frac{a^2}{12}}(\mathrm{\Delta }^2)^2\psi (n),`$ (5)
while the fields $`\stackrel{~}{E}_i=\stackrel{~}{F}_{4i}`$ and $`\stackrel{~}{B}_i=\frac{1}{2}ϵ_{ijk}\stackrel{~}{F}_{jk}`$ are taken from an improved gauge field tensor ,
$`\stackrel{~}{F}_{\mu \nu }(n)`$ $`=`$ $`{\displaystyle \frac{5}{3}}F_{\mu \nu }(n){\displaystyle \frac{1}{6}}[U_\mu (n)F_{\mu \nu }(n+\widehat{\mu })U_\mu ^{}(n)+U_\mu ^{}(n\widehat{\mu })F_{\mu \nu }(n\widehat{\mu })U_\mu (n\widehat{\mu })`$ (6)
$`U_\nu (n)F_{\mu \nu }(n+\widehat{\nu })U_\nu ^{}(n)+U_\nu ^{}(n\widehat{\nu })F_{\mu \nu }(n\widehat{\nu })U_\nu (n\widehat{\nu })].`$
All lattice operators are tadpole improved , by dividing all instances of the link operators $`U_\mu (n)`$ by the tadpole correction factor $`u_0`$,
$$U_\mu (n)\frac{U_\mu (n)}{u_0}.$$
(7)
This means, for example, that the gauge E and B fields are adjusted by a factor of $`u_0^4`$. Much evidence exists for the superiority of the *Landau* definition of the tadpole factor,
$$u_0^L=\frac{1}{3}\text{Tr}U_\mu _{_\mu A_\mu =0},$$
(8)
over the *plaquette* definition,
$$u_0^P=\frac{1}{3}\text{Tr}P_{\mu \nu }^{1/4}.$$
(9)
For example, $`u_0^L`$ leads to smaller corrections to hyperfine splittings, and better scaling of quarkonium masses ; it restores rotational invariance to a greater degree in the static quark potential ; and it results in closer agreement between the tadpole-improved value and the perturbative value for the ‘clover’ coefficient $`c_{sw}`$ in the Sheikholeslami-Wohlert action . We have used both the Landau and plaquette definition in our simulations.
Since the quarks and antiquarks are decoupled in the non-relativistic theory, the heavy-quark Green’s function may be found from an evolution equation,
$$G_{t+1}=\left(1\frac{aH_0}{2s}\right)^sU_4^{}\left(1\frac{aH_0}{2s}\right)^s\left(1a\delta H\right)G_t,$$
(10)
with the initial time-step given by
$$G_1=\left(1\frac{aH_0}{2s}\right)^sU_4^{}\left(1\frac{aH_0}{2s}\right)^s\delta _{x,0}.$$
(11)
The $`\left(1aH\right)`$ factors are linear approximations to the continuum evolution operator $`e^{Ht}`$. The ‘stabilisation parameter’ $`s`$ appearing in Equations (1) and (10) improves the approximation to the time evolution operator $`e^{aH}`$.
To complement the use of a highly-improved quark Hamiltonian, we use a tadpole and ‘rectangle’ improved action for the gauge fields ,
$$S_G=\beta \underset{n,\mu >\nu }{}\left(\frac{5}{3u_0^4}P_{\mu \nu }(n)\frac{1}{12u_0^6}(R_{\mu \nu }+R_{\nu \mu })\right),$$
(12)
where $`P_{\mu \nu }(n)`$ and $`R_{\mu \nu }`$ represent the traces of $`1\times 1`$ plaquettes and $`2\times 1`$ rectangles of link operators respectively.
Operators for the various quarkonium states have the form
$$M(t)=\underset{n}{}\psi ^{}(n,t)\mathrm{\Gamma }(n)\chi ^{}(n,t),$$
(13)
where $`\psi ^{}`$ and $`\chi ^{}`$ are the quark and antiquark creation operators, and $`\mathrm{\Gamma }(n)`$ provides the appropriate spin and spatial wavefunction quantum numbers. Operators for the lowest-lying $`S`$, $`P`$ and $`D`$ states are given in a number of references ; using these, we have constructed propagators for each of the $`{}_{}{}^{2S+1}L_{J}^{}=^1S_0`$, $`{}_{}{}^{3}S_{1}^{}`$, $`{}_{}{}^{1}P_{1}^{}`$, $`{}_{}{}^{3}P_{0}^{}`$, $`{}_{}{}^{3}P_{1}^{}`$ and $`{}_{}{}^{3}P_{2}^{}`$ states. Only one spin polarisation of each of the triplet states was used.
To reduce the effects of excited-state contamination and improve the operators’ overlap with the true meson ground-state wavefunctions, we have used a gauge-invariant smearing function, replacing
$$\mathrm{\Gamma }(n)\mathrm{\Gamma }(n)\varphi _{sm}(n).$$
(14)
A simple and effective choice for $`\varphi _{sm}`$ is
$$\varphi _{sm}(ϵ,n_s)=\left(1+ϵ\mathrm{\Delta }^2\right)^{n_s}.$$
(15)
The weighting factor $`ϵ`$ and number of smearing iterations $`n_s`$ were tuned to optimise the overlap with the ground state.
## 3 Details of the simulations
We have performed a number of different simulations of the charm system, to compare the magnitudes of various systematic effects on the spectrum. We obtained results with the NRQCD Hamiltonian in Equation (1) truncated to $`𝒪(v^4)`$ and $`𝒪(v^6)`$, with both the Landau and plaquette definitions for the tadpole factor $`u_0`$.
To examine the size of dynamical quark effects, we obtained an ensemble of 200 unquenched gauge field configurations, generously provided by the MILC collaboration . The configurations were created with the Wilson gluon action at $`\beta =5.415`$, with two flavours of staggered dynamical quarks at $`m=0.025`$. This light quark mass corresponds to a pseudoscalar-to-vector meson mass ratio of $`m_{ps}/m_v0.45`$. The lattice volume of these configurations is $`16^3\times 32`$—with a spacing of $`a0.16`$ fm (determined from the charmonium spectrum as described below), this corresponds to a lattice extending roughly 2.5 fermi in each spatial direction.
We produced an ensemble of quenched configurations with both the Landau and plaquette tadpole definitions, using the improved action in Equation 12. We found that, using Landau and plaquette tadpoles respectively, $`\beta =2.1`$ and $`\beta =2.52`$ give almost the same lattice spacing as the unquenched configurations. These results agree with the spacings given in Reference at the same values of $`\beta `$. We created 100 configurations in each case, with lattice volume $`12^3\times 24`$, the largest we were able to manage with our computational resources. Given the small physical size of the heavy mesons, however, the difference in volume between the quenched and unquenched configurations should not have an effect on our results.
The lattice spacing was determined for each ensemble using the spin-averaged $`P`$$`S`$ splitting, for charmonium $`E(PS)=458`$ MeV. This splitting is known to be quite independent of the heavy quark mass, falling only slightly to 440 MeV for bottomonium, and so serves as a stable quantity for determining the physical lattice spacing. We have collected the parameters of our simulations together in Table 1.
The *kinetic mass* $`M_k`$ of a boosted state with momentum p is defined by
$$E(𝐩)=E(0)+\frac{𝐩^2}{2M_k}+𝒪(𝐩^4).$$
(16)
The bare charm quark mass $`M_0`$ is tuned by requiring that the kinetic mass of the $`{}_{}{}^{1}S_{0}^{}`$ charmonium state agrees with the experimental mass of the $`\eta _c`$, $`M_{\eta _c}=2.98`$ GeV. We created correlators for a boosted state with $`𝐩=(\frac{2\pi }{L},0,0)`$, where $`L`$ is the spatial extent of the lattice. The tuned bare masses, and their corresponding physical (kinetic) masses for the $`{}_{}{}^{1}S_{0}^{}`$, are shown in Table 1.
Meson correlators were calculated for the various charmonium states, using smeared meson operators with $`n_s=8`$ and $`ϵ=1/12`$ in Equation (15) at both the source and sink. To decrease statistical uncertainties, we calculated more than one meson correlator for each gauge field configuration. Meson sources were situated at four different spatial origins—$`(0,0,0)`$, $`(L/2,L/2,0)`$, $`(L/2,0,L/2)`$ and $`(0,L/2,L/2)`$—and starting from two time slices, at $`t=0`$ and $`t=12`$, for a total of 800 meson correlator measurements for each state.
Statistical correlations will exist between the multiple measurements of the propagators within each configuration, however the small size of $`Q\overline{Q}`$ systems (the $`c\overline{c}`$ is roughly 0.5 fm in radius) is some justification for this practice. The correlations are expected to be small, as noted in other charmonium studies with similar lattice spacings .
Masses for the various $`c\overline{c}`$ states were found by fitting the correlators with a single exponential,
$$G_M(t>t_{min})=c_Me^{E_Mt}$$
(17)
after a minimum time $`t_{min}`$, allowing for suitable suppression of excited state contributions. Energy splittings between correlated states, such as the $`S`$-state hyperfine splitting $`\mathrm{\Delta }E=E(^3S_1)E(^1S_0)`$, can often be extracted more precisely by fitting to a ratio of their two correlators,
$$(t)=\frac{G_B(t)}{G_A(t)}\frac{c_Be^{E_Bt}}{c_Ae^{E_At}}=\frac{c_Be^{(E_A+\delta E)t}}{c_Ae^{E_At}}=c_{}e^{\delta Et}.$$
(18)
We used ratio fits to extract the $`S`$-state hyperfine splitting, and the kinetic mass from the boosted $`{}_{}{}^{1}S_{0}^{}`$ state. Attempts to extract $`P`$-state hyperfine splittings in this manner were unsuccessful, as statistical noise overtook the very small signal before a reasonable plateau emerged. Single-exponential fits, however, resolved the three $`{}_{}{}^{3}P`$ levels. We have not employed a bootstrap analysis for the fit results, which may suggest we have overestimated the statistical uncertainties.
In the following sections, we present the results for a range of simulations, incorporating all combinations of quenched and unquenched gauge configurations, $`𝒪(v^4)`$ and $`𝒪(v^6)`$ correction terms, and Landau and plaquette tadpole factors.
### 3.1 Quenched results
An example of the quality of the correlator data is shown in Figures 1 and 2, plots of the effective masses for the $`{}_{}{}^{1}S_{0}^{}`$, $`{}_{}{}^{1}P_{1}^{}`$ and $`{}_{}{}^{3}P_{0}^{}`$ from the simulation using the Landau tadpole factor. The meson propagators were fit with single exponentials over a range of time intervals $`(t_{min}:t_{max})`$. An indication of the convergence of these fits is given in Table 2, where the fit results are shown for the $`𝒪(v^6)`$ simulations using the plaquette tadpole factor. The results presented in this table are representative of all of the charmonium spectra we present here. The two $`S`$-states had a much cleaner signal than the four $`P`$-states, evident in the lower value for $`t_{max}`$ used for the $`P`$-state fits.
Table 3 contains the final results for the quenched charmonium mass fits. We considered the ground-state for each meson propagator to have properly emerged when three consecutive $`t_{min}:t_{max}`$ intervals gave results that agreed within statistical errors; the meson mass was then taken as the middle of these three values. The masses are given in both lattice units and physical units, using the values for $`a^1`$ in Table 1 to provide the physical energy scale. The simulated spectra are displayed in Figures 3 and 4, shown against the experimental data.
### 3.2 Unquenched Results
Given the similar lattice spacings of the MILC configurations and our own quenched ensembles, we have used almost the same parameter set for the unquenched charmonium simulations—the lower half of Table 1 shows the specific parameters used. The results of the unquenched simulations are contained in Table 4, with the physical energy scale set by $`a^1=1.21(2)`$ GeV, again from the spin-averaged $`P`$$`S`$ splitting. The spectra are shown in Figures 5 and 6.
## 4 Discussion of the Spectra
A cursory comparison of the quenched and unquenched results shows that, while the qualitative structure of the spectrum appears, precision NRQCD simulations of the charmonium system have a number of issues yet to be resolved. This is most readily seen in the hyperfine splittings, which are collected in Figures 7 and 8, and compared in Table 5.
Consider first the quenched results. The $`𝒪(v^6)`$ corrections lead to a disturbingly large decrease in the hyperfine splittings, taking them further away from the experimental values by as much as $`60\%`$. The situation for the plaquette-tadpole simulations is strikingly bad, where the $`{}_{}{}^{3}P`$ states appear in the wrong order. This reversal is corrected in the Landau-tadpole simulations, though the hyperfine splittings are still badly underestimated.
These difficulties are not new—Trottier first drew attention to the large $`𝒪(v^6)`$ corrections to the $`S`$-state hyperfine splitting in 1996, and noted a possible problem with the $`{}_{}{}^{3}P`$-state ordering. Trottier and Shakespeare examined the effects of the different tadpole definitions $`u_0^P`$ and $`u_0^L`$ on the $`S`$-state hyperfine splitting. They performed $`𝒪(v^6)`$-improved NRQCD simulations using both tadpole schemes, across a wide range of lattice spacings, and drew a number of important conclusions; most notably, the $`𝒪(v^6)`$ hyperfine corrections with Landau tadpoles were significantly smaller than the plaquette tadpole results.
We have confirmed a number of these results here, and in particular clearly resolved the extremely poor $`{}_{}{}^{3}P`$-state behaviour, most notably when $`u_0^P`$ is used. This may simply be a problem due to the bare charm mass falling below one in these simulations. However, the $`u_0^L`$ simulations lead to a higher bare $`c`$-quark mass for a given lattice spacing, and the very low $`P`$-state hyperfine splitting even with $`aM_0>1`$ suggests that these problems extend beyond the size of the bare mass.
### 4.1 Evidence for Quenching Effects?
The large discrepancies in spin-dependent splittings would be less worrisome if quenching were seen to have a considerable effect on the spectrum, as suggested in . Sadly, this does not seem to be the case. There is some evidence for a difference between the quenched and unquenched simulations in the $`𝒪(v^6)`$ $`S`$-hyperfine data, perhaps as much as ten percent. However, given the apparent size of other systematic uncertainties, no great significance can be attached to these differences.
We must address the difference between the quenched and unquenched gluon actions—the MILC configurations were created with the Wilson plaquette action, while we have employed the rectangle-improved action for the quenched lattices. We therefore anticipate an $`𝒪(a)`$ error entangled with the effects of the dynamical quarks. Our quenched $`𝒪(v^4)`$ results can be compared with the results from Reference , where the plaquette action was used at roughly the same lattice spacing. We see a $`10`$ MeV difference between the $`S`$-hyperfine splittings in the two simulations.
We wish to reiterate our goal, however, to see whether the dynamical quark effects are *large* or *small*. The $`S`$-hyperfine splitting, even in relativistic simulations, falls short of experiment by 40 to 50 MeV. An unquenching effect of this magnitude would be visible, even taking differences in gluon action into account. No such effect was observed in these simulations, and we therefore suggest that quenching effects are small in this sense.
This conclusion is supported by results in high-precision $`\mathrm{{\rm Y}}`$ simulations , where the $`P`$-state hyperfine splitting is still somewhat underestimated in unquenched simulations of this highly-nonrelativistic system, despite the use of the $`𝒪(v^6)`$ improved NRQCD action. Very recently, a $`10\%`$ sea-quark effect was seen in the hyperfine splittings of the charmonium and bottomonium system in Reference , but differences between the $`n_f=0`$ and $`n_f=2`$ $`P`$-state splittings were not significant compared with other systematic uncertainties. Recent results with unquenched lattices in the $`B`$ meson spectrum have also shown no significant differences between $`n_f=0`$ and $`n_f=2`$ dynamical quark flavours .
### 4.2 Other Systematic Errors
The preceding results suggest that agreement between lattice simulations and experiment in quarkonium systems will likely not improve through the effects of dynamical quarks alone. In the remainder of this section we explore various other systematic errors that impact on heavy-quark simulations, as a contrast to the small quenching effects found above.
#### 4.2.1 The Choice of the Tadpole Factor
We have seen, as others have previously, large differences between results using the Landau tadpole factor $`u_0^L`$, and those with the plaquette definition $`u_0^P`$. In our own simulations, the size of the $`𝒪(v^6)`$ corrections with $`u_0^L`$ are significantly smaller than the plaquette tadpole results. This is not surprising: the E and B fields are each multiplied by a factor of $`u_0^4`$ in the tadpole-improved theory. On our lattices,
$$\left(\frac{u_0^P}{u_0^L}\right)^4=\{\begin{array}{c}1.24\text{(Quenched)}\hfill \\ 1.30\text{(MILC)}\hfill \end{array}$$
(19)
Terms in the NRQCD Hamiltonian linear in E or B will differ by as much as $`30\%`$ between the different tadpole improvement schemes.
As noted earlier, the evidence in favour of Landau tadpoles is strong. Our simulations offer further support, particularly in the $`{}_{}{}^{3}P`$-state behaviour, though the more salient issue here is that tadpole effects are at least as important as quenching effects in our simulations.
#### 4.2.2 Radiative Corrections
We expect some effect on the spectrum from high-momentum modes that are cut off by the finite lattice spacing. These high-energy effects may be calculated in perturbative QCD as $`𝒪(\alpha _s)`$ radiative corrections to the coefficients of the NRQCD expansion, and there are indications that these may be large for the charm quark. Lattice perturbation theory calculations of corrections to $`c_1`$ and $`c_5`$, the ‘kinetic’ terms in Equation (1), have been completed by Morningstar . The corrections are roughly $`10\%`$ or less for the bottom quark, but rise dramatically as the bare quark mass falls below one (in lattice units). In typical simulations, the bare charm quark mass sits close to unity, and so these corrections may become quite significant.
It is possible to find these radiative corrections without performing long calculations in lattice perturbation theory, by using Monte Carlo simulations at very high values of $`\beta `$ . Such ‘non-perturbative’ perturbative results have been obtained by Trottier and Lepage for the spin-dependent $`c_4`$ term in the $`𝒪(v^4)`$ NRQCD Hamiltonian, Equation (1). Unfortunately, radiative corrections to the remaining terms in the NRQCD Hamiltonian have not been calculated to date.
We performed a ‘toy’ simulation to roughly estimate the effects of $`𝒪(\alpha _s)`$ corrections to all terms in the NRQCD Hamiltonian, replacing the tree level coefficients $`c_i=1`$ with $`c_i=1\pm \alpha _s`$. A rough estimate of $`\alpha _s`$ can be made from the (tadpole-improved) parameters of our simulations,
$$\alpha _s(\pi /a)\alpha _{lat}^{TI}+𝒪(\alpha ^2)\frac{g^2}{4\pi }=\frac{6}{4\pi \beta u_0^4}.$$
(20)
For our values of $`\beta `$ and $`u_0`$, this gives $`\alpha _s0.15`$$`0.2`$. For the three terms in the Hamiltonian where perturbative analysis has been performed, we used the calculated values ; for the remaining terms, we varied the coefficients between 0.8 and 1.2.
Altering the coefficients in this way, we found that the charmonium $`S`$\- and $`P`$-hyperfine splittings changed by as much as $`10`$$`40\%`$, depending on the sign of the corrections for each individual $`c_i`$. While this is only a crude estimate, it is clear that the effects of radiative corrections may be as important as quenching effects for heavy-quark systems. Accurate determinations of the remaining $`𝒪(\alpha _s)`$ corrections are sorely needed.
#### 4.2.3 Improving the Evolution Equation
The evolution equation we presented in Section 2 for the heavy-quark propagator, Equation (10), contains better-than-linear approximations to the exponential $`e^{Ht}`$ for the terms involving the zeroth-order Hamiltonian $`H_0`$, but only a linear approximation for the correction terms $`\delta H`$. Noting that the high-order corrections are quite large for charmonium, it is conceivable that this lowest-order approximation is too severe. A similar conclusion was made by Lewis and Woloshyn of their NRQCD simulations of the $`D`$ meson spectrum . The authors were able to remove some spurious effects due to large vacuum expectation values of one of the high-order terms in their NRQCD Hamiltonian , by improving the exponential approximation for the $`\delta H`$ terms in the evolution equation.
The coefficients of the $`𝒪(v^6)`$ terms include high powers of $`M_0^1`$ and $`u_0^1`$, and it is conceivable that for the charm quark, with $`aM_01`$, the $`(1a\delta H_{v^6})`$ approximation is poor. We examined this possibility for the $`𝒪(v^6)`$ terms, by using an improved form for the evolution equation that incorporates a ‘stabilisation’ parameter for the correction terms, with the replacement
$$(1a\delta H)\left(1\frac{a\delta H}{s_\delta }\right)^{s_\delta }.$$
(21)
We have performed a simulation with this alteration to the evolution equation, with $`s_\delta =4`$. Otherwise, all other parameters were kept the same as the previous Landau-tadpole quenched simulations. In general, altering the evolution equation will lead to a change in the bare charm quark mass $`M_0`$. In this case we found that $`M_0=1.15`$ once again gave a value of $`3.0(1)`$ GeV for the $`{}_{}{}^{1}S_{0}^{}`$ mass.
The the improved evolution equation altered the $`S`$-hyperfine splitting significantly, increasing it by roughly $`40\%`$ to $`70`$ MeV. The statistical uncertainties in the $`P`$-hyperfine splittings were large, though a similar increase seems likely. These results suggest the linear approximation $`(1a\delta H)`$ typically used in NRQCD simulations is not sufficiently accurate for the large corrections encountered at the charm quark mass.
## 5 Conclusions
Lattice NRQCD simulations of heavy-quark systems have evolved greatly over the last decade. By incorporating high-order interaction terms to counter relativistic and discretisation errors, simulations now routinely produce results that agree with experiment at the $`10`$$`30\%`$ level. However, stubborn discrepancies remain in highly-improved simulations, typically performed in the quenched approximation, or at tree-level in the $`𝒪(\alpha _s)`$ expansion, or both. To proceed further, all remaining systematic errors must be addressed.
Past studies strongly suggest that the NRQCD expansion converges slowly for the charm quark, with the leading and next-to-leading order corrections apparently oscillating in sign. To $`𝒪(v^6)`$, the hyperfine spin-splittings fall short of experimental values by $`50\%`$ or more. Without knowing the magnitude of the next-order corrections in the velocity expansion, the question of reducing the disparities in the charmonium spectrum seems academic.
While the NRQCD approach appears to be problematic for the charmonium system, relativistic lattice formalisms have their share of difficulties. Simulations of charmonium with a variety of quark actions—NRQCD, the Fermilab actions, the D234 action—all underestimate the $`S`$-hyperfine splitting by at least $`40`$ MeV (see for a good summary).
There are sound reasons for estimating the size of dynamical quark effects in the charmonium system. Some have suggested the remaining hyperfine discrepancy is due to quenching; estimates of the effects of dynamical quark loops range as high as $`40\%`$ . Our results indicate this is unlikely to be the case—we find at most a $`10\%`$ difference between our quenched and unquenched hyperfine splittings. As the quenching effects are apparently small for the range of different quark interactions present in the NRQCD action, we suggest that they will also be small across other quark actions.
We recognise several shortcomings in our study: we have used different gluon actions for quenched and unquenched simulations, we have only examined the effects of unquenching at a single dynamical quark mass and a single lattice spacing, and we have not attempted to extrapolate to the physical case of three light sea quark flavours. The first of these issues was discussed in Section 4.1 above. To address the other objections with further simulations is beyond our present computational resources. In any case, such efforts are perhaps justified in simulations of the $`b`$ quark, where systematic uncertainties are under better control and quenching effects are probably of comparable size to discretisation and radiative effects. For the charm system however, the much larger high-order relativistic errors and the large tadpole corrections dominate the effects of quenching.
The sensitivity of the NRQCD corrections to the choice of tadpole factor is well established. This sensitivity should disappear with a higher-order treatment of the tadpole loops (and other radiative corrections) in lattice simulations. In practice, such a treatment is not yet available, and some choice for the tadpole factor is required. Our results add to the growing list of evidence if favour of calculating the tadpole correction factor from the mean link in the Landau gauge, in preference to the plaquette definition.
The large effects we have encountered due to instabilities in the evolution equation should also be investigated further. These instabilities are doubtless amplified in simulations of the charm quark, where the convergence of the NRQCD expansion is already questionable. Using an improved evolution equation, as we have demonstrated, may bring the NRQCD approach into agreement with other quenched relativistic results for charmonium.
Further, we have shown that $`𝒪(\alpha _s)`$ radiative corrections may shift the spin-splittings by as much as $`40\%`$. While this is a crude estimate, the possibility of such sizeable corrections in comparison with the small quenching effect gives us pause for consideration. Of particular note are unquenched results for the $`\mathrm{{\rm Y}}`$ spectrum in , using the $`𝒪(v^6)`$ Hamiltonian, which indicate that remaining discrepancies with experiment are at the ten percent level—conceivably within the reach of radiative corrections. Perturbative calculations of the remaining radiative corrections to the NRQCD coefficients, and those in other actions as well, will likely be necessary in the near future.
We thank Howard Trottier and Randy Lewis for stimulating discussions. This work has been supported by the Natural Sciences and Engineering Research Council of Canada.
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# Systematic analysis of Δ𝐼=4 bifurcation in 𝐴∼150 superdeformed nuclei: Active orbitals
\[
## Abstract
A simple criterion for the $`\mathrm{\Delta }I=4`$ bifurcation is applied to thirty superdeformed bands in the $`A`$150 mass region. The analysis allows to differentiate between active and inactive for staggering single-particle states. The consideration is based on additivity of the nonaxial hexadecapole moment, which plays a key role in the phenomenon.
\]
The $`\mathrm{\Delta }I=4`$ bifurcation, or the $`\mathrm{\Delta }I=2`$ staggering, is a well known mysterious phenomenon in the physics of superdeformed (SD) bands. It is observed as regular deviations of the $`\gamma `$–ray energy differences $`\mathrm{\Delta }E_\gamma `$ from the smooth behavior. In spite of having the energy scale of tens of electron-volt the phenomenon has been a topic of considerable interest and a lot of the experimental and theoretical works have been devoted to it. The reason lies in the unusual for the nuclear physics period of oscillations, $`\mathrm{\Delta }I=4`$, and their long and regular character. Recently attention to this problem is called again by the versatile analysis of the experimental data undertaken in Refs. .
Among other explanations it has been suggested in Ref. that the $`\mathrm{\Delta }I=4`$ bifurcation is caused by the coupling of rotation with single-particle motion in an axially symmetric nuclei. The coupling is effected by the quadrupole and hexadecapole two-body interactions through the term
$$V_{\mathrm{cpl}}=\underset{\lambda =2,4}{}\chi _\lambda \underset{\mu 0}{}Q_{\lambda \mu }𝒬_{\lambda \mu },$$
(1)
where $`\chi _\lambda `$ are the interaction strengths, $`Q_{\lambda \mu }`$ and $`𝒬_{\lambda \mu }`$ are respectively the perturbative and nonperturbative nonaxial multipole moments induced by rotation. The separation of multipole moments in the two parts is explained by the two types of the single-particle states involved in superdeformation: low-$`j`$ natural and high-$`j`$ intruder orbitals. The former have the natural parity and usually low angular momentum $`j`$. They are weakly sensitive to the nuclear rotation and may be described by the perturbation theory. Therefore, the perturbative parts are proportional to nonaxial components of the total angular momentum $`𝐈`$:
$$𝒬_{\lambda \pm 2}=\alpha _{\lambda 2}I_\pm ^2,𝒬_{4\pm 4}=\alpha _{44}I_\pm ^4,$$
(2)
where $`I_\pm =I_1\pm I_2`$, $`I_1`$ and $`I_2`$ are the projections on the coordinate axes perpendicular to the symmetry axis 3. On the other hand, the intruder orbitals are very sensitive to rotation and, in addition, each intruder state is affected differently by rotation. Therefore the quantity $`Q_{\lambda \mu }`$ cannot be treated by the perturbation theory.
The coupling (1) distorts the rotational motion of an axially symmetric nucleus. Accordingly, the effective rotational Hamiltonian for an isolated band,
$$H_{\mathrm{eff}}=𝒜𝐈^2+𝐈^4+d(I_+^2+I_{}^2)+c(I_+^4+I_{}^4),$$
(3)
contains nonadiabatic terms. The last two nonaxial terms split the single band characteristic of an axial nucleus in a series of bands, which correspond to the different orientation of the vector $`𝐈`$. The Hamiltonian (3) is adequate only for the description of the yrast band, for which the angular momentum $`𝐈`$ is perpendicular to the symmetry axis. All other bands are separated from the yrast one by a large energetic gap caused by the small nonaxial deformation induced by rotation. They are of no concern for us. The last term in (3) proportional to $`\chi _4`$ is an essential ingredient for the $`\mathrm{\Delta }I=2`$ staggering. On the other hand, the term with the operator $`I_+^2+I_{}^2`$ breaks a fourfold symmetry and makes the staggering pattern irregular (see also Ref.). In other words, the nonaxial terms of the Hamiltonian (3) crimp the rotational energy surface. The short wave crimps near the stationary point (i.e., the axis of rotation) are important for staggering. However the crimped surface does not yet solve the problem. The staggering may exist if the stationary point is a minimum, which happens for $`c>0`$. For the negative value of this parameter, the staggering is absent in the yrast band, but exists in the uppermost one because the transformation $`cc`$ results in the inversion of multiplet levels. The sign of $`d`$ is not important for staggering.
The parameter $`c`$ involves the perturbative and nonperturbative factors. For a nucleus with $`Z`$ protons ($`\pi `$) and $`N`$ neutrons ($`\nu `$) the necessary condition for the existence of the $`\mathrm{\Delta }I=4`$ bifurcation has the form
$$c=\chi _4\left[\left(\frac{2Z}{A}\right)^{2/3}Q_{44}(\pi )+\left(\frac{2N}{A}\right)^{2/3}Q_{44}(\nu )\right]$$
$$\times [\left(\frac{2Z}{A}\right)^{2/3}\alpha _{44}(\pi )+\left(\frac{2N}{A}\right)^{2/3}\alpha _{44}(\nu )]>0,$$
(4)
where
$$Q_{44}(\tau )=\underset{n,\alpha }{}n\alpha \tau |q_{44}|n\alpha \tau ,\tau =\pi ,\nu ,$$
(5)
and the summation extends over all the occupied nonperturbative or, active single-particle states having the quantum numbers $`n`$ and the signature $`\alpha `$. We will use for $`n`$ the asymptotic quantum numbers $`[Nn_z\mathrm{\Lambda }]\mathrm{\Omega }`$ of a nonrotating nucleus. The pairing effects are neglected and the simplest shell model is used. It should be noted that the condition (4) is necessary but insufficient, due to a rather general form of the rotation-single-particle interaction obtained from Eq. (1). More detailed information concerning this interaction is required to get a sufficient condition and to reproduce staggering patterns.
The inequality (4) was used in Ref. to check the theory on eight SD bands in the $`A150`$ nuclei. The quantity $`\alpha _{44}`$ was calculated in the anisotropic harmonic oscillator (HO) potential and $`Q_{44}`$ was calculated in the limit of an isolated intruder shell. It has been shown that the later quantity is the fluctuating function of the number of the nucleons occupying intruder orbitals. Thus the parameter $`c`$ may change sign and staggering appears or disappears with the change of the intruder configuration of a SD band. The correlation between the $`\mathrm{\Delta }I=4`$ bifurcation and the sign of $`c`$ has been found. However, the model in use is not reliable for the SD bands because the intruder shells cannot be treated as isolated.
In this paper, we report the results of the analysis of thirty SD bands, including those from Ref. . We employ the realistic modified oscillator (MO) potential by using the code GAMPN . The $`\kappa `$ and $`\mu `$ parameters have been taken from Ref.. The expectation values $`q_{44}(n\alpha \tau )=n\alpha \tau |q_{44}|n\alpha \tau `$ involved in $`Q_{\lambda \mu }`$ and $`𝒬_{\lambda \mu }`$ are calculated with the wave functions,
$$\psi _{n\alpha }=\underset{lj\mathrm{\Omega }}{}a_{lj\mathrm{\Omega }}^{n\alpha }|N_{rot}lj\mathrm{\Omega },$$
(6)
of the cranking potential, where $`N_{rot}`$ is the principal quantum number in the stretched rotating basis. The small coupling between different $`N_{rot}`$-shells is neglected.
In order to single out the active orbitals we have investigated how the expectation values $`q_{44}(n\alpha )`$ depend on the rotational frequency $`\omega `$. All the single-particle states occupied by neutrons and protons in the $`A`$150 nuclei have been considered. The following conclusions can be drawn. (i) There are the three different patterns of the $`\omega `$-dependence, which are shown in Fig.1. Each of the patterns is associated with the specific types of orbitals. (ii) The perturbative (Fig.1a) and nonperturbative (Fig.1b) dependencies associate with inactive and active for staggering orbitals respectively. These orbitals are the same for neutrons and protons. (iii) The set of active orbitals is supplemented with the inactive states interacting with the active ones. The orbitals $`\nu `$5/2 and $`\nu `$1/2 placed above the neutron gap at $`N=80`$ represent a classical example of an avoid level crossing . Figure 1c shows $`q_{44}`$ as a function of $`\omega `$ for the two signature branches of these orbitals. According to the Strutinsky prescription, we use the renormalization factor 1.27 to have the possibility of a comparison with observed rotational frequencies. At low frequency the 1/2 orbitals are inactive, whereas at high frequency they involve large admixture of the active, 5/2, orbitals and have the nonperturbative dependence of $`q_{44}`$. These interacting orbitals with the positive signature carry considerable hexadecapole moment. Thus the removal or addition of a neutron in these states may change the inequality (4). The orbitals $`\pi `$1/2 and $`\nu `$1/2 with both signatures are active practically for all frequencies because their avoid crossings with $`\pi `$3/2 and $`\nu `$5/2 occur at low frequencies. Some other active states induced by an avoid crossing are presented in Table IV. Their contribution to the quantities $`Q_{44}(\pi )`$ or $`Q_{44}(\nu )`$ is moderate. All the considered pairs of interacting states belong to the same $`N_{rot}`$-shell. The coupling between different $`N_{rot}`$-shells generates the avoid crossings, which change the moments $`q_{44}`$ only slightly due to a small interaction.
The set of active for staggering states is shown in Fig.2. These are not necessary the intruder orbitals but all have the asymptotic quantum numbers $`\mathrm{\Omega }`$=3/2 and 5/2. Compared to inactive orbitals, active ones carry the large values of $`q_{44}`$, which are nevertheless smaller that the hexadecapole $`q_{40}`$ moment of the single-particle states around the SD core of <sup>152</sup>Dy . As a rule, the moments of the states with the same asymptotic quantum numbers have close absolute values and opposite signs for different signatures. They almost offset each other. Thus, the contribution of the active orbitals to the total nonaxial hexadecapole moment has the same order as that of inactive ones. Consequently the equilibrium deformation $`\epsilon _{44}`$ at a high rotational frequency is small .
To explain these findings let us consider a cranking isolated $`j`$-shell with the Hamiltonian
$$H_j=H_{j0}+ϵj_3^2\omega j_1,$$
(7)
where $`H_{j0}`$ is a spherical part and $`ϵ`$ is proportional to an axial deformation. Assuming the small rotational frequency $`\omega `$, we shall use perturbation theory with the unperturbed function in the signature representation
$$u_{j\mathrm{\Omega }\alpha }=(|j\mathrm{\Omega }+e^{i\pi (j\alpha )}|j\mathrm{\Omega })/\sqrt{2}.$$
(8)
It is easy to see that the only expectation values of $`q_{44}(𝐫)`$, which are proportional to $`\omega `$, are those for the states with $`\mathrm{\Omega }`$=3/2 and 5/2:
$$q_{44}(j3/2\alpha )=q_{44}(j5/2\alpha )=\frac{\omega }{ϵ}q_4f(j)e^{i\pi (j\alpha )},$$
(9)
where the form of the positive function $`f`$ is inessential for us. For other $`\mathrm{\Omega }`$, the first non-vanishing contribution to $`q_{44}(j\mathrm{\Omega }\alpha )`$ is proportional to the higher powers of $`\omega `$. The equality (9) explains the signature dependence of the values $`q_{44}`$ for almost all active orbitals shown in Fig.2 because they are mostly high-j intruder or high-j ones with a rather good quantum number $`j`$. An anomalous signature dependence is observed for the five states with small $`j`$.
In a similar way as in Ref. , we now check the criterion (4) bearing in mind that the third multiplier is negative for all the bands considered below. Really, straightforward calculation within the MO potential show that the perturbative quantity $`\alpha _{44}`$ is negative as in the case of the HO model . Thus we will be interested only in the sign of the second multiplier
$$Q_{44}=\left(\frac{2Z}{A}\right)^{2/3}Q_{44}(\pi )+\left(\frac{2N}{A}\right)^{2/3}Q_{44}(\nu ),$$
(10)
where the moments $`Q_{44}(\pi )`$ and $`Q_{44}(\nu )`$ are calculated by using additivity of contributions from individual orbitals according to Eq. (5). The main difficulty in their calculation is the nuclear equilibrium deformation, since the shape trajectories in the $`(\epsilon ,\epsilon _4)`$-plane are known for limited number of SD bands. Starting with these bands, we give in Table I the estimated values of $`Q_{44}`$ for three rotational frequencies. The corresponding parameters $`\epsilon `$ and $`\epsilon _4`$ have been taken from Refs. (<sup>147</sup>Gd), (<sup>148</sup>Gd), (<sup>149</sup>Gd), (<sup>150</sup>Gd, <sup>151</sup>Tb, <sup>152,153</sup>Dy), and (<sup>150</sup>Gd(4a,4b)). The present analysis has an advantage because the sign of $`c`$ can be compared with the staggering significance $`Y`$ found in Ref. . According to this work, the significance is equal to the mean staggering amplitude divided by its uncertainty. It is highly unlikely that all the bands with the significance $`Y>2`$ exhibit the $`\mathrm{\Delta }I=4`$ bifurcation only because of statistical fluctuations in the $`\gamma `$-ray energy measurements. In particular, the independent measurements of the <sup>149</sup>Gd(1) staggering conclusively demonstrate the existence of the effect. We use this band as a reference for the single-particle structures of all the bands studied in this work.
Table II presents the bands without calculated deformation. In their analysis we have used the observation of Ref. that the filling of any particular orbital always induced the same deformation change in different nuclei. Subsequently this feature has been explained by the additivity of quadrupole and hexadecapole moments for SD bands in the $`A`$150 mass region . In a similar way as in the cited works we find the deformation changes $`\mathrm{\Delta }\epsilon `$ and $`\mathrm{\Delta }\epsilon _4`$, induced by a nucleon in the given state. The corresponding values are presented in Table III for the two rotational frequencies. They are used to evaluate the parameters $`\epsilon `$ and $`\epsilon _4`$ of the bands in Table II. The bands <sup>150</sup>Gd(6a,6b) are not given in the last table because the deformation changes induced by the orbital $`\nu [514]9/2`$ are not known.
Tables I and II help to understand which property of the single-particle structure is responsible for the $`\mathrm{\Delta }I=4`$ bifurcation. First of all we would like to emphasize that the necessary condition (4) is not violated in either of the bands with the known staggering significance. This is not a trivial fact because of the double cancellations in the expression $`Q_{44}`$: the partial cancellation of the $`q_{44}(n\alpha )`$ values with different signatures and the partial cancellation of the quantities $`Q_{44}(\pi )`$ and $`Q_{44}(\nu )`$ for almost all these bands. As a direct consequence of these cancellations, the value $`Q_{44}`$ for some bands with the small significance $`Y`$ changes sign and turns negative for high rotational frequencies. Besides the deformed shell model potential, the zero point of $`Q_{44}`$ depends also on the frequency renormalization factor, for which we take the conventional value 1.27. With such scaling the criterion (4) seems to be not reliable for small frequencies. Thus we use the high frequencies ($`\omega =`$ 0.6 and 0.8 MeV) to compare the staggering criterion with the experimental significance.
While Tables I and II exhibit definitely the correlations between the sign of the parameter $`c`$
and the significance $`Y`$, they also show some discrepancies. The high positive value of $`Q_{44}`$ in the bands <sup>147</sup>Gd(2) and <sup>148</sup>Gd(4) is the consequence of the neutron hole in the state $`\nu [642]5/2(\alpha =1/2)`$, which has according to Table IV the large negative $`q_{44}`$. The same effect produces the orbital $`\pi 6_3`$ in the bands <sup>147</sup>Eu(3) and <sup>150</sup>Gd(8a). The discrepancies observed in the bands <sup>147</sup>Gd(3) and <sup>148</sup>Gd(3) are less evident. Among the bands under study only these bands have the empty state $`\nu 7_1`$. It is possible that the first intruder plays a crucial role in the phenomenon (let us recall that the criterion (4) is only necessary). This tentative conclusion is confirmed by the nonstaggering bands <sup>150</sup>Gd(1,2), <sup>151</sup>Tb(1), and <sup>152</sup>Dy(1), but it disagrees with the staggering bands <sup>148</sup>Gd(5) and <sup>151</sup>Gd(1a). The first intruder is blocked up by the second one, $`\nu 7_2`$, in these bands (see also Ref. ). Let us note also that the first proton intruder $`\pi 6_1`$ is blocked up in all the bands under study.
From a strictly logical point of view, better test of the inequality (4) is provided by the pairs of the bands with configurations that differ by a single nucleon occupying an active or inactive orbitals. The filling of the inactive orbital $`\pi [301]1/2(\alpha =1/2)`$ does not change $`Q_{44}`$. Therefore any pair of the identical bands <sup>147</sup>Eu(1)/<sup>148</sup>Gd(1), <sup>148</sup>Eu(1)/<sup>149</sup>Gd(1) has identical staggering properties. The same is true for the identical bands <sup>147</sup>Gd(4)/<sup>148</sup>Gd(1) and <sup>148</sup>Gd(6)/<sup>149</sup>Gd(1), which configurations are distinguished by a neutron in the state $`[411]1/2(\alpha =1/2)`$. Identical staggering properties have the pair of bands <sup>150</sup>Gd(2)/<sup>152</sup>Dy(1) differing in two protons in the state $`[301]1/2`$.
This finding explains the observation of the staggering effect in identical SD bands . An exception is the band <sup>148</sup>Gd(5), which exhibits clear evidence for staggering. Its configuration is the same as those for the bands <sup>150</sup>Gd(2) or <sup>152</sup>Dy(1) apart from two neutron holes in the state $`[411]1/2`$ or the two $`[411]1/2`$ neutron and two $`[301]1/2`$ proton holes correspondingly. Nevertheless, statistically significant staggering has not been observed in the latter bands. One would suppose that the superposition principle does not work in this case. This suggestion is confirmed by the large nonaxial deformation of <sup>148</sup>Gd(5) found in the calculations of Ref. .
The active orbitals give us a more rigorous verification of the theory. A nucleon occupying this state adds significantly to the quantity $`Q_{44}`$ and may change its sign. Table IV shows some active orbitals involving in the configurations of almost all the studied bands and the estimated values $`q_{44}`$ for them. The orbital $`\nu [651]1/2(\alpha =1/2)`$ is one such example. Starting with the staggering bands <sup>148</sup>Eu(1), <sup>148</sup>Gd(6), <sup>149</sup>Gd(1) and removing a neutron from this orbital we get correspondingly the bands <sup>147</sup>Eu(1), <sup>147</sup>Gd(4), <sup>148</sup>Gd(1), which do not stagger. Thus this active orbital explains the remarkable property of the $`\mathrm{\Delta }I=4`$ bifurcation observed in Ref. .
In the next step we consider the signature partner bands based on the state $`\nu [402]5/2`$, which are associated with the generation of identical bands. This active orbital has the reasonably large values of the moment $`q_{44}`$ to modify the inequality (4). Consequently a pair of identical bands may have different staggering properties. The example is the band <sup>150</sup>Gd(4b), which is identical to <sup>149</sup>Gd(1) but does not exhibit staggering because the state $`\nu [402]5/2(\alpha =1/2)`$ has the large negative value $`q_{44}`$. Its signature partner, <sup>150</sup>Gd(4b), should stagger. Other examples of the signature partner bands involving this state are shown in Table I and II.
We extend now this procedure to the bands involving the configurations that differ by an arbitrary number of particles and holes in active and inactive orbitals. For the fixed rotational frequency, the $`Q_{44}`$ values of the two bands A and B are connected by the equality
$$Q_{44}(\mathrm{A})=Q_{44}(\mathrm{B})+\delta Q_{44}+\delta Q_{44}^{\mathrm{def}},$$
(11)
where $`\delta Q_{44}`$ is the contribution resulting from difference in active orbitals, while $`\delta Q_{44}^{\mathrm{def}}`$ represents the contribution due to the deformation change, which induced both active and inactive orbitals. According to additivity of multipole moments, the former quantity can be written as
$$\delta Q_{44}=\underset{\lambda }{}q_{44}(\lambda ),$$
(12)
where $`\lambda `$ runs over the active particle and/or hole states, which define the intrinsic configuration of the band A with respect to the band B (the reference band).
Since the contributions $`\delta Q_{44}`$ and $`\delta Q_{44}^{\mathrm{def}}`$ may be comparable, we have used the values $`Q_{44}`$ listed in Tables I and II to evaluate the relative nonaxial moment of active orbitals
$$\mathrm{\Delta }Q_{44}=\delta Q_{44}+\delta Q_{44}^{\mathrm{def}}=Q_{44}(\mathrm{A})Q_{44}(\mathrm{B}).$$
(13)
These quantities along with the staggering significances $`Y_\mathrm{A}`$ and $`Y_\mathrm{B}`$ allow us to get the more sophisticated test of the staggering criterion.
We have first selected twelve bands having the proper staggering significances to deal with the sample involving staggering ($`Y1.8`$) and nonstaggering ($`Y0.25`$) bands with a reasonable high likelihood. According to Eqs. (4) and (10), the former are characterized by the value $`Q_{44}>0`$ and the latter have $`Q_{44}<0`$. To compare the staggering properties of the bands A and B, we consider the two strong inequalities
$$Q_{44}(\mathrm{B})+\mathrm{\Delta }Q_{44}>0,\mathrm{if}Y_\mathrm{B}1.8,\mathrm{\Delta }Q_{44}>0,$$
$$Q_{44}(\mathrm{B})+\mathrm{\Delta }Q_{44}<0,\mathrm{if}Y_\mathrm{B}0.25,\mathrm{\Delta }Q_{44}<0.$$
(14)
The transformation $`\mathrm{\Delta }Q_{44}\mathrm{\Delta }Q_{44}`$ makes the sign of the sum $`Q_{44}(\mathrm{B})+\mathrm{\Delta }Q_{44}`$ in the inequalities (14) indefinite unless the absolute value of $`\mathrm{\Delta }Q_{44}`$ is small compared to $`Q_{44}(\mathrm{B})`$. Let us consider, for example, the bands <sup>148</sup>Gd(1) with $`Y_\mathrm{A}=0.23`$ and <sup>147</sup>Gd(3) with $`Y_\mathrm{B}=0.25`$, for which $`\mathrm{\Delta }Q_{44}=0.71`$. According to Eq. (11) and the second inequality (14), we have $`Q_{44}(\mathrm{A})<0`$ that is in agreement with absence of staggering in the band <sup>148</sup>Gd(1). On the other hand, considering <sup>148</sup>Gd(1) as the reference band, we cannot find the staggering behavior of the band <sup>147</sup>Gd(3) because the sign of the right hand side of Eq. (11) is indefinite.
The result of such comparison for 132 pairs of bands is presented in Table V. The columns of this table involve the reference bands B, whereas lines represent the bands A. The symbol $`+`$ ($``$) means that Eq. (11) and the inequalities (14) determine the staggering behavior of the band A correctly (incorrectly). A blank space is used when the sign of $`Q_{44}(\mathrm{A})`$ is indefinite and its comparison with the significance $`Y_\mathrm{A}`$ is impossible. The three groups of bands are clearly visible in Table V. (i) The nonstaggering bands <sup>147</sup>Gd(4), <sup>148</sup>Gd(1), <sup>151</sup>Gd(1b). There is no contradiction in the staggering behavior inside this group of bands. Such contradiction has not been found also between these bands and the bands of other groups. (ii) The four bands with clear evidence of staggering <sup>148</sup>Eu(1), <sup>148</sup>Gd(6), <sup>149</sup>Gd(1), and <sup>151</sup>Gd(1a). Whether or not the staggering behavior of the last band contradicts with that of the band <sup>148</sup>Eu(1) or <sup>149</sup>Gd(1) is not clear. (iii) The most striking feature of Table V is the third group of bands, which contradict with all the staggering bands of the second group. The bands <sup>147</sup>Gd(3), <sup>148</sup>Gd(3) with the empty first intruder $`\nu 7_1`$ and the band <sup>150</sup>Gd(2) with the blocked first intruder belong to this group. It should be noted that the band <sup>148</sup>Gd(5), being included in the sample, contradicts with the first and third groups.
We have presented a systematic study of the $`\mathrm{\Delta }I=4`$ bifurcation in the SD bands of the $`A150`$ mass region. The analysis is based on the necessary condition of the staggering phenomenon obtained in previous works and the MO potential in the pure single-particle approach. The results show that the criterion is working surprisingly well and is in a reasonably agreement with the statistical analysis of Haslip et al. We also have revealed the set of the single-particle states with a large nonaxial hexadecapole moment (active orbitals). They clearly highlight the role of the hexadecapole degree of freedom in the phenomenon and allows to answer the main question of why the staggering is not a universal feature of SD bands. Another class of states (inactive orbitals) allows to explain the observation of staggering in the identical bands <sup>148</sup>Eu(1), <sup>148</sup>Gd(6), and <sup>149</sup>Gd(1), which are the only bands clearly exhibiting the effect. Some discrepancies and contradictions established by the analysis indicate the need for further experimental and theoretical study of this interesting phenomenon.
The authors wish to thank Anatoli Afanasjev, Ingemar Ragnarsson for supplying the information concerning the structure and deformations of some bands, and Duncan Appelbe for giving the <sup>150,151</sup>Gd band configurations.
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# The Universe behind the Milky Way
## 1. Introduction
The unveiling of the galaxy distribution behind the Milky Way has turned into a research field of its own in the last ten years (e.g. “Unveiling Large-Scale Structures behind the Milky Way”, ASP Conf. Ser. 67, eds. Balkowski & Kraan-Korteweg, 1994, and “Mapping the Hidden Universe”, ASP Conf. Ser., eds. Kraan-Korteweg et al. 2000, in press). Why is it of interest to know the galaxy distribution behind the Milky Way, and why is it not sufficient to study galaxies and their large-scale distribution away from the foreground “pollution” of the Milky Way? To understand the dynamics in the nearby Universe and answer the question whether the the dipole in the Cosmic Microwave Background (CMB) and other velocity flow fields (e.g. towards the Great Attractor) can be fully explained by the clumpy galaxy/mass distribution, whole-sky coverage is essential. The lack of data in large areas of the sky – where the size of the gap due to the Milky Way depends on the wavelength at which galaxies are sampled – constitutes a severe restriction in solving these questions.
Based on various dedicated observational programs, using nearly all the bands of the electromagnetic spectrum, and the charting of large-scale structures and flow fields in the Zone of Avoidance (ZOA) from statistical reconstructions, a lot of progress has been achieved. In this review, we will give a status report on all of the above approaches. After a historic perspective on the ZOA (Sect. 2.), the cosmological questions for which the unveiling of the ZOA are most relevant are described (Sect. 3.). In Sect. 4., a discussion on the current knowledge of the foreground extinction is presented. This is followed by a description of the various observational multi-wavelength techniques that are currently being employed to uncover the galaxy distribution in the ZOA such as deep optical searches (Sect. 5.), near-infrared and far-infrared surveys (Sect. 6., and 7.), systematic blind radio surveys (Sect. 8.) and searches for hidden massive X-ray clusters (Sect. 9.). For each method, the different limitations and selection effects and results are presented. The various statistical reconstruction methods are reviewed and results of the density field in the ZOA compared with observational data (Sect. 10.).
## 2. The Zone of Avoidance
A first reference to the Zone of Avoidance (ZOA), or the “Zone of few Nebulae” was made in 1878 by Proctor, based on the distribution of nebulae in the “General Catalogue of Nebulae” by Sir John Herschel (1864). This zone becomes considerably more prominent in the distribution of nebulae presented by Charlier (1922) using data from the “New General Catalogue” by Dreyer (1888, 1895). These data also reveal first indications of large-scale structure: the nebulae display a very clumpy distribution. Currently well-known galaxy clusters such as Virgo, Fornax, Perseus, Pisces and Coma are easily recognizable even though Dreyer’s catalog contains both Galactic and extragalactic objects as it was not known then that the majority of the nebulae actually are external stellar systems similar to the Milky Way. Even more obvious in this distribution, though, is the absence of galaxies around the Galactic Equator. As extinction was poorly known at that time, no connection was made between the Milky Way and the “Zone of few Nebulae”.
A first definition of the ZOA was proposed by Shapley (1961) as the region delimited by “the isopleth of five galaxies per square degree from the Lick and Harvard surveys” (compared to a mean of 54 gal./sq.deg. found in unobscured regions by Shane & Wirtanen, 1967). This “Zone of Avoidance” used to be “avoided” by astronomers interested in the extragalactic sky because of the lack of data in that area of the sky and the inherent difficulties in analyzing the few obscured galaxies known there.
Merging data from more recent galaxy catalogs, i.e. the Uppsala General Catalog UGC (Nilson 1973) for the north ($`\delta 2.^{}5`$), the ESO Uppsala Catalog (Lauberts 1982) for the south ($`\delta 17.^{}5`$), and the Morphological Catalog of Galaxies MCG (Vorontsov-Velyaminov & Archipova 1963-74) for the strip inbetween ($`17.^{}5<\delta <2.^{}5`$), a whole-sky galaxy catalog can be defined. To homogenize the data determined by different groups from different survey material, the following adjustments have to be applied to the diameters: $`D=1.15D_{\mathrm{UGC}},D=0.96D_{\mathrm{ESO}}`$ and $`D=1.29D_{\mathrm{MCG}}`$ (see Fouqué & Paturel 1985, Lahav 1987). According to Hudson & Lynden-Bell (1991) this “whole-sky” catalog then is complete for galaxies larger than $`D=1.^{}3`$.
The distribution of these galaxies is displayed in Galactic coordinates in Fig. 1 in an equal-area Aitoff projection centered on the Galactic Bulge ($`\mathrm{}=0{}_{}{}^{},b=0^{}`$). The galaxies are diameter-coded, so that structures relevant for the dynamics in the local Universe stand out accordingly. Figure 1 clearly displays the irregularity in the distribution of galaxies in the nearby Universe such as the Local Supercluster visible as a great circle (the Supergalactic Plane) centered on the Virgo cluster at $`\mathrm{}=284{}_{}{}^{},b=74^{}`$, the Perseus-Pisces chain (PP) bending into the ZOA at $`\mathrm{}=95^{}`$ and $`\mathrm{}=165^{}`$, the general overdensity in the Cosmic Microwave Background dipole direction ($`\mathrm{}=276{}_{}{}^{},b=30^{}`$; Kogut et al. 1993) and the general galaxy overdensity in the Great Attractor region (GA) centered on $`\mathrm{}=320{}_{}{}^{},b=0^{}`$ (Kolatt et al. 1995) with the Hydra ($`270{}_{}{}^{},27^{}`$), Antlia ($`273{}_{}{}^{},19^{}`$), Centaurus ($`302{}_{}{}^{},22^{}`$) and Pavo ($`332{}_{}{}^{},24^{}`$) clusters.
Most conspicuous in this distribution is, however, the very broad, nearly empty band of about 20 width. As optical galaxy catalogs are limited to the largest galaxies they become increasingly incomplete close to the Galactic Equator where the dust thickens. This diminishes the light emission of the galaxies and reduces their visible extent. Such obscured galaxies are not included in diameter- or magnitude-limited catalogs because they appear small and faint – even though they might be intrinsically large and bright. A further complication is the growing number of foreground stars close to the Galactic Plane (GP) which fully or partially block the view of galaxy images.
Comparing this “band of few galaxies” with the currently available 100$`\mu `$m dust extinction maps of the DIRBE experiment (Schlegel et al. 1998; see Sect. 4.), we can see that the ZOA – the area where the galaxy counts become severely incomplete – is described almost perfectly by the absorption contour in the blue $`A_B`$ of $`1.^\mathrm{m}0`$ (where $`A_B=4.14E(BV)`$; Cardelli et al. 1989). This contour matches the by Shapley (1961) defined ZOA closely.
## 3. The ZOA as an obstacle for cosmological studies
### 3.1. Large-scale structures
Enormous effort and observation time has been devoted in the last decades to establish the extragalactic large-scale structure. From the various redshift slices, 3-dimensional pictures evolve that distribute galaxies predominantly in clusters, sheets and filaments, leaving large areas devoid of luminous matter. These filaments, respectively their sizes, carry information on the conditions and formation processes of the early Universe, providing important constraints which must be reproduced in cosmological models.
Many of the known nearby large-scale structures are, however, bisected by the Galactic Plane - such as the Local Supercluster, the Perseus-Pisces chain, and the Great Attractor (see Fig. 1). What is their true extent and their mass? It is curious that the two major superclusters in the local Universe, i.e. Perseus-Pisces and the Great Attractor overdensity, lie at similar distances on opposite sides of the Local Group (LG) – both partially obscured by the ZOA. Which one of the two is dominant in the tug-of-war on the LG? Do these features continue across the Galactic Plane and are there other massive structures hidden in the ZOA for which so far no indication exists?
What is the size of the largest coherent structures? Are, for instance, the Great Wall and the Perseus-Pisces chain connected across the ZOA as suggested in 1982 by Giovanelli & Haynes (but see Marzke et al. 1996), indicating structures of over $`200h_{50}^1`$ Mpc. The latter would be incompatible with the angular extent over which fluctuations – the seeds of current large-scale structures – have been measured in the CMB. To answer these questions, superclusters need to be fully mapped across the ZOA.
### 3.2. Dipole determinations
Filling the ZOA also is paramount with respect to the evaluation of the peculiar velocity of the Local Group. The dipole anisotropy in the CMB radiation is explained by a peculiar motion of the LG relative to the CMB of $`\stackrel{}{v}_p=627`$ km s<sup>-1</sup> towards $`\mathrm{}=276{}_{}{}^{},b=30^{}`$ (Kogut et al. 1993). This motion arises from the net gravitational attraction on the LG due to the irregular distribution of matter in the Universe. Reproducing the vector measured in the CMB radiation with the LG motion determined from the matter distribution (direction as well as convergence distance) will lead to constraints on the cosmological parameter $`\mathrm{\Omega }_0`$.
The determination of the gravity field at the position of the LG, i.e. velocity and direction of the peculiar motion, requires whole-sky coverage. Kolatt et al. (1995) have shown, for instance, that the mass distribution within $`\pm 20^{}`$ of the ZOA – as derived from theoretical reconstructions of the density field (see Sect. 10.) – is crucial to the derivation of the gravitational acceleration of the LG: the direction of the motion measured within a volume of 6000 km s<sup>-1</sup> will change by $`31^{}`$ when the (reconstructed) mass within the ZOA is included. The results derived so far for the apex of the LG motion, as well as the distance at which convergence is attained, still are controversial. The lack of data for the ZOA remains one of the main uncertainty in current dipole determinations (e.g. Rowan-Robinson et al. 2000).
Most dipole determinations have assumed a uniformly filled ZOA (e.g. by Poissonian statistics) or have used cloning methods which transplant the fairly well-mapped regions adjacent to the ZOA into the ZOA, or a spherical harmonic analysis (Sect. 10.). All procedures are unsatisfactory, because inhomogenous data coverage, incorrect assumptions on the galaxy distribution in the ZOA, and/or false assumptions on the ZOA mask to be filled will introduce nonexisting flow fields. Care should therefore be taken on how to extrapolate the galaxy density field across the ZOA. Obviously, a reliable consensus on the galaxy distribution in the ZOA is important to minimize these uncertainties.
In this context, not only the identification of unknown and suspected clusters, filaments and voids are relevant, but also the detection of nearby smaller entities. In linear theory, the peculiar velocity of the LG is propertional to the net gravity field which can be determined from the sum of the masses of all galaxies divided by the distance squared:
$$\stackrel{}{v}_p\frac{\mathrm{\Omega }_0^{0.6}}{b}\frac{_i}{r_i^2}\widehat{𝐫_𝐢},$$
where $`\mathrm{\Omega }_0`$ is the density parameter and $`b`$ the bias parameter. Since gravity as well as the flux of a galaxy decrease with $`r^2`$, the direction and amplitude of the peculiar velocity can be determined directly from, for instance, the sum of the apparent magnitudes of the galaxies in the sky under the assumption of constant mass-to-light ratio
$$\stackrel{}{v}_p\underset{\mathrm{i}}{}10^{0.4\mathrm{m}}\widehat{𝐫_𝐢}.$$
This has important implications and suggests, for instance, that the galaxy CenA with an absorption-corrected magnitude of $`B^o=6.^\mathrm{m}1`$ exerts a stronger luminosity-indicated gravitational attraction on the Local Group than the whole Virgo cluster. The problem whether galaxies trace the mass is inherent to all cumulative dipole determinations. These calculations also predict that the 8 apparently brightest galaxies – which are all nearby ($`v<300`$ km s<sup>-1</sup>) – are responsible for 20% of the total dipole as determined from optically known galaxies within $`v<6000`$ km s<sup>-1</sup>. Hence, a major part of the peculiar motion of the LG is generated by a few average, but nearby galaxies. Note, however, that for nearby objects non-linear dynamics has to be taken into account.
In this sense, the detection of nearby galaxies or galaxy groups hidden by the obscuration layer of the Galaxy can be as important as the detection of entire clusters at larger distances. The expectation of finding additional nearby galaxies in the ZOA is not unrealistic. Six of the nine apparently brightest galaxies (extinction-corrected) are located in the ZOA: IC342, Maffei 1, Maffei 2, NGC4945, CenA and the recently discovered galaxy Dwingeloo 1 (Kraan-Korteweg et al. 1994b). In the absence of Galactic extinction both Maffei 1 and IC 342 would subtend angles as large as the full Moon (McCall & Buta 1996).
### 3.3. Dynamics of the Local Group
It is commonly believed that the Local Group (LG) of galaxies is dominated by the Milky Way and Andromeda (M31). A discovery of an unknown Andromeda-like galaxy behind the Milky Way may therefore dramatically change our understanding of the LG dynamics.
Kahn & Woltjer (1959) and Lynden-Bell (1982) have shown that the distances and motions of nearby galaxies may be used to constrain the total mass of the pair Milky Way and M31 and the dynamical age of the Local Group. The method, known as “Local Group timing”, is based upon the assumption that M31 and the Milky Way separated from the Hubble flow soon after formation and behaved dynamically like an isolated binary ever since, with the consequence that the Milky Way has turned in its orbit and is now falling back toward M31. However, it seems that the dominant members of the IC 342/Maffei Group may be massive enough and near enough to the Local Group to have had an influence on its dynamical history (McCall 1989; Valtonen et al. 1993), thereby calling into question the binary hypothesis of LG timing. Peebles (1990) and Dunn & Laflamme (1995) extended this approach by tracing the orbits of the LG galaxies back in time, under the “least action principle”. These analyses rely heavily on having a full census of LG galaxies and on accurate distances to them. Therefore, better mapping of the ZOA is important in two ways: (a) discovering new members of the LG and neighbouring groups (or ruling out their existence); (b) measuring the intrinsic properties of galaxies, properly corrected for extinction, to accurately estimate their distances.
### 3.4. Cosmic flow fields
Density enhancements locally decelerate the uniform expansion field, as observed within our own Local Supercluster, resulting in systematic streaming motions over and above the uniform expansion field. These flows, on the other hand, can be exploited to map the mass density field independent of the galaxy distribution and/or an assumption on the mass-to-light ratio using peculiar velocities of galaxies, $`\stackrel{}{v}_p=\stackrel{}{v}_{\mathrm{obs}}\stackrel{}{v}_\mathrm{H}`$, for which distance determinations independent of redshift are available. The latter can be obtained via the Tully-Fisher relation for spiral galaxies (Tully and Fisher 1977) or the $`D_n\sigma `$ relation for elliptical galaxies (Lynden-Bell et al. 1988). Note though, that only the radial component of the peculiar motion of a galaxy can be measured. The reconstruction of such potential fields and density fields have the advantage that they can locate hidden mass overdensities even if “unseen”. These methods are therefore of particular interest for ZOA research, as these potential fields provide information on the mass distribution behind the Milky Way without having access to the real data on the galaxies hidden there.
Based on these considerations, Dressler et al. (1987) interpreted a systematic infall pattern from the peculiar velocities of about 400 elliptical galaxies as being due to a hypothetical Great Attractor with a mass of $`5\times 10^{16}_{}`$ at a position in redshift space of $`(\mathrm{},b,v)=(307{}_{}{}^{},9{}_{}{}^{},4500`$ km s<sup>-1</sup>) (Lynden-Bell et al. 1988). A more recent study by Kolatt et al. (1995), based on a larger data set (elliptical and spiral galaxies) and the potential reconstruction method POTENT (see Sect. 10. and Fig. 18) place the center of the GA right behind the Milky Way. Recent consensus is that the GA is an extended region ($`40{}_{}{}^{}\times 40^{}`$) of moderately enhanced galaxy density. Although there is a considerable excess of optical galaxies and IRAS-selected galaxies in this region (see Fig. 1 and Fig. 8), no dominant cluster or central peak can been seen. But the central part of the GA is hidden by the Milky Way.
Large surveys (e.g. Mathewson et al. 1992 for field galaxies, Mould et al. 1991, Han 1992 for clusters) have resulted in a large collection of peculiar velocities of galaxies, put together in the Mark III catalog of peculiar velocities (Willick et al. 1997). A large effort to obtain peculiar motions of a sample of uniform sky coverage (excluding, however, the ZOA at $`|b|10^{}`$) is presented by Giovanelli et al. (1997) and da Costa et al. (1996). This has opened up a new field in cosmology, namely the dynamics of cosmic flows or large-scale dynamics.
## 4. Extinction by the Milky Way
A crucial step in exploring the extragalactic sky behind the ZOA is a detailed understanding of our own Galaxy, since it would otherwise remain impossible to disentangle newly unveiled clustering from the patchiness of the foreground extinction and, moreover, to correct the observed parameters of detected galaxies for the absorption effects. This requires (a) a high-resolution, well-calibrated map of the foreground extinction and (b) a clear understanding how Galactic extinction affects the observed galaxy parameters. Cameron (1990) investigated the latter in the optical by artificially obscuring high-latitude galaxies, an approach which really needs to be refined and to be explored at other wavelengths.
The Galactic foreground extinction is a function of wavelength. According to Cardelli et al. (1989), the mean extinction at a given wavelength, $`A_\lambda `$, compared to the visual extinction, $`A_V`$, is:
$$A_\lambda /A_V=a(1/\lambda )+b(1/\lambda )/R_V$$
(see Mathis 1990, for a comprehensive overview on the interstellar dust in the Galaxy). The interstellar extinction is a function of the ratio of total to selective extinction, e.g. $`R_V`$ \[$`A_V/E(BV)`$\] and the actual value depends on the environment along the line of sight through the Galaxy. For the diffuse interstellar medium a standard value of $`R_V=3.1`$ applies, whereas a higher value of $`R_V`$ is evident for lines of sight through dense molecular clouds ($`R4`$). But this ratio of total to selective extinction generally is founded upon observations of stars in the Milky Way, ignoring shifts in effective wavelengths which are known to depend upon reddening and intrinsic color. For galaxies heavily reddened by dust in the Milky Way, this approach can lead to significant errors in both magnitudes and distances (McCall & Armour 2000). So the values for $`R_V`$ are by no means generally applicable.
For years, the extinction maps by Burstein & Heiles (1982) were the standard. But they do not cover the ZOA ($`|b|>10^{}`$). These maps can, however, be extrapolated towards lower latitudes, following the precepts of Burstein & Heiles (1978, 1982) using the Galactic H I column densities $`N_{\mathrm{HI}}`$ (e.g. from Kerr et al. 1986 for the south, and Hartmann & Burton 1997 for the north) assuming a constant gas-to-dust ratio:
$$E(BV)=\left(\frac{N_{\mathrm{HI}}}{2.2310^{18}}\right)\times 4.4310^40.055.$$
However, the gas-to-dust ratio does vary. Burstein et al. (1987) report, for instance, an increase in the gas-to-dust ratio for a region in the southern Milky Way ($`230{}_{}{}^{}\mathrm{}310^{}`$ and $`20{}_{}{}^{}b20^{}`$) of up to a factor of 2, implying severe overstimates of the extinction. Moreover, close to the Galactic Plane ($`|b|<2^{}`$), the Galactic H I line might be saturated, leading to an underestimate of the true extinction. At these latitudes though, the Galactic CO (Dame et al. 1987) can be used as a tracer of extinction.
The recently published 100$`\mu `$m extinction maps from the DIRBE experiment (Schlegel et al. 1998) give an improved estimate of the foreground extinction because they provide a direct measure of the dust column density, and because these maps have better angular resolution (6$`.^{}`$1 compared to $`2030^{}`$ of the H I maps). According to Schlegel et al., they are a factor 2 better at low and moderate extinction compared to the Burstein & Heiles maps. However, as stated by Schlegel et al. themselves, the accuracy of the DIRBE maps still needs to be established at low Galactic latitudes ($`|b|10^{}`$). Woudt (1998) found – from photometry and measurements of the Mg<sub>2</sub>-index of 18 early type galaxies in the ZOA – that the extinction for moderate to high DIRBE reddenings is systematically underestimated (by a factor of f=0.86). As his new calibration so far is based only on few galaxies in a small region of the ZOA, it seems too early to incorporate these adjustments to the DIRBE maps. However, this study clearly illustrates the need for a careful calibration of the DIRBE maps in the ZOA.
Although the knowledge on the Galactic foreground extinction has improved enormously in the last 20 years, many open points still need to be resolved before we can properly analyze the galaxy distribution uncovered in the ZOA. On the other hand, a lot can be learned about total extinction from ZOA research itself. As mentioned above, Woudt (1998) has used photometry and spectroscopy of early-type galaxies in the ZOA to obtain a first calibration of the DIRBE maps in the ZOA. He now has started a program to pursue this systematically for the southern Milky Way using a sample of about 300 early-type galaxies distributed within the southern ZOA. Furthermore, Saito et al. (2000) have used $`BI`$ colors of over a hundred galaxy candidates at extremely low latitudes ($`|b|<1^{}`$) in combination with CO emission to determine extinction estimates, and Temporin et al. (2000) determined total extinctions from $`BVRI`$ photometry towards galaxies in the region $`29{}_{}{}^{}<\mathrm{}<14{}_{}{}^{},|b|<11^{}`$. Similarly, the colors of galaxies identified in the ZOA with the near infrared surveys DENIS and 2MASS (see Sect. 6.) will provide a huge data base to calibrate the DIRBE maps at low Galactic latitudes.
## 5. Optical surveys
Systematic optical galaxy catalogs are generally limited to the largest galaxies (typically with diameters D $`>1^{}`$, e.g. Lauberts 1982). These catalogs become, however, increasingly incomplete as the dust thickens, creating a “Zone of Avoidance” in the distribution of galaxies of roughly 25% of the sky. Systematically deeper searches for partially obscured galaxies – down to fainter magnitudes and smaller dimensions compared to existing catalogs – have been performed on existing sky surveys with the aim of reducing this ZOA. Meanwhile, through the efforts of various collaborations, nearly the whole ZOA has been surveyed and over 50000 previously unknown galaxies were discovered in this way. These surveys are not biased with respect to any particular morphological type and were able to identify important new large-scale structures in and across the Milky Way.
### 5.1. Early searches and results
One of the first attempts to detect galaxies in the ZOA was carried out by Böhm-Vitense in 1956. She did follow-up observations in selected fields in the GP in which Shane and Wirtanen (1954) found objects that ”looked like extragalactic nebulae” but were not believed to be galaxies because they were so close to the dust equator. She confirmed many galaxies and concluded that the obscuring matter in the GP must be extremely thin and full of holes between $`\mathrm{}=125^{}`$ and $`130^{}`$.
Because extinction was known to be low in Puppis, Fitzgerald (1974) performed a galaxy search in one field there ($`\mathrm{}245^{}`$) and discovered 18 small and faint galaxies. Two years later Dodd & Brand (1976) examined 3 fields adjacent to this area ($`\mathrm{}243^{}`$) and detected another 29 galaxies. Kraan-Korteweg & Huchtmeier (1992) observed these galaxies in H I at Effelsberg and identified an unknown nearby cluster at ($`\mathrm{},b,v)=(245{}_{}{}^{},0{}_{}{}^{},1500`$ km s<sup>-1</sup>). Including IRAS data in the analysis of this cluster, it could be shown that its density is comparable to the Virgo cluster and that this Puppis cluster may contribute a significant component to the motion of the LG (Lahav et al. 1993).
During a search for infrared objects, Weinberger et al. (1976) detected two galaxy candidates near the Galactic Plane ($`\mathrm{}88^{}`$) which Huchra et al. confirmed in 1977 to be the brightest members of a galaxy cluster at 4200 km s<sup>-1</sup>. This discovery led Weinberger (1980) to start the first systematic galaxy search. Using the POSS E prints, he covered the whole northern GP ($`\mathrm{}=33{}_{}{}^{}213^{}`$) in a thin strip $`(|b|2{}_{}{}^{})`$. He found 207 galaxies, the distribution of which is highly irregular: large areas disclose no galaxies and the ”hole” pointed out by Böhm-Vitense was verified, but most conspicuous was a huge excess of galaxies around $`\mathrm{}=160{}_{}{}^{}165^{}`$. In 1984, Focardi et al. made the connection with large-scale structures: they interpreted the excess as the possible continuation of the Perseus-Pisces cluster across the GP to the cluster A569. Radio-redshift measurements by Hauschildt (1987) established that the PP cluster at a mean redshift of $`v=5500`$ km s<sup>-1</sup> extends to the cluster 3C129 in the GP ($`\mathrm{}=160{}_{}{}^{},b=0.^{}1`$). Additional H I and optical redshift measurements of Zwicky galaxies by Chamaraux et al. (1990) indicate that this chain can be followed even further to the A569 cloud at $`v6000`$ km s<sup>-1</sup> on the other side of the ZOA.
These early searches proved that large-scale structure can be traced to much lower Galactic latitudes despite the foreground obscuration and its patchy nature which suggests clustering in the galaxy distribution independent of large-scale structure. The above investigations did confirm suspected large-scale features across the GP through searches in selected regions and follow-up redshift observations. To study large-scale structure systematically broader latitude strips covering the whole Milky Way, respectively the whole ZOA (see Fig. 1) are required.
### 5.2. Status of systematic optical searches
Using existing sky surveys such as the first and second generation Palomar Observatory Sky Surveys POSS I and POSS II in the north, and the ESO/SRC (United Kingdom Science Research Council) Southern Sky Atlas, various groups have performed systematic deep searches for “partially obscured” galaxies, i.e. they catalogued galaxies down to fainter magnitudes and smaller dimensions ($`D>0.^{}1`$) than existing catalogs. Here, examination by eye remains the best technique. A separation of galaxy and star images can not be done as yet on a viable basis below $`|b|<10{}_{}{}^{}15^{}`$ by automated measuring machines such as e.g. COSMOS (Drinkwater et al. 1996) or APM (Lewis & Irwin 1996) and sophisticated extraction algorithms, nor with the application of Artificial Neural Networks \[ANN\]. The latter was tested by Naim (1995) who used ANN to identify galaxies with diameters above 25<sup>′′</sup> at low Galactic latitudes ($`b5^{}`$). Galaxies could be identified using this algorithm, and although an acceptable hit rate for galaxies of 80 – 96% could be attained when ANN was trained on high latitude fields, the false alarms were of equal order. Using low latitude fields as training examples, the false alarms could be reduced to nearly zero but then the hit rate was low ($``$ 30 - 40%). The first attempts of using ANN in the ZOA are encouraging but clearly need further development. So, although surveys by eye are both tiring and time consuming, and maybe not as objective, they currently still provide the best technique to identify partially obscured galaxies in crowded star fields.
Meanwhile, nearly the whole ZOA has been visually surveyed for galaxies. The various surveyed regions are displayed in Fig. 2. Details and results on the uncovered galaxy distributions and the respective references are described below:
Region A: In order to trace the possible continuity across the Galactic Plane of the southwestern spur of the PP complex Pantoja et al. (1994, 1997) have searched for galaxy candidates on 29 POSS E prints with a 12 $`x`$ magnification. They identified 1480 galaxy candidates in the region delimited by $`4^h\alpha 8^h,0{}_{}{}^{}\delta 37^{}`$, where the declination range was optimized for H I follow-up redshift observations with the 305 m Arecibo radio telescope.
Regions B<sub>1</sub>-B<sub>3</sub>: Using the first generation POSS I prints, and more recently also the deeper POSS II films, Weinberger and collaborators in Austria have expanded their optical galaxy searches in such a way that they now cover nearly the whole northern Milky Way. They also searched the Puppis region for comparative purposes. The information and data of the B<sub>1</sub> region from POSS I plates can be found in Seeberger et al. (1994), Seeberger et al. (1996), Lercher et al. (1996), Saurer et al. (1997), Seeberger et al. (1998), and Marchiotto et al. 1999 for the B<sub>2</sub> regions (also POSS I), and the B<sub>3</sub> region from POSS II fields in Weinberger et al. 1999). In total, they uncovered about 9500 galaxies in the northern ZOA. Their distribution shows a marked overdensity at the suspected connection of the PP supercluster across the Galactic Plane ($`ł165^{}`$). A comparison between the old and new generation POSS fields found a dramatic increase in galaxy numbers of about a factor of 3 for the deeper POSS II fields.
Regions C<sub>1</sub>-C<sub>3</sub>: Using the infrared film copies of the ESO/SRC survey Japanese groups led by Saito investigated the ZOA in the longitude range $`205<\mathrm{}<260{}_{}{}^{},|b|<10^{}`$ (C<sub>1</sub>: Saito et al. 1990, 1991). On 32 fields they charted over 7000 galaxies with $`D0.1`$ mm. As anticipated, a large number of galaxies was detected in the low opacity region of Puppis ($`\mathrm{}245^{}`$). Correlating the density with the H I-column density, they found indications of the existence of various clusters. They then continued studying regions close to the Galactic bulge, i.e. the Sagittarius/Galactic region (C<sub>2</sub>: $`7<\mathrm{}<16{}_{}{}^{},19<b<1^{}`$ by Roman et al. 1998), and the Aquila and Sagittarius region (C<sub>3</sub>: $`8<\mathrm{}<47{}_{}{}^{},|b|<17^{}`$ by Roman et al. 1996) uncovering a further 12500 galaxies in this very opaque region. This data set is less homogeneous as a variety of survey material had to be used to cover this area.
Regions D<sub>1</sub>-D<sub>5</sub>: Since 1988, various groups led by Kraan-Korteweg have searched the southern Milky Way between the Puppis region (C<sub>1</sub>) and the Galactic bulge region (C<sub>2</sub> & E) using the IIIaJ film copies of the ESO/SRC and a magnification of 50. The surveys are divided into the Hydra to Puppis region (D<sub>1</sub>: Salem & Kraan-Korteweg, in prep.), the Hydra/Antlia Supercluster region (D<sub>2</sub>: Kraan-Korteweg 2000), the Crux region (D<sub>3</sub>: Woudt 1998, Woudt & Kraan-Korteweg 2000a), the GA region (D<sub>4</sub>: Woudt 1998, Woudt & Kraan-Korteweg 2000b), and the Scorpius region (D<sub>5</sub>: Fairall & Kraan-Korteweg, 2000). Slightly over 17 000 galaxies were identified in these regions of which $`97\%`$ were previously unknown. Folding the galaxy distribution with extinction maps revealed various unknown clumpings, the most impressive recognition being that the cluster A3627 at ($`\mathrm{},b,v)=(325{}_{}{}^{},7{}_{}{}^{},4848`$ km s<sup>-1</sup>) within the Great Attractor region (see Fig. 3) would be the most prominent galaxy overdensity in the southern sky were it not for the diminishing effects of the foreground extinction (Kraan-Korteweg et al. 1996).
Region E: Motivated by the chance discovery of two clusters behind the Galactic bulge, i.e. the Ophiuchus cluster at $`(\mathrm{},b,v)=(0.^{}5,9.^{}0,8600`$ km s<sup>-1</sup>) by Johnston et al. (1981) and Wakamatsu & Malkan (1981), as well as the Sagittarius cluster closeby in redshift space ($`359.^{}8,8.^{}0,8400`$ km s<sup>-1</sup>), Wakamatsu and collaborators surveyed this region in more detail. They performed a deep survey of six ESO/SRC fields centered on these clusters and a shallow survey from $`16^h10^m<\alpha <17^h50^m,32.^{}5<\delta <0^{}`$ to search for wall-like connections with the Hercules cluster (region E). In the former region close to 4000 galaxies were charted with $`D>0.1`$ mm, revealing two new clusters and 4 galaxy groups all at the same redshift range (Wakamatsu et al. 1994, Hasegawa et al. 2000).
Region F: In 1995, Hau et al. (1995) searched 12 red POSS plates at $`\mathrm{}135^{}`$ because of the increased likelyhood of detecting galaxies along the Supergalactic Plane. They indeed identified a signficant number of galaxies (N = 2575), though this relatively high number is also due to the fact that this search includes higher latitude ($`|b|<25^{}`$) regions, hence lower extinction levels compared to the other searches. To confirm the nature of these galaxy candidates, follow-up observations using various techniques were performed on a sample of suspected nearby galaxies (Lahav et al. 1998).
A comparison of the surveyed regions (Fig. 2) with the ZOA as outlined in Fig. 1 clearly demonstrates that nearly the whole ZOA has been covered by systematic deep optical galaxy searches. All these searches have similar characteristics and reveal that galaxies can easily be traced through obscuration layers of 3 magnitudes, narrowing therewith the ZOA considerably. This is illustrated with Fig. 3 which shows an area of the sky centered on the southern Milky Way with all the Lauberts galaxies larger than $`D1.^{}3`$ (diameter-coded as in Fig. 1) plus all the galaxies with $`D12^{\prime \prime }`$ from the deep optical galaxy searches by Kraan-Korteweg and collaborators (D<sub>1</sub>-D<sub>5</sub> in Fig. 2). DIRBE extinction contours equivalent to $`A_B=1.^\mathrm{m}0`$ and $`3.^\mathrm{m}0`$ are also drawn.
A few galaxies still are recognizable up to extinction levels of $`A_B=5.^\mathrm{m}0`$ and a handful of very small galaxy candidates have been found at even higher extinction levels. The latter ones most likely indicate holes in the dust layer. Overall, the mean number density follows the dust distribution remarkably well. The contour level of $`A_B=5.^\mathrm{m}0`$, for instance, is nearly indistinguishable from the galaxy density contour at 0.5 galaxies per square degree.
Analyzing the galaxy density as a function of galaxy size, magnitude and/or morphology in combination with the foreground extinction has led to the identification of various important large-scale structures and their approximate distances. In Fig. 3, for instance, the most extreme overdensity is found at $`(\mathrm{},b)(325{}_{}{}^{},7^{}`$). It is at least a factor 10 denser compared to regions at similar extinction levels. This galaxy excess is centered on the cluster A3627, now recognized as the most massive cluster in the nearby Universe (Kraan-Korteweg et al. 1996, see also Sect. 5.4.). To trace these structures in detail, an understanding of the completeness of these searches is required and follow-up observations must be obtained to map the large-scale structures in redshift space (see Sect. 5.4.).
### 5.3. Completeness of optical galaxy searches
In order to merge the various deep optical ZOA surveys with existing galaxy catalogs, Kraan-Korteweg (2000) and Woudt (1998) have analyzed the completeness of their ZOA galaxy catalogs – the Hyd/Ant \[D<sub>2</sub>\], Crux \[D<sub>3</sub>\] and GA \[D<sub>4</sub>\] region – as a function of the foreground extinction.
By studying the apparent diameter distribution as a function of the extinction $`E(BV)`$ (Schlegel et al. 1998) as well as the location of the flattening in the slope of the cumulative diameter curves $`(\mathrm{log}D)(\mathrm{log}N)`$ for various extinction intervals (cf. Fig. 5 and 6 in Kraan-Korteweg 2000), they conclude that their optical ZOA surveys are complete to an apparent diameter of $`D=14^{\prime \prime }`$ – where the diameters correspond to an isophote of 24.5 mag/arcsec<sup>2</sup> – for extinction levels less than $`A_B=3.^\mathrm{m}0`$.
How about the intrinsic diameters, i.e. the diameters galaxies would have if they were unobscured? A spiral galaxy seen through an extinction of $`A_B=1.^\mathrm{m}0`$ will, for example, be reduced to $`80\%`$ of its unobscured size. Only $`22\%`$ of a (spiral) galaxy’s original dimension is seen when it is observed through $`A_B=3.^\mathrm{m}0`$. In 1990, Cameron derived analytical descriptions to correct for the obscuration effects by artificially absorbing the intensity profiles of unobscured galaxies. These corrections depend quite strongly on morphological type due to the mean surface brightness and difference in brightness profiles between early-type galaxies and spiral galaxies. Applying these corrections, Kraan-Korteweg (2000) and Woudt (1998) found that at $`A_B=3.^\mathrm{m}0`$, an obscured spiral or an elliptical galaxy at their apparent completeness limit of $`D=14^{\prime \prime }`$ would have an intrinsic diameter of $`D^o60^{\prime \prime }`$, respectively $`D^o50^{\prime \prime }`$. At extinction levels higher than $`A_B=3.^\mathrm{m}0`$, an elliptical galaxy with $`D^o=60^{\prime \prime }`$ would appear smaller than the completeness limit $`D=14^{\prime \prime }`$ and might have gone unnoticed. These optical galaxy catalogs should therefore be complete to $`D^o60^{\prime \prime }`$ for galaxies of all morphological types down to extinction levels of $`A_B3.^\mathrm{m}0`$ with the possible exception of extremely low-surface brightness galaxies. Only intrinsically very large and bright galaxies – particularly galaxies with high surface brightness – will be recovered in deeper extinction layers. This completeness limit could be confirmed by independently analyzing the diameter vs. extinction and the cumulative diameter diagrams for extinction-corrected diameters.
One can thus supplement the ESO, UGC and MCG catalogs – which are complete to $`D=1.^{}3`$ – with galaxies from optical ZOA galaxy searches that have $`D^o1.^{}3`$ and $`A_B3.^\mathrm{m}0`$. As the completeness limit of the optical searches lies well above the ESO, UGC and MGC catalogs, one can assume that the other similarly performed optical galaxy searches in the ZOA should also be complete for galaxies with extinction-corrected diameters $`D^o1.^{}3`$ to extinction levels of $`A_B3.^\mathrm{m}0`$. In Fig. 4, we have then taken the first step in arriving at an improved whole-sky galaxy distribution with a reduced ZOA. In this Aitoff projection, we plot all the UGC, ESO, MGC galaxies that have extinction-corrected diameters $`D^o1.^{}3`$ (remember that galaxies adjacent to the optical galaxy search regions are also affected by absorption though to a lesser extent: $`A_B1.^\mathrm{m}0`$), and added all the galaxies from the various optical surveys with $`D^o=1.^{}3`$ and $`A_B3.^\mathrm{m}0`$ for which positions and diameters were available. The regions for which these data are not yet available are marked in Fig. 4. As some searches were performed on older generation POSS I plates, which are less deep compared to the second generation POSS II and ESO/SERC plates, an additional correction was applied to those diameters, i.e. the same correction as for the UGC galaxies which also are based on POSS I survey material ($`D_{25}=1.15D_{POSSI}`$).
A comparison of Fig. 1 with Fig. 4 demonstrates convincingly how the deep optical galaxy searches realize a considerable reduction of the ZOA: we can now trace the large-scale structures in the nearby Universe to extinction levels of $`A_B=3.^\mathrm{m}0`$. Inspection of Fig. 4 reveals that the galaxy density enhancement in the GA region is even more pronounced and a connection of the Perseus-Pisces chain across the Milky Way at $`\mathrm{}=165^{}`$ more likely. Hence, these supplemented whole-sky maps certainly should improve our understanding of the velocity flow fields and the total gravitational attraction on the Local Group.
### 5.4. Redshift follow-ups of optical surveys
The analysis of the galaxy density as a function the foreground extinction revealed various large-scale structures in the ZOA. These need to be mapped in redshift space. So far, the Perseus-Pisces supercluster, the Puppis region, the Ophiuchus supercluster behind the Galactic Bulge area, and the southern ZOA have been intensively observed. The most prominent new galaxy structures revealed in this way are summarized below. Their approximate positions (ordered in Galactic latitude) are given as ($`\mathrm{},b,v`$):
Based on a redshift survey of over 2500 galaxies using the multi-spectrograph “Flair” on the UKST (AAO), Wakamatsu et al. (2000) confirmed that the Ophiuchus cluster behind the Galactic bulge ($`0.^{}5,9.^{}5,8500`$ km s<sup>-1</sup>) is the central component of a supercluster including two more clusters and four groups of galaxies. There seems to be a wall-like structure connecting the Ophiuchus cluster to the Hercules supercluster at 11000 km s<sup>-1</sup>. This wall runs orthogonal to the Great Wall.
At ($`\mathrm{},b`$)$``$ ($`33{}_{}{}^{},5{}_{}{}^{}15^{}`$), Marzke et al. (1996) and Roman et al. (1998) found evidence for a nearby cluster close to the Local Void at 1500 km s<sup>-1</sup>, as well as a prominent cluster behind the Local Void at 7500 km s<sup>-1</sup>. The nearby cluster is independently supported by data from blind H I-surveys (see Sect. 8.).
The connection of the Perseus-Pisces supercluster across the ZOA to the cluster A569, suspected by Focardi et al. (1984), was confirmed by Chamaraux et al. (1990) and Pantoja et al. (1997). The Perseus-Pisces chain seems to fold back into the ZOA at higher redshifts at ($`95{}_{}{}^{},10{}_{}{}^{},7500`$ km s<sup>-1</sup>), Marzke et al. (1996), Pantoja et al. (1997).
In 1992, Kraan-Korteweg & Huchtmeier uncovered a nearby cluster in Puppis ($`245{}_{}{}^{},0{}_{}{}^{},1500`$ km s<sup>-1</sup>) which was later shown by Lahav et al. (1993) to contribute a non-negligible component to the peculiar z-motion of the Local Group. Since then, this region has been investigated in further detail. Chamaraux et al. (1999) found further evidence for this cluster. It lies within a long narrow filament (Masnou & Chamaraux, in prep.) which extends from the Antlia to the Fornax cluster (see also Fig. 12).
Kraan-Korteweg et al. (1994a) presented evidence for a continous filamentary structure extending over $`30^{}`$ on the sky from the Hydra and Antlia clusters across the ZOA, intersecting the Galactic Plane at ($`280{}_{}{}^{},0{}_{}{}^{},3000`$ km s<sup>-1</sup>). At the same longitudes, they noted significant clustering at $``$ 15000 km s<sup>-1</sup>, indicative of a connection between the Horologium and Shapley clusters a hundred degrees apart in the sky.
Kraan-Korteweg & Woudt (1993) uncovered a shallow but extended supercluster in Vela at ($`285{}_{}{}^{},6{}_{}{}^{},6000`$ km s<sup>-1</sup>).
Next to the massive cluster A3627 at the core of the Great Attractor (clustering in the Great Attractor region is discussed in the next section), Woudt (1998) discovered a cluster at ($`306{}_{}{}^{},6{}_{}{}^{},6200`$) called the Cen-Crux cluster, and a more distant cluster, the Ara cluster at ($`329{}_{}{}^{},9{}_{}{}^{},15000`$ km s<sup>-1</sup>). The latter might be connected to the Triangulum-Australis cluster.
#### Clustering within the Great Attractor region
Based on a deep optical galaxy search and subsequent redshift follow-ups, Kraan-Korteweg et al. (1996) and Woudt (1998) have clearly shown that the Norma cluster, A3627, at ($`325{}_{}{}^{},7^{}`$, 4848 km s<sup>-1</sup>) is the most massive galaxy cluster in the GA region known to date and probably marks the previously unidentified but predicted density-peak at the bottom of the potential well of the GA overdensity. The prominence of this cluster has independently been confirmed by ROSAT observations: the Norma cluster ranks as the 6<sup>th</sup> brightest X-ray cluster in the sky (Böhringer et al. 1996). It is comparable in size, richness and mass to the well-known Coma cluster. Redshift-independent distance determinations (R<sub>C</sub> and I<sub>C</sub> band Tully – Fisher relation analysis) of the Norma cluster have shown it to be at rest with respect to the rest frame of the Cosmic Microwave Background (Woudt 1998).
One cannot, however, exclude the possibility that other unknown rich clusters reside in the GA region as the ZOA has not been fully unveiled with optical searches. Finding a hitherto uncharted, rich cluster of galaxies at the heart of the GA would have serious implications for our current understanding of this massive overdensity in the local Universe. Kraan-Korteweg & Woudt (1999) found various indications that PKS1343$``$601, the second brightest extragalactic radio source in the southern sky ($`f_{20cm}=79`$ Jy, McAdam 1991, and references therein) might form the center of yet another highly obscured rich cluster, particularly as it also shows significant X-ray emission (Tashiro et al. 1998): extended diffuse hard X-ray emission at the position of PKS1343$``$601 has been detected with ASCA. The radiation, kT = 3.9 keV, is far too large for it being associated with a galactic halo surrounding the host galaxy, hence it might be indicative of emission from a cluster – if it is not due to the Inverse Compton process.
At ($`\mathrm{},b)(310{}_{}{}^{},2^{}`$), this radio galaxy lies behind an obscuration layer of about 12 magnitudes of extinction in the B-band, as estimated from the DIRBE extinction maps (Schlegel et al. 1998). Its observed diameter of 28 arcsec in the Gunn-z filter (West & Tarenghi 1989) translates into an extinction-corrected diameter of 232 arcsec (following Cameron 1990). With a recession velocity of $`v=3872`$ km s<sup>-1</sup> (West & Tarenghi 1989), this galaxy can be identified with a giant elliptical galaxy.
Since PKS1343$``$601 is so heavily obscured, little data are available to substantiate the existence of this prospective cluster. In Fig. 5, the A3627 cluster at a mean extinction $`A_B=1.^\mathrm{m}5`$ as seen in deep optical searches is compared to the prospective PKS1343 cluster at ($`309.^{}7,+1.^{}7,3872`$ km s<sup>-1</sup>) with an extinction of 12<sup>m</sup>. One can clearly see that at the low Galactic latitude of the suspected cluster PKS1343, the optical galaxy survey could not retrieve the underlying galaxy distribution, especially not within the Abell radius of the suspected cluster (the inner circle in Fig. 5). To verify this cluster, other observational approaches are necessary. We have imaged the prospective cluster within its Abell radius in the near infrared (Woudt et al. in progress). These observations will allow us to determine whether or not PKS1343$``$601 is embedded in a centrally condensed overdensity of galaxies. Interestingly enough, deep H I observations did uncover a significant excess of galaxies at this position in velocity space (see Sect. 8.2.) although a “finger of God”, the characteristic signature of a cluster in redshift space, is not seen. Hence, the Norma cluster A3627 remains the best candidate for the center of the extended GA overdensity.
### 5.5. The Sagittarius dwarf
A remarkable discovery was made by Ibata et al. (1994). They found a nearby ‘dwarf’ galaxy with diameter of about 3 kpc. This new galaxy, named the Sagittarius dwarf, is on the far side of the Galactic center, about 25 kpc away from us, but well inside the Milky Way. This galaxy is most probably undergoing some tidal disruption, before being absorbed by the Milky Way. An interesting feature of the discovery of this galaxy is that it was based on velocities, not direct detection on plates. In fact, the Sagittarius dwarf galaxy is some twenty degrees from end to end, making it the largest structure in the sky after the Milky Way itself. Nonetheless, since the new galaxy lies directly behind the central bulge of the Milky Way, it cannot be seen in direct images - even with hindsight - through the dust and against the very much larger number of Galactic stars.
The Sagittarius dwarf galaxy provides an important clue to the formation process of the Milky Way. Many popular models of galaxy formation suggest that large galaxies are formed by a long process of aggregation of many smaller galaxies, possibly with some merger events being disruptive of the normal disk structure of spiral galaxies. Such a process should still be common today, yet it had been observed previously only in extremely rare cases. The Sagittarius dwarf merger with the Milky Way provides the ‘smoking gun’, showing that such mergers do happen, they happen today, and they are not destructive.
### 5.6. Conclusions
Deep optical galaxy searches have succesfully reduced the solid angle of the ZOA by a factor of about $`22.5`$ down to extinction levels of $`A_B=3.^\mathrm{m}0`$ and have identified a number of important unknown structures. However, they fail in the most opaque part of the Milky Way, the region encompassed by the $`A_B=3.^\mathrm{m}0`$ contour in Fig. 4 – a sufficiently large region to hide further dynamically important galaxy densities. Here, systematic surveys in other wavebands can be applied to reduce the current ZOA even further. The success and status of these approaches are discussed in the following sections.
## 6. Near infrared surveys
Observations in the near infrared (NIR) can provide important complementary data to other surveys. With extinction decreasing as a function of wavelength, NIR photons are up to 10 times less affected by absorption compared to optical surveys – an important aspect in the search and study of galaxies behind the obscuration layer of the Milky Way. The NIR is sensitive to early-type galaxies – tracers of massive groups and clusters – which are missed in IRAS and H I surveys (Sect. 7. and 8.). In addition, confusion with Galactic objects is considerably lower compared to the FIR surveys. Furthermore, because recent star formation contributes only little to the NIR flux of galaxies (in contrast to optical and FIR emission), NIR data give a better estimation of the stellar mass content of galaxies. It is therefore well suited for the application of the Tully – Fisher relation either through pointed H I observations of galaxies detected in the NIR or a merging of detections from systematic blind H I surveys with NIR observations (Sect. 6.3.).
### 6.1. The NIR surveys DENIS and 2MASS
Two systematic near infrared surveys are currently being performed: DENIS, the DEep Near Infrared Southern Sky Survey, is imaging the southern sky from $`88{}_{}{}^{}<\delta <+2^{}`$ in the $`I_c`$ (0.8$`\mu `$m), $`J`$ ($`1.25\mu `$m) and $`K_s`$ ($`2.15\mu `$m) bands. 2MASS, the 2 Micron All Sky Survey, is covering the whole sky in the $`J`$ ($`1.25\mu `$m), $`H`$ ($`1.65\mu `$m) and $`K_s`$ ($`2.15\mu `$m) bands. The mapping of the sky is performed in declination strips, which are $`30^{}`$ in length and 12 arcmin wide for DENIS, and $`6{}_{}{}^{}\times 8.^{}5`$ for 2MASS. Both the DENIS and 2MASS surveys are expected to complete their observations by the end of 2000. The main characteristics of the 2 surveys and their respective completeness limits for extended sources are given in Table 1 (Epchtein 1997, 1998, Skrutskie et al. 1997, Skrutskie 1998).
Details and updates on completeness, data releases and data access for DENIS and 2MASS can be found on the websites http://www-denis.iap.fr, and http://www.ipac.caltech.edu/2mass, respectively.
The DENIS completeness limits (total magnitudes) for highly reliable automated galaxy extraction (determined away from the ZOA, i.e. $`|b|>10^{}`$) are $`I=16.^\mathrm{m}5`$, $`J=14.^\mathrm{m}8`$, $`K_s=12.^\mathrm{m}0`$ (Mamon 1998). The number counts per square degrees for these completeness limits are 50, 28 and 3 respectively. For the 2MASS, the completeness limits are $`J=15.^\mathrm{m}0`$, $`H=14.^\mathrm{m}2`$, $`K_s=13.^\mathrm{m}5`$ (isophotal magnitudes), with number counts of 48, $``$40 and 24. Because no $`H`$ counts have been published we estimated their counts from low, mid and high-density fields and the $`JH`$ and $`HK`$ color distribution (cf. http://spider.ipac.caltech.edu/staff/jarrett/2mass). The numbers agree well with the 2MASS predictions of the detection of 1 million galaxies on the whole sky for their $`K_s`$ band completeness limit. In all wavebands, except $`I_c`$, the number counts are still imprecise due to the low number statistics and the strong dependence on the star crowding in the analyzed fields. Still, these numbers suffice to reveal the promise of NIR surveys for probing the galaxy distribution at very low Galactic latitudes.
The decrease in number counts is considerably slower in the $`I_c`$, $`J`$, $`H`$ and $`K_s`$ bands compared to the optical because extinction is only 45%, 21%, 14% and 9% compared to the $`B`$ band. This dependence makes NIR surveys very powerful at low Galactic latitudes even though they are not as deep as the POSS and ESO/SERC sky surveys. As illustrated in Fig. 6, the galaxy density in the $`B`$ band in unobscured regions is 110 galaxies per square degree for the completeness limit of $`B_J19.^\mathrm{m}0`$ (Gardner et al. 1996). But the counts in the blue decrease rapidly with increasing obscuration: $`N(A_B)110\times \mathrm{dex}(0.6[A_B])`$deg<sup>-2</sup>. The decrease in detectable galaxies due to extinction is slower in the NIR, and the counts of the shallower NIR surveys overtake the optical counts at extinction levels of $`A_B>13^m`$. The location of the reversal in efficiency is particularly opportune because the NIR surveys become more efficient where deep optical galaxy searches become incomplete, i.e. at $`A_B>3.^\mathrm{m}0`$ (see Sect. 5.3.).
The above predictions do not take into account any dependence on morphological type, surface brightness, intrinsic color, orientation and crowding, which may lower the counts of actually detectable galaxies counts. In practice, $`B`$ was found to be superior for identification of galaxies on DENIS images to extinction levels of at least $`A_B=2.^\mathrm{m}0`$. And even though 2MASS appears more powerful in Fig. 6 for ZOA research compared to DENIS, the higher sensitivity of 2MASS also results in higher star densities at low latitudes, making galaxy identifications more difficult. This problem becomes apparent from, for instance, Fig. 28 and Fig. 29a in the ’List of Figures’ accessible from
http://spider.ipac.caltech.edu/staff/jarrett/2mass/3chan/basic/paper\_ I.html
which show galaxy images in $`J`$, $`H`$ and $`K_s`$ in the ZOA at extinction levels which are not yet very severe ($`0.^\mathrm{m}8`$ and $`2.^\mathrm{m}6`$ in the optical).
### 6.2. Pilot studies with DENIS data in the ZOA
To compare the above predictions with real data, Schröder et al. (1997, 1999) and Kraan-Korteweg et al. (1998b) examined the efficiency of uncovering galaxies at high extinctions using DENIS images. The analyzed regions include the rich cluster A3627 ($`\mathrm{},b)=(325.^{}3,7.^{}2)`$ at the heart of the GA (Norma) supercluster, as well as its suspected extension across the Galactic Plane.
Three high-quality DENIS strips cross the cluster A3627. 66 images on these strips lie within the Abell-radius ($`R_A=1.^{}75`$) and were inspected by eye (Kraan-Korteweg et al. 1998a). This covers about one-eighth of the cluster area. The extinction over the regarded cluster area varies as $`1.^\mathrm{m}2A_B2.^\mathrm{m}0`$.
On these 66 images, 151 galaxies had previously been identified in the deep optical ZOA galaxy search (Woudt & Kraan-Korteweg 2000b). Of these, 122 were recovered in the $`I_c`$, 100 in the $`J`$, and 74 in the $`K_s`$ band. Most of the galaxies not re-discovered in $`K_s`$ are low surface brightness spiral galaxies.
Surprisingly, the $`J`$ band provides better galaxy detection than the $`I_c`$ band. In the latter, the severe star crowding makes identification of faint galaxies very difficult. At these extinction levels, the optical survey does remain the most efficient in identifying obscured galaxies.
The search for more obscured galaxies was made in the region $`320{}_{}{}^{}\mathrm{}325^{}`$ and $`|b|5^{}`$, i.e. the suspected crossing of the GA. Of the 1800 images in that area, 385 of the then available DENIS images were inspected by eye (308 in $`K_s`$). 37 galaxies at higher latitudes were known from the optical survey. 28 of these could be re-identified in $`I_c`$, 26 in $`J`$, and 14 in the $`K_s`$ band. In addition, 15 new galaxies were found in $`I_c`$ and $`J`$, 11 of which also appear in the $`K_s`$ band. The ratios of galaxies found in $`I_c`$ compared to $`B`$, and of $`K_s`$ compared to $`I_c`$ are higher than in the A3627 cluster. This is due to the higher obscuration level (starting with $`A_B2.^\mathrm{m}33.^\mathrm{m}1`$ at the high-latitude border).
On average, about 3.5 galaxies per square degree were found in the $`I_c`$ band. This roughly agrees with the predictions of Fig. 6. Because of star crowding, we do not expect to find galaxies below latitudes of $`b1{}_{}{}^{}2^{}`$ in this longitude range (Mamon 1994). Low-latitude images substantiate this – the images are nearly fully covered with stars. Indeed, the lowest Galactic latitude galaxies were found at $`b1.^{}2`$ and $`A_B11^\mathrm{m}`$ (in $`J`$ and $`K_s`$ only).
Figure 7 shows a few characteristic examples of highly obscured galaxies found in the DENIS blind search. $`I_c`$ band images are at the top, $`J`$ in the middle and $`K_s`$ at the bottom. The first galaxy located at $`(l,b)=(324.^{}6,4.^{}5`$) is viewed through an extinction layer of $`A_B=2.^\mathrm{m}0`$ according to the DIRBE extinction maps (Schlegel et al. 1998). It is barely visible in the $`J`$ band. The next galaxy at $`(l,b)=(324.^{}7,3.^{}5`$) is subject to heavier extinction ($`A_B=2.^\mathrm{m}7`$), and indeed easier to recognize in the NIR. It is most distinct in the $`J`$ band. The third galaxy at even higher extinction $`(l,b,A_B)=(320.^{}1,+2.^{}5,5.^\mathrm{m}7`$) is – in agreement with the prediction of Fig. 6 – not visible in the $`B`$ band. Neither is the fourth galaxy at $`b=+1.^{}9`$ and $`A_B=9.^\mathrm{m}6`$: this galaxy can not be seen in $`I_c`$ band either and is very faint only in $`J`$ and $`K_s`$.
The conclusions from this pilot study are that at intermediate latitudes and extinction ($`|b|>5^{}`$, $`A_B<45^\mathrm{m}`$) optical surveys are superior for identifying galaxies. But despite the extinction and the star crowding at these latitudes, $`I_c`$, $`J`$ and $`K_s`$ photometry from the survey data could be performed successfully at these low latitudes. The NIR data (magnitudes, colors) of these galaxies can therefore add important data in the analysis of these obscured galaxies. They led, for instance, to the preliminary $`I_c^o`$, $`J^o`$ and $`K_s^o`$ galaxy luminosity functions in A3627 (Fig. 2 in Kraan-Korteweg et al. 1998a).
At the lowest latitudes and at high extinction ($`|b|<5^{}`$ and $`A_B>45^\mathrm{m}`$), the search for ‘invisible’ obscured galaxies on existing DENIS-images imply that NIR-surveys can trace galaxies down to about $`|b|>1{}_{}{}^{}1.^{}5`$. The $`J`$ band was found to be optimal for identifying galaxies up to $`A_B7^\mathrm{m}`$. NIR surveys can hence further reduce the width of the ZOA.
### 6.3. Systematic exploitation of NIR surveys and redshift follow-ups
At the Observatoire de Lyon (France), DENIS images are being processed routinely as they come out of the pipeline (cf. Vauglin et al. 1999 for a catalog of over 20000 galaxies found on the $`I_c`$ band images in 25% of the southern sky). Concerning the ZOA, $`J`$ and $`K_s`$ band DENIS images are systematically being inspected by eye for galaxies below $`|b|15^{}`$ since March 1997. This so far has led to the detection of 1500 extended objects most of which were unknown previously. For these galaxies, magnitudes, isophotal diameters, axis ratios, position angles and a rough estimate of the morphological type are recorded. This catalog will build the basis for systematic H I line redshift follow-ups with the goal of applying the Tully – Fisher relation to determine the peculiar velocity field. Using a morphological classification based upon a concentration parameter (Theureau et al. 1997), they want to obtain a complete sample of late-type inclined spirals within a volume of $`v<10000`$ km s<sup>-1</sup>.
A similar project, but starting out from an H I-selected sample, is being pursued by cross-identifying all galaxies detected in the the systematic deep blind H I survey (Sect. 8.2.) in the ZOA ($`|b|5.^{}5`$) on NIR images (Schröder et al. in prep.) – also with the aim of determining the density field from the peculiar velocity field through the Tully – Fisher relation.
A collaboration involving UMASS, IPAC and the CfA is working on a redshift survey (the 2MASS Redshift Survey, 2MRS) of the whole sky that will contain $`250000`$ galaxies to a limiting magnitude of $`K_s`$ $`=13.^\mathrm{m}5`$. The first phase will be directed at galaxies with $`K_s`$ $`12.^\mathrm{m}2`$. This survey will include the ZOA and although the shallow redshift survey so far is only 25% complete, it already shows considerable structure through the Galactic Plane out to high velocities (http://cfa-www.harvard.edu/~huchra/2mass/).
Two spectroscopic follow-ups of DENIS and 2MASS galaxies are planned on the 6dF robotic multi-object spectroscopic unit, currently under construction at the AAO: a redshift survey of roughly 120000 NIR-selected galaxies for which a total of 300 nights are guaranteed between 2001 and 2003 and a peculiar velocity survey of roughly 12000 early-type galaxies. The latter will be complemented by an H I peculiar velocity survey of over 5000 inclined spirals (Mamon 1998, 1999). In this project, however, only selected areas in the ZOA will be probed.
### 6.4. Conclusions
First results from NIR data are very promising for ZOA research – and complementary to other approaches. NIR surveys become more efficient in revealing galaxies at extinction levels where deep optical searches become increasingly incomplete, i.e. at $`A_B3^\mathrm{m}`$. It was found that galaxies can be traced to Galactic latitudes of $`|b|>11.^{}5`$.
A independent advantage of the NIR surveys is the fact that the NIR colors, in particular of early-type galaxies, might help in the calibration of the DIRBE extinction maps at low Galactic latitudes (see e.g. Fig. 5 in Schröder et al. 1997).
The NIR surveys are particularly useful for the mapping of massive early-type galaxies – tracers of density peaks in the mass distribution – as these can not be detected with any of the techniques that are efficient in tracing the spiral population in more opaque regions (Sect. 7. and 8.).
Nevertheless, NIR surveys are also important with regard to the blue and low surface-brightness spiral galaxies because a significant fraction of them are detectable in the near infrared. This is confirmed, for instance, with the serendipitous discovery in the ZOA of a large, nearby ($`v_{LSR}=750`$ km s<sup>-1</sup>) edge-on spiral galaxy by 2MASS (Hurt et al. 1999): with an extension in the $`K_s`$ band of 5 arcmin, this large galaxy is – not unexpectedly for its extinction of $`A_B=6.^\mathrm{m}6`$ at the position of $`(\mathrm{},b)=(236.^{}8,1.^{}8)`$ – not seen in the optical (Saito et al. 1991). A more systematic analysis of low-latitude spiral galaxies in the NIR was undertaken by Schröder et al. (priv. comm): about 80% of the spiral galaxies detected in the shallow systematic H I survey performed with the Parkes Multibeam Receiver (Sect. 8.2.) could be reidentified on $`I_c`$ band DENIS images, a few more in $`J`$ and $`K_s`$ only. The above nearby NIR galaxy with a flux of 33.7 Jy km s<sup>-1</sup> in the 21 cm line was discovered independently in this survey (HIZSS012, Henning et al. 2000).
The overlap of galaxies found in NIR and H I surveys is important. With the combination of H I data and NIR data one can study the peculiar velocity field via the NIR Tully – Fisher relation “in the ZOA” compared to earlier interpolations of data “adjacent to the ZOA”. This will provide important new input for density field reconstructions in the ZOA (Sect. 10.).
## 7. Far infrared surveys
In 1983, the Infrared Astronomical Satellite IRAS surveyed 96% of the whole sky in the far infrared bands at 12, 25, 60 and 100 $`\mu `$m, resulting in a catalog of 250 000 point sources, i.e. the IRAS Point Source Catalogue (IRAS PSC; Joint IRAS Science Working Group 1988). The latter has been used extensively to quantify extragalactic large-scale structures. The identification of the galaxies from the IRAS data base is quite different compared to the optical: only the fluxes at the 4 far infrared (FIR) IRAS passbands are available but no images, and the identification of galaxies is strictly based on the flux ratios. For instance, Yamada et al. (1993) used the criteria: 1. $`f_{60}>0.6`$Jy, 2. $`f_{60}^2>f_{12}f_{25}`$, 3. $`0.8<f_{100}/f_{60}<5.0`$, to select galaxy candidates from the IRAS PSC.
With these flux and color criteria, mainly normal spiral galaxies and starburst galaxies are identified. Hardly any dwarf galaxies enter the IRAS galaxy sample, nor the dustless elliptical galaxies, as they do not radiate in the far infrared. The upper cut-off in the third criterion is imposed to minimize the contamination with cool cirrus sources and young stellar objects within our Galaxy. This, however, also makes the IRAS surveys less complete for nearby galaxies (e.g. Woudt 1998, Kraan-Korteweg 2000).
The advantage of using IRAS data for large-scale structure studies is its homogeneous sky coverage (all data from one instrument) and the negligible effect of the extinction on the flux at these long wavelengths. Even so, it remains difficult to probe the inner part of the ZOA with IRAS data because of cirrus, high source counts of Galactic objects in the Galaxy, and confusion with these objects – most of them have the same IRAS characteristics as external galaxies. The difficulty in obtaining unambiguous galaxy identifications at these latitudes was demonstrated by Lu et al. (1990), who found that the detection rate of IRAS galaxy candidates decreases strongly as a function of Galactic latitude (from $`|b|=16^{}`$ to $`|b|=2^{}`$). This can only be explained by the increase in faulty IRAS galaxy identifications. Yamada et al. (1993) also found a dramatic and unrealistic increase in possible galaxies close to the Galactic Plane in their systematic IRAS galaxy survey of the southern Milky Way ($`|b|15^{}`$).
So, despite the various advantages given with IRAS data, the sky coverage in which reliable IRAS galaxy identifications can be made (84%) provides only a slight improvement over optical galaxy catalogs (compare e.g. the light-grey mask in Fig. 8 with the optical ZOA-contour as displayed in Fig. 1). In addition to that, the density enhancements are very weak in IRAS galaxy samples because (a) the IRAS luminosity function is very broad, which results in a more diluted distribution since a larger fraction of distant galaxies will enter a flux-limited sample compared to an optical galaxy sample, and (b) IRAS is insensitive to elliptical galaxies, which reside mainly in galaxy clusters, and mark the peaks in the mass density distribution of the Universe. This is quite apparent when comparing the IRAS galaxy distribution (Fig. 8) with the optical galaxy distribution (Fig. 1 and Fig. 4).
Nevertheless, dedicated searches for large-scale clustering within the whole ZOA ($`|b|15^{}`$) have been made by various Japanese collaborations (see Takata et al. 1996 for a summary). They used IRAS color criteria to select galaxy candidates which were subsequently verified through visual examination on sky surveys, such as the POSS for the northern hemisphere and the ESO/SRC for the southern sky. Because of their verification procedure, this data-set suffers, however, from the same limitations in highly obscured regions as optical surveys.
Based on redshift follow-ups of these ZOA IRAS galaxy samples, they established various filamentary features and connections across the ZOA. Most coincide with the structures uncovered in optical work. In the northern Milky Way, both crossings of the Perseus-Pisces arms into the ZOA are very prominent – considerably stronger in the FIR than at optical wavelengths – and they furthermore identified a new structure: the Cygnus-Lyra filament at ($`60{}_{}{}^{}90{}_{}{}^{},0{}_{}{}^{},4000`$ km s<sup>-1</sup>). Across the southern Milky Way they confirmed the three general concentrations of galaxies around Puppis ($`\mathrm{}=245^{}`$), the Hydra–Antlia extension ($`\mathrm{}=280^{}`$; Kraan-Korteweg et al. 1995) and the Centauraus Wall ($`\mathrm{}=315^{}`$). However, the cluster A3627 is not seen, nor is the Great Attractor very prominent compared to the optical or to the POTENT reconstructions described in Sect. 10..
Besides the search for the continuity of structures across the Galactic Plane, the IRAS galaxy samples have been widely used for the determination of the peculiar motion of the Local Group, as well as the reconstructions of large-scale structure across the Galactic Plane (see Sect. 10.). These analyses have been performed on the two-dimensional IRAS galaxy distribution and, in recent years, as well as on their distribution in redshift space through the availability of redshift surveys for progressively deeper IRAS galaxy samples, i.e. 2658 galaxies to f$`{}_{60\mu m}{}^{}=1.9`$ Jy (Strauss et al. 1992), 5321 galaxies to f$`{}_{60\mu m}{}^{}=1.2`$ Jy (Fisher et al. 1995), and lately the PSCz catalog of 15411 galaxies complete to f$`{}_{60\mu m}{}^{}=0.6`$ Jy with 84% sky coverage and a depth of 20000 km s<sup>-1</sup> (Saunders et al. 2000b).
The PSCz is deep enough to test the convergence of the dipole. The most recent analysis of the IRAS PSCz dipole (Schmoldt et al. 1999; see also Rowan-Robinson et al. 2000) finds that the acceleration vector points about $`15^{}`$ away from the CMB dipole. Assuming full convergence at the sample boundary, about 2/3 of the measured acceleration is generated within 4000 km s<sup>-1</sup>. There is a non-negligible contribution out to 14000 km s<sup>-1</sup>, after which the acceleration amplitude seems to have converged.
Saunders and collaborators realized, however, that the 16% of the sky missing from the survey causes significant uncertainty, particularly because of the location behind the Milky Way of many of the prominent large-scale structures (superclusters as well as voids). In 1994, they therefore started a longterm program to increase the sky coverage of the PSCz. Optimizing their color criteria to minimize contamination by Galactic sources ($`f_{60}/f_{25}>2`$, $`f_{60}/f_{12}>4`$, and $`1.0<f_{100}/f_{60}<5.0`$), they extracted a further 3500 IRAS galaxy candidates at lower Galactic latitudes (light-grey area of Fig. 8). Taking $`K^{}`$ band snapshots of all the galaxy candidates of their ‘Behind The Plane’ \[BTP\] survey, they could add a thousand galaxies to the PSCz sample and reduce the coverage gap to a mere 7% (dark-grey area).
The resulting sky map of 16,400 galaxies (PSCz plus BTP) is shown in Fig. 8 (from Saunders et al. 2000a). The BTP survey has reduced the “IRAS ZOA” dramatically. Some incompleteness remains towards the Galactic Center, but large-scale structures can easily be identified across most of the Galactic Plane. In the Great Attractor region, the galaxies can be traced (for the first time with IRAS data) to the rich cluster A3627 – the suspected core of the GA (Kraan-Korteweg et al. 1996). The BTP collaboration is currently working hard on obtaining redshifts for these new and heavily obscured galaxies and exciting new results on large-scale structure across the Milky Way and dipole determinations can be expected in the near future.
## 8. HI surveys
In the regions of the highest obscuration and infrared confusion, the Galaxy is fully transparent to the 21cm line radiation of neutral hydrogen. H I-rich galaxies can readily be found at lowest latitudes through the detection of their redshifted 21cm emission, though early-type galaxies – tracers of massive groups and clusters – are gas-poor and will not be identified in these surveys. Furthermore, low-velocity extragalactic sources (blue- and red-shifted) within the strong Galactic H I emission will be missed, and galaxies close to radio continuum sources may also be missed because of baseline ripples.
The advantage of blind H I surveys in the Milky Way is not only the transparency of the 21cm radiation to the thickest dust layers: with the detection of an H I signal, the redshift and rotational properties of an external galaxy are immediately known, providing insight not only on its location in redshift space but also on the intrinsic properties of such obscured galaxies. The rotational velocity can furthermore be used in combination with e.g. NIR photometry (Sect. 6.3.) to determine the distribution in real space from Tully – Fisher relation distances and the density field independent of interpolations across the Milky Way from the peculiar velocity field.
Until recently, radio receivers were not sensitive and efficient enough to attempt large systematic surveys of the ZOA. But in a pilot survey, Kerr & Henning (1987) pointed the late 300-ft telescope of Green Bank to 1900 locations in the ZOA (1.5% coverage) and detected 19 previously unknown spiral galaxies, proving therewith the effectiveness of this approach.
Since then two systematic blind H I searches for galaxies behind the Milky Way were initiated. The first – the Dwingeloo Obscured Galaxies Survey (DOGS) – used the 25 m Dwingeloo radio to survey the whole northern Galactic Plane for galaxies out to 4000 km s<sup>-1</sup> (cf. Kraan-Korteweg et al. 1994b, Henning et al. 1998, Rivers et al. 1999). Although the sensitivity was fairly low (40 mJy for a 1 hr integration), the advantage of the small telescope aperture is rapid areal coverage.
A more sensitive survey, probing a considerably larger volume (out to 12700 km s<sup>-1</sup>), is being performed for the southern Milky Way with the 64 m radio telescope of Parkes (Kraan-Korteweg et al. 1998a, Staveley-Smith et al. 1998, Henning et al. 1999, 2000). A Multibeam (MB) receiver with 13 beams in the focal plane array (Staveley-Smith 1996) was specifically constructed to efficiently search for galaxies not identified in optical surveys because of low optical surface brightness or high optical extinction.
In the following, the observing techniques of these two surveys as well as the first results will be discussed.
### 8.1. The Dwingeloo Obscured Galaxies Survey (DOGS)
Since 1994, the Dwingeloo 25 m radio telescope has been dedicated to a systematic search for galaxies in the northern Zone of Avoidance ($`30{}_{}{}^{}\mathrm{}220^{}`$, $`|b|5.^{}25`$). The last few patches of the survey were completed early 1999 using the Westerbork array in total power mode. The 20 MHz bandwidth was tuned to cover the velocity range $`0v_{\mathrm{LSR}}4000`$ km s<sup>-1</sup>. Negative velocities were excluded because the Leiden/Dwingeloo Galactic H I survey (Hartmann 1994, Hartmann & Burton 1997) had already covered the velocity range $`450v_{\mathrm{LSR}}400`$ km s<sup>-1</sup>, albeit with higher rms.
The 25 m Dwingeloo telescope has a half-power-beamwidth (HPBW) of 36 arcmin. With a DAS-1000 channel autocorrelator spectrometer at the telescope backend, the coverage over the 20 MHz bandwidth resulted in a velocity resolution of 4 km s<sup>-1</sup>. With this resolution even the galaxies with the narrowest linewidth are covered by several channels. The 15000 survey points are ordered in a honeycomb pattern with a grid spacing of $`0.^{}4`$. Galaxies are generally detected in various adjacent pointings, facilitating a more accurate determination of their positions through interpolations. Each DOGS observation consisted of a sequence of 5 contiguous pointings at constant Galactic latitude. From this, 5 On-Off pairs were created in such a way that a real galaxy will appear once as a positive and once as a negative signal in two independent scans. The rms noise per channel typically was $`\sigma _{ch}=40`$ mJy for a 1 hr integration (12 x 5min).
Because of the duration of the project (15000 hours not including overhead and downtime) the strategy was to first conduct a fast search of 5min integrations ($`\mathrm{rms}=175`$ mJy) to uncover possible massive nearby galaxies whose effect might yield important clues to the dynamics of the Local Group. In the following, the results from the Dwingeloo shallow and deep surveys are discussed.
#### The shallow Dwingeloo survey and the discovery of Dwingeloo 1
The shallow Dwingeloo search ($`\mathrm{rms}=175`$ mJy) has been completed in 1996 yielding five objects (Henning et al. 1998), three of which were known previously. The most exciting discovery was the barred spiral galaxy Dwingeloo 1 (Kraan-Korteweg et al. 1994b).
This galaxy candidate was detected early on in the survey through a strong signal (peak intensity of 1.4 Jy) at the very low redshift of $`v_{\mathrm{LSR}}=110`$ km s<sup>-1</sup> in the spectra of four neighbouring pointings, suggestive of a galaxy of large angular extent. The optimized position of $`(\mathrm{},b)=(138.^{}5,0.^{}1)`$ coincided with a very low surface brightness feature on the Palomar Sky Survey plate of $`2.^{}2`$, detected earlier by Hau et al. (1995) in his optical galaxy search of the northern Galactic/Supergalactic Plane crossing (cf. Sect. 5.2.). Despite foreground obscuration of about 6<sup>m</sup> in the optical, follow-up observations in the $`V`$, $`R`$ and $`I`$ band at the INT (La Palma) confirmed this galaxy candidate as a barred, possibly grand-design spiral galaxy of type SBb of 4.2 x 4.2 arcmin (cf. Fig. 9).
Dwingeloo 1 has been the subject of many follow-up observations (optical: Loan et al. 1996; Buta & McCall 1999; HI-synthesis: Burton et al. 1996; CO observations: Kuno et al. 1996; Li et al. 1996; Tilanus & Burton 1997; X-ray: Reynolds et al. 1997). To summarize, Dwingeloo 1 is a barred spiral, with a rotation velocity of 130 km s<sup>-1</sup>, implying a dynamical mass of roughly one-third the mass of the Milky Way (within the same fixed radius). Its approximate distance of $``$ 3 Mpc and angular location place it within the IC342/Maffei group of galaxies. The follow-up HI synthesis observations (Burton et al. 1996) furthermore revealed a counterrotating dwarf companion, Dwingeloo 2. Since then various further dwarf galaxies have been discovered in this nearby galaxy group (next section).
#### The deep Dwingeloo survey
Currently, 60% of the deeper Dwingeloo survey ($`\mathrm{rms}=40`$ mJy) have been analyzed resulting in 36 detected galaxies, 23 of which were previously unknown (Rivers et al. 1999). Five of the 36 sources were originally identified in the shallow survey. Based on the survey sensitivity, the registered number of galaxies is in agreement with the Zwaan et al. (1997) HI mass function which predicts 50 to 100 detections for the full survey.
The distribution of the Dwingeloo galaxies is shown in Fig. 10 together with other known galaxies out to 4000 km s<sup>-1</sup> for visualization of connectivity of structures across the Galactic Plane. Indeed, various known structures appear continuous across the GP, although two galaxies were also found in the Local Void ($`\mathrm{}30^{}`$). The latter two were independently detected in the Parkes MB ZOA survey (Sect. 8.2., and Henning et al. 1999, 2000) and seem to be part of the by Roman et al. (1998) recently discovered clustering at 1500 km s<sup>-1</sup>(Sect. 5.4.).
11 galaxies were discovered in the Supergalactic Plane crossing region ($`\mathrm{}137.^{}4`$), five of which in the nearby IC342/Maffei group: Dwingeloo 1, Maffei 2 from the shallow survey, plus three dwarf members. Further group members have been detected through infrared photometry (McCall & Buta 1995), and pointed observations with the 100 m radiotelescope at Effelsberg of optically identified dwarf candidates (Huchtmeier et al. 2000). The total number of obscured IC342/Maffei group members has meanwhile grown to 19 members (for recent updates on this group see Huchtmeier et al. 2000; Buta & McCall 1999). Photometric distances have been derived for 10 of these galaxies (Karachentsev & Tikhonov 1993, 1994; Karachentsev et al. 1997), putting the group at a mere distance of $`2.2\pm 0.5`$ Mpc. At such a close distance, this group might have played a significant role in the dynamical history of the Local Group (McCall 1986, 1989; Zheng et al. 1991; Valtonen et al. 1993; Peebles 1994; see Sect. 3.3.).
Suprisingly three dwarf galaxies were detected close to the nearby isolated galaxy NGC 6946 at ($`\mathrm{},b,v_{\mathrm{LSR}})=(95.^{}7,11.^{}7,46`$ km s<sup>-1</sup>). One of these had earlier been cataloged as a compact high-velocity cloud (Wakker 1990). Burton et al. (1999), in their search for compact isolated high-velocity clouds in the Dwingeloo/Leiden Galactic H I survey, discovered a further member of this galaxy concentration. Now, seven galaxies with recessional velocities of $`v_{_{\mathrm{LSR}}}250`$ km s<sup>-1</sup> have been identified within $`15^{}`$ of the galaxy NGC 6946. More might be discovered as the DOGS data in this region have not yet been fully analyzed (Fig. 10). The agglomeration of these various galaxies might indicate a group or cloud of galaxies in the nearby Universe. As such it would be the only galaxy group in the nearby Universe that is strongly offset (by $`40^{}`$) from the Supergalactic Plane (Tammann & Kraan-Korteweg 1978, Kraan-Korteweg 1979).
The most significant nearby, previously unknown galaxy identified with DOGS was Dwingeloo 1. Given the 80% coverage of the survey region by the shallow survey (Henning et al. 1998), chances are low that a massive nearby spiral was missed, since nearby galaxies appear in many adjacent pointings, all of which would have to be missed for the galaxy to escape detection. Thus, it is fairly unlikely that there exists another previously unidentified massive spiral galaxy in the area covered by the survey.
### 8.2. The Parkes Multibeam ZOA blind HI survey
In March 1997, the systematic blind H I survey in the southern Milky Way ($`212{}_{}{}^{}\mathrm{}36^{}`$; $`|b|5.^{}5`$) began with the Multibeam receiver at the 64 m Parkes telescope. The instrument has 13 beams, each detecting orthogonal linear polarization. The beams of $`\mathrm{FWHP}=14.^{}4`$ are arranged in a hexagonal grid in the focal plane (cf. Staveley-Smith et al. 1996), allowing rapid sampling of large areas. The average system temperature is about 20 K.
The observations are being performed in driftscan mode. 23 contiguous fields of length $`\mathrm{\Delta }\mathrm{}=8^{}`$ have been defined. Each field is being surveyed along constant Galactic latitudes with latitude offsets of 35 arcmin until the final width of $`|b|5.^{}5`$ has been attained (17 passages back and forth). The ultimate goal is to have 25 repetitions per field where each repetition will furthermore be offset in latitude by $`\mathrm{\Delta }b=1.^{}5`$ for homogeneous sampling. With an effective integration time of 25 min/beam, a 3 $`\sigma `$ detection limit of 25 mJy is obtained. The survey covers the velocity range $`1200<v<12700`$ km s<sup>-1</sup> with a channel spacing of 13.2 km s<sup>-1</sup> per channel, and will be sensitive to normal spiral galaxies well beyond the Great Attractor region. As a byproduct, the survey will produce a high resolution integrated H I column density map of the southern Milky Way and a detailed catalog of high velocity clouds (cf. Putman et al. 1998).
So far, a shallow survey (next section) covering the whole southern Milky Way, based on 2 out of the 25 foreseen driftscan passages, has been analyzed (cf. Kraan-Korteweg et al. 1998a; Henning et al. 1999, 2000). A detailed study of the Great Attractor region ($`308{}_{}{}^{}\mathrm{}332^{}`$) based on 4 scans has been made (Juraszek 1999; Juraszek et al. 2000). The first four full-sensitivity cubes are available for that region as well (Sect. 8.2.).
#### The Parkes ZOA MB shallow survey
In the shallow survey, 110 galaxies were cataloged with peak H I-flux densities of $`>`$80 mJy ($`\mathrm{rms}=15`$ mJy after Hanning smoothing). The detections show no dependence on Galactic latitude, nor on the amount of foreground obscuration through which they have been detected. Though galaxies up to 6500 km s<sup>-1</sup> were identified, most of the detected galaxies (80%) are quite local ($`v<3500`$ km s<sup>-1</sup>) due to the (yet) low sensitivity. About one third of the detected galaxies have a counterpart either in NED or in the deep optical surveys.
The distribution of the 110 H I-detected galaxies is displayed in the lower panel of Fig. 11. It demonstrates convincingly that galaxies can be traced through the thickest extinction layers of the Galactic Plane. The fact that hardly any galaxies are found behind the Galactic bulge ($`\mathrm{}=350^{}`$ to $`\mathrm{}=30^{}`$) is due to local structure: this is the region of the Local Void (see discussion below and top panel of Fig. 12).
For comparative purposes, the top panel of Fig. 11 shows the distribution of all galaxies with known velocities v $`10000`$ km s<sup>-1</sup> to date (extracted from the LEDA database). Although this constitutes an uncontrolled sample it traces the main structures in the nearby Universe in a representative way. Note the increasing incompleteness for extinction levels of $`A_B>1.^\mathrm{m}0`$ (outer contour) – reflecting the growing incompleteness of standard galaxy catalogs (see Sect. 5. and Fig. 1) – and the almost complete lack of galaxy data for extinction levels $`A_B>3.^\mathrm{m}0`$ (inner contour). The middle panel shows galaxies with v$`<`$10000 km s<sup>-1</sup> from the follow-up observations of the deep optical galaxy search by Kraan-Korteweg and collaborators (Sect. 5.4.). Various new overdensities are apparent at low latitudes but the innermost part of our Galaxy remains obscured with this approach. Here, the blind H I data (lower panel) finally can provide the missing link for LSS studies.
In Fig. 12, the data of Fig. 11 are combined in redshift slices. The achieved sensitivity of the shallow MB H I-survey fills in structures all the way across the ZOA for the upper panel ($`v<`$ 3500 km s<sup>-1</sup>) for the first time. Note the continuity of the thin filamentary sine-wave-like structure that dominates the whole southern sky and crosses the Galactic equator twice. This structure snakes over $`180^{}`$ through the southern sky. Taking a mean distance of $`30h^1`$ Mpc, this implies a linear size of $`100h^1`$ Mpc, with thickness of ’only’ $`5h^1`$ Mpc or less. Various other filaments spring forth from this dominant filament, always from a rich group or small cluster at the junction of these interleaving structures. This feature is very different from the thick, foamy Great Wall-like structure, the GA, in the middle panel.
Also note the prominence of the Local Void which is very well delineated in this presentation. No low redshift galaxies were found within the Local Void. But three newly identified galaxies at $`\mathrm{}30^{}`$ help to define the boundary of the Void.
The full-sensitivity ZOA MB-survey will fill in the LSS in the more distant panels of Fig. 12. First results of the full sensitivity survey have been obtained in the Great Attractor region (next section).
Three nearby, very extended ($`20^{}`$ to $`>1^{}`$) galaxies were discovered in the shallow survey. Being likely candidates of dynamically important galaxies, immediate follow-up observations were initiated at the ATCA. These objects did not turn out to be massive perturbing monsters, however. Two were seen to break up into H I complexes and have unprecedented low H I column densities (Staveley-Smith et al. 1998). Systematic synthesis observations are being performed to investigate the frequency of these interacting and/or low H I column density systems in this purely H I-selected sample.
#### The Parkes ZOA MB deep survey
Four cubes centered on the Great Attractor region ($`300{}_{}{}^{}\mathrm{}332^{}`$, $`|b|5.^{}5`$) of the full-sensitivity survey have been analyzed (Juraszek et al. 2000). 236 galaxies above the $`3\sigma `$ detection level of 25 mJy have been uncovered. 70% of the detections had no previous identification.
In the left panel of Fig. 13, a sky distribution centered on the GA region displays all galaxies with redshifts $`v10000`$ km s<sup>-1</sup>. Next to redshifts from the literature (LEDA), redshifts from the follow-up observations of Kraan-Korteweg and collaborators in the Hyd/Ant-Crux-GA ZOA surveys (dashed area) are plotted. They clearly reveal the prominence of the cluster A3627 at $`(\mathrm{},b,v)=(325{}_{}{}^{},7{}_{}{}^{},4848`$ km s<sup>-1</sup>; Kraan-Korteweg et al. 1996) close to the core of the GA region at $`(320{}_{}{}^{},0{}_{}{}^{},4500`$ km s<sup>-1</sup>), (Kolatt et al. 1995). Adding now the new detections from the systematic blind H I MB-ZOA survey (box), structures can be traced all the way across the Milky Way. The new picture seems to support that the GA overdensity is a “great-wall” like structure starting close to the Pavo cluster, having its core at the A3627 cluster and then bending over towards shorter longitudes across the ZOA.
This becomes even clearer in the right panel of Fig. 13 where the galaxies are displayed in a redshift cone out to $`v10000`$ km s<sup>-1</sup> for the longitude range $`300{}_{}{}^{}\mathrm{}332^{}`$ analyzed so far of the MB full-sensitivity data. The A3627 cluster is clearly the most massive galaxy cluster uncovered by the combined surveys in the GA region and therefore the most likely candidate for the previously unidentified but predicted density-peak at the bottom of the potential well of the GA overdensity.
The new data do not unambigously confirm the existence of the suspected further cluster around the bright elliptical radio galaxy PKS1343$``$601 (Sect. 5.4.). Although the MB data reveal an excess of galaxies at this position in velocity space ($`b=+2{}_{}{}^{},v=4000`$ km s<sup>-1</sup>) a “finger of God” is not seen. It could be that many central cluster galaxies are missed by the H I observations because spiral galaxies generally avoid the cores of clusters. The reality of this possible cluster still remains a mystery. A first glimpse of the $`I`$-band images obtained by Woudt et al. (in progress) reveal various early-type galaxies. The forthcoming analysis should then unambiguously settle the question whether another cluster forms part of the GA overdensity.
### 8.3. Conclusions
The systematic probing of the galaxy distribution in the most opaque parts of the ZOA with H I surveys have proven very powerful. For the first time LSS could be mapped without hindrance across the Milky Way (Figs. 10, 12 and 13). This is the only approach that easily uncovers the galaxy distribution in the ZOA, allows the confirmation of implied connections and uncovers new connections behind the Milky Way.
From the analysis of the Dwingeloo survey and the shallow Parkes MB ZOA survey, it can be maintained that no Andromeda-like or other H I-rich galaxy is lurking undetected behind the deepest extinction layers of the Milky Way (although gas-poor, early-type galaxies might, of course, still remain hidden, like the recently discovered very local dwarf behind the Galactic bulge, Sect. 5.5.). The census of dynamically important, H I-rich nearby galaxies whose gravitational influence could significantly impact peculiar motion of the Local Group or its internal dynamics is now complete – at least for objects whose signal is not drowned within the strong Galactic H I emission. Simulations are currently being devised to investigate what kind of H I galaxies – whose signals lie within the frequency range of the Milky Way’s H I – could still have been missed.
## 9. X-ray surveys
The X-ray band potentially is an excellent window for studies of large-scale structure in the Zone of Avoidance, because the Milky Way is transparent to the hard X-ray emission above a few keV, and because rich clusters are strong X-ray emitters. Since the X-ray luminosity is roughly proportional to the cluster mass as $`L_XM^{3/2}`$ or $`M^2`$, depending on the still uncertain scaling law between the X-ray luminosity and temperature, massive clusters hidden by the Milky Way should be easily detectable through their X-ray emission.
This method is particularly attractive, because clusters are primarily composed of early-type galaxies which are not recovered by IRAS galaxy surveys (Sect. 7.) or by systematic H I surveys (Sect. 8.). Even in the NIR, the identification of early-type galaxies becomes difficult or impossible at the lowest Galactic latitudes because of the increasing extinction and crowding problems (Sect. 6.). Rich clusters, however, play an important role in tracing large-scale structures because they generally are located at the center of superclusters and Great Wall-like structures. They mark the density peaks in the galaxy distribution and – with the very high mass-to-light ratios of clusters – the deepest potential wells within these structures. Their location within these overdensities will help us understand the observed velocity flow fields induced by these overdensities.
The X-ray all-sky surveys carried out by Uhuru, Ariel V, HEAO-1 (in the $`210`$ keV band) and ROSAT ($`0.12.4`$ keV) provide an optimal tool to search for clusters of galaxies at low Galactic latitude. However, confusion with Galactic sources such as X-ray binaries and Cataclysmic Variables may cause serious problems, especially in the earlier surveys (Uhuru, Ariel V and HEAO-1) which had quite low angular resolution. And although dust extinction and stellar confusion are unimportant in the X-ray band, photoelectric absorption by the Galactic hydrogen atoms – the X-ray absorbing equivalent hydrogen column density – does limit detections close to the Galactic Plane. The latter effect is particularly severe for the softest X-ray emission, as e.g. observed by ROSAT ($`0.12.4`$ keV) compared to the earlier $`210`$ keV missions. On the other hand, the better resolution of the ROSAT All Sky Survey (RASS) compared to the HEAO-1 survey will reduce confusion problems with Galactic sources as happened, for example, in the case of the cluster A3627 (see below).
Until recently, the possibility of searching for galaxy clusters behind the Milky Way through their X-ray emission has not been pursued in a systematic way, even though a large number of X-ray bright clusters are located at low Galactic latitudes (Fabian 1994): for instance, four of the seven most X-ray luminous clusters in the 2–10 keV range, the Perseus, Ophiuchus, Triangulum Australis, and PKS 0745$``$191 clusters ($`L_\mathrm{X}>10^{45}`$ erg s<sup>-1</sup>) lie at latitudes below $`|b|<20^{}`$ (Edge et al. 1990).
A first attempt to identify galaxy clusters in the ZOA through their X-ray emission had been made by Jahoda and Mushotzky in 1989. They used the HEAO-1 all-sky data to search for X-ray-emission of a concentration of clusters or one enormous cluster that might help explain the shortly before discovered large-scale deviations from the Hubble flow that were associated with the Great Attractor. Unfortunately, this search missed the 6<sup>th</sup> brightest cluster A3627 in the ROSAT X-ray All Sky Survey (Böhringer et al. 1996, Tamura et al. 1998) which had been identified as the most likely candidate for the predicted but unidentified core of the Great Attractor (Kraan-Korteweg et al. 1996). A3627 was not seen in the HEAO-1 data because of the low angular resolution and the confusion with the neighbouring X-ray bright, Galactic X-ray binary 1H1556$``$605 (see Fig. 8 and 9 in Böhringer et al. 1996).
### 9.1. CIZA: Clusters in the Zone of Avoidance
Since 1997, a group led by Ebeling (Ebeling et al. 1999, 2000) have systematically searched for bright X-ray clusters of galaxies at $`|b|<20^{}`$. Starting from the ROSAT Bright Source Catalog (BSC, Voges et al. 1999) which lists the 18811 X-ray brightest sources detected in the RASS, they apply the following criteria to search for clusters: (a) $`|b|<20^{}`$, (b) an X-ray flux above $`S>5\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup> (the flux limit of completeness of the ROSAT BCS), and (c) a spectral hardness ratio. Ebeling et al. demonstrated in 1998 that the X-ray hardness ratio is very effective in discriminating against softer, non-cluster X-ray sources. With these criteria they selected a candidate cluster sample which, although at this point still highly contaminated by non-cluster sources, contains the final CIZA cluster sample.
They first cross-identified their 520 cluster candidates against NED and SIMBAD, and checked unknown ones on the Digitized Sky Survey. The new cluster candidates, including known Abell clusters without photometric and spectroscopic data, were imaged in the $`R`$ band, respectively in the $`K^{}`$ band in regions of high extinctions. With the subsequent spectroscopy of galaxies around the X-ray position, the real clusters could be confirmed.
Time and funding permitting, the CIZA team plans to extend their cluster survey to lower X-ray fluxes ($`23\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup>), the aim being a total sample of 200 X-ray selected clusters at $`|b|<20^{}`$.
So far, 76 galaxy clusters were identified within $`|b|<20^{}`$ of which 80% were not known before. Their distribution is displayed in Fig. 14 (reproduced from Ebeling et al. 2000). 14 of these clusters are relatively nearby ($`z0.04`$), and one was uncovered within the Perseus-Pisces chain at a latitude of only $`b=0.^{}3`$.
### 9.2. Conclusions
With the discovery of 76 clusters so far, of which only 20% were known before, Ebeling et al. (2000) have proven the strength of the method to use X-ray criteria to search for galaxy clusters in the ZOA. As mentioned in the introduction to this section, this approach is complementary to searches at other wavelengths which all fail to uncover galaxy clusters at very low Galactic latitudes.
Having used the ROSAT BSC to select their cluster candidates, the CIZA collaboration wants to combine their final cluster sample with other X-ray selected cluster samples from the RASS, such as the ROSAT Brightest Cluster Sample at $`|b|20^{}`$ and $`\delta 0^{}`$ (Ebeling et al. 1998) and the REFLEX sample at $`|b|20^{}`$ and $`\delta 2.5^{}`$ (Böhringer et al. in prep.). The resulting, all-sky cluster list will be ideally suited to study large-scale structure and the connectivity of clusters across the Galactic Plane.
## 10. Statistical reconstructions
Where the Zone of Avoidance cannot be observed directly, the alternative is to reconstruct the structure in a statistical way. Corrections for unobserved regions in catalogs were done, somewhat ad-hoc, by populating the ZOA uniformly according to the mean density, or by interpolating the structure below and above the Galactic Plane (e.g. Lynden-Bell et al. 1989; Yahil et al. 1991; Strauss et al. 1992; Hudson 1992). Other authors utilized statistical methods such as Wiener filtering (see below) to recover an all-sky density field (Lahav et al. 1994; Hoffman 1994; Fisher et al. 1995; Zaroubi et al. 1995). The recovery of a signal from noisy and incomplete data is a classic problem of inversion, common in problems of image processing. A straightforward inversion is often unstable, and a regularization scheme of some sort is essential in order to interpolate where data are missing or noisy. In the Bayesian spirit, one can use raw data and a prior model to produce an “optimal reconstruction”. Using the above principle, one can derive the Wiener filter (the ratio of signal to signal+noise), which also follows from requiring minimum variance (e.g. Rybicki & Press 1992). The formalism of Wiener filtering is given in Appendix A.
### 10.1. Mask inversion using Wiener filtering in spherical harmonic analysis
It is convenient to expand the galaxy distribution in a nearly whole-sky survey in spherical harmonics. This was applied to 2-D (i.e. projected on the sky) samples (e.g. Peebles 1973; Scharf et al. 1992) and to redshift and peculiar velocity surveys (e.g. Regös & Szalay 1989; Scharf & Lahav 1993; Lahav et al. 1994; Fisher et al. 1994; Nusser & Davis 1994; Fisher et al. 1995; Heavens & Taylor 1995). In projection, the density field over $`4\pi `$ is expanded as a sum:
$$𝒮(\theta ,\varphi )=\underset{l}{}\underset{m=l}{\overset{m=+l}{}}a_{lm}Y_{lm}(\theta ,\varphi ),$$
where the $`Y_{lm}`$’s are the orthonormal set of spherical harmonics and $`\theta `$ and $`\varphi `$ are the spherical polar angles.
The problem of reconstructing large-scale structure behind the ZOA can be formulated as follows: what are the full-sky, noise-free harmonic coefficients $`a_{lm}`$ given the observed harmonics, the mask describing the unobserved region, and a prior model for the power-spectrum of fluctuations? In a projected catalog the observed harmonics $`c_{lm,obs}`$ (with the masked regions filled in uniformly according to the mean) are related to the underlying ‘true’ whole-sky harmonics $`a_{lm}`$ by (see Peebles 1980)
$$c_{lm,obs}=\underset{l^{}}{}\underset{m^{}}{}W_{ll^{}}^{mm^{}}[a_{l^{}m^{}}+\sigma _a],$$
where the monopole term ($`l^{}=0`$) is excluded. The Poisson shot-noise $`\sigma _a`$ is added to the number-weighted harmonics $`a_{lm}`$’s. The noise variance is estimated as $`\sigma _a^2=N`$ (the mean number of galaxies per steradian, independent of $`l`$) The harmonic transform of the mask, $`W_{ll^{}}^{mm^{}}`$, introduces ‘cross-talk’ between the different harmonics.
It can be shown (Lahav et al. 1994; Zaroubi et al. 1995) that the solution of this inversion problem is
$$\widehat{𝐚}=\mathrm{𝐅𝐖}^1𝐜_{obs},$$
where the vectors $`𝐚`$ and $`𝐜_{obs}`$ represent the sets of harmonics $`\{a_{lm}\}`$ and $`\{c_{lm,obs}\}`$, with the diagonal Wiener matrix
$$𝐅=diag\left\{\frac{a_l^2_{th}}{a_l^2_{th}+\sigma _a^2}\right\}.$$
Here $`a_l^2_{th}`$ is the cosmic variance in the harmonics, which depends on the power-spectrum. In the special case of an underlying Gaussian field the most probable field, the mean field, and the minimum variance Wiener filter all are identical. The scatter in the reconstruction can be written analytically for Gaussian random fields.
Even if the sky coverage is $`4\pi `$ ($`𝐖=𝐈`$), the Wiener filter is essential to reveal the optimal underlying continuous density field, cleaned of noise. In the absence of other prior information on the location of clusters and voids, the correction factor is ‘isotropic’ per $`l`$, i.e. independent of $`m`$. So, in the case of full sky coverage, only the amplitudes are affected by the correction, but not the relative phases. For example, the dipole direction is not affected by the shot-noise, only its amplitude. But if the sky coverage is incomplete, both the amplitudes and the phases are corrected. The reconstruction also depends on number of observed and desired harmonics. Note also that the Wiener method is non-iterative.
### 10.2. Reconstruction of the projected IRAS 1.2 Jy galaxy distribution
Lahav et al. (1994) applied the method to the sample of IRAS galaxies brighter than 1.2 Jy which includes 5313 galaxies, and covers 88% of the sky. This incomplete sky coverage is mainly due to the Zone of Avoidance, which they modelled as a ‘sharp mask’ at Galactic latitude $`|b|<5^{}`$. The mean number of galaxies is $`N400`$ per steradian, which sets the shot-noise, $`\sigma _a^2`$. As the model for the cosmic scatter $`a_l^2_{th}`$, they adopted a fit to the observed power spectrum of IRAS galaxies.
Fig. 16 shows the reconstruction of the projected IRAS 1.2 Jy sample. The ZOA was left empty, and it clearly ‘breaks’ the possible chain of the Supergalactic Plane and other structures. Fig. 16 shows our optimal reconstruction for $`1l15`$. Now the structure is seen to be connected across the ZOA, in particular in the regions of Centaurus/Great Attractor ($`\mathrm{}315^{}`$), Hydra ($`\mathrm{}275^{}`$) and Perseus-Pisces ($`\mathrm{}315^{}`$), confirming the connectivity of the Supergalactic Plane. Note that Figs. 16 and 16 are shifted by $`180^{}`$ compared to the earlier sky projections. We also see the Puppis cluster ($`\mathrm{}245^{}`$) recovered behind the Galactic Plane. This cluster was noticed in earlier harmonic expansion (Scharf et al. 1992) and further studies (Lahav et al. 1993). The other important feature of this reconstruction is the suppression of shot noise all over the sky. This is particularly important for judging the reality of clusters and voids.
Comparison of this Wiener reconstruction with the one applied (using a $`4\pi `$ Wiener filter) to the IRAS sample, in which the ZOA was filled in ‘by hand’ across the Galactic Plane (Yahil et al. 1991), shows good agreement. By testing the method on $`N`$-body simulations (where the whole ‘sky’ true harmonics are known) it was found, that for masks larger than $`|b|=15^{}`$, it is difficult to recover the unobserved structure. In this case extra-regularization was required. The success of the method depends on the interplay of three angular scales: the width of the mask, the desired resolution ($`\pi /l_{max}`$) and the physical correlation of structure. It is important to note that, as opposed to ad-hoc smoothing schemes, the smoothing due to a Wiener filter is determined by the sparseness of data relative to the expected signal.
### 10.3. 3-D reconstruction
Fisher et al. (1995) generalized the 2-D Wiener reconstruction to 3-D (i.e. for redshift catalogs) by expanding the density field $`\rho `$ in redshift-space in terms of spherical harmonics, $`Y_{lm}`$, and radial Bessel functions $`j_l`$:
$$\rho (𝐬)=\underset{l=0}{\overset{l_{\mathrm{max}}}{}}\underset{m=l}{\overset{+l}{}}\underset{n=1}{\overset{n_{\mathrm{max}}(l)}{}}C_{ln}\rho _{lmn}j_l(k_ns)Y_{lm}(\widehat{𝐬}).$$
The discrete $`k_n`$’s are chosen according to the boundary conditions, so as to make the set orthogonal. This process is analogous to Fourier decomposition, but uses instead a set of spherical basis functions. The data from a redshift catalog can be seen as a set of $`N`$ discrete points, $`𝐬_i`$, each giving the direction and redshift of a galaxy. These are used to estimate the underlying density field in redshift space, $`\widehat{\rho }^\mathrm{S}(𝐬)`$, expanded as in the above equation. Here, $`C_{ln}`$ are normalization constants, while the harmonic coefficients are given by
$$\widehat{\rho }_{lmn}^\mathrm{S}=\underset{i=1}{\overset{N}{}}\frac{1}{\varphi (s_i)}j_l(k_ns_i)Y_{lm}^{}(\widehat{𝐬}_𝐢),$$
where $`\varphi (s_i)`$ is the spherical selection function of the survey, evaluated at the radius of the $`i^{th}`$ galaxy. The real-space density, velocity and potential fields are then reconstructed using linear theory and a Wiener filter.
In analogy with the 2-D case, a mask (e.g. to describe the ZOA) can be formulated in the 3-D harmonic formalism (Heavens & Taylor 1995; Schmoldt et al. 2000). However, for the simplicity of the analysis, Fisher et al. (1995) first populated the ZOA with mock galaxies according to the procedure of Yahil et al. (1991), and then applied the reconstruction over the full sky. Webster et al. (1997) and Schmoldt et al. (2000) extended this analysis and presented detailed maps of the reconstructed fields, as well as optimal determinations of the Local Group dipole and bulk flows. Lahav et al. (2000) used Wiener filtering to study the the extent of the Supergalactic Plane. Bistolas (1998) has done a similar Wiener reconstruction in Cartesian coordinates (see discussion below). Saunders et al. (2000a) have extended the spherical harmonics and Wiener approach to represent non-linear clustering by describing the density field as drawn from a log-normal probability distribution function. They recently applied it to the PSCz IRAS catalog (taking into account the detailed IRAS mask, see Fig. 8). This type of analysis, when applied to other surveys at low latitude, can potentially provide the most detailed and self-consistent map of the ZOA.
### 10.4. POTENT reconstruction of the ZOA
The POTENT method (Bertschinger & Dekel 1989; Dekel 1994) recovers the smoothed fluctuations field of potential, velocity and mass density from observed radial peculiar velocities of galaxies. The velocity field is recovered under the assumption of potential flow, $`𝐯(𝐫)=\mathrm{\Phi }(𝐫)`$. The potential can thus be calculated by integrating the radial velocity along radial rays. Differentiating $`\mathrm{\Phi }`$ in the transverse directions recovers the two missing velocity components. The underlying mass-density fluctuation $`\delta `$ is then derived in linear theory from $`\delta (𝐫)=\mathrm{\Omega }^{0.6}𝐯(𝐫)`$ (or from non-linear extensions). This method is very useful for exploring the ZOA as the peculiar velocity field responds to the entire mass distribution, regardless of the ‘unseen’ distribution of light. However, due to the heavy smoothing required by this method, only structures on large scales (e.g. superclusters) can be mapped. Individual (massive) nearby galaxies that may perturb the dynamics in the vicinity of the Local Group cannot be uncovered in this manner. Kolatt et al. (1995) used the POTENT method to specifically predict the mass distribution behind the ZOA (see Fig. 18). Some of their predictions are summarized below.
Zaroubi et al. (1999) analyzed the peculiar velocity field using Wiener filtering. The reconstructed structures are consistent with those extracted by the POTENT method. A comparison with the structures in the distribution of IRAS 1.2 Jy galaxies yields a general agreement.
Figure 17 shows an Aitoff projection in Galactic coordinates of reconstrucion at $`R=40h_1^100`$ Mpc from the three velocity catalogs: Mark III (Willick et al. 1997), SFI (Survey of Field Spirals in the I-band; da Costa et al. 1996; Giovanelli et al. 1998) and ENEAR (the ESO Nearby Early-type Galaxies Survey; da Costa et al. 2000a, 2000b). For further discussion of the Wiener reconstruction of the ZOA from velocities see Hoffman (2000).
### 10.5. Predicted structures behind the ZOA
Early reconstructions on relatively sparse data galaxy catalogs have been performed within volumes out to $`v`$ 5000 km s<sup>-1</sup>. Despite heavy smoothing, they have been quite successful in pinpointing a number of important features.
Scharf et al. (1992) applied spherical harmonics to the 2-dimensional IRAS PSC and noted a prominent cluster behind the ZOA in Puppis ($`\mathrm{}245^{}`$) which was simultanously discovered as a nearby cluster through H I-observations of obscured galaxies in that region by Kraan-Korteweg & Huchtmeier (1992). It was analyzed further by Lahav et al. (1993).
Hoffman (1994) predicted the Vela supercluster at ($`280{}_{}{}^{},6{}_{}{}^{},6000`$ km s<sup>-1</sup>), using 3-dimensional Wiener filter reconstructions on the IRAS 1.9 Jy redshift catalog (Strauss et al. 1992). It was discovered observationally just a bit earlier by Kraan-Korteweg & Woudt (1993).
Using POTENT analysis, Kolatt et al. (1995) predicted the center of the Great Attractor overdensity – its density peak – to lie behind the ZOA at ($`320{}_{}{}^{},0{}_{}{}^{},4500`$ km s<sup>-1</sup>; see Fig. 18). Shortly thereafter, Kraan-Korteweg et al. (1996) unveiled the cluster A3627 as being very rich and massive and at the correct distance. It hence is the most likely candidate for the central density peak of the GA.
POTENT reconstructions have been applied to denser galaxy samples covering larger volumes (8-10000 km s<sup>-1</sup>) with smoothing scales of the order of 500 km s<sup>-1</sup> (compared to 1200 km s<sup>-1</sup>). It therefore seemed of interest to see whether these reconstructions find evidence for unknown major galaxy structures in the ZOA at higher redshifts.
The currently most densely-sampled, well-defined galaxy redshift catalog is the Optical Redshift Survey (Santiago et al. 1995). However, this catalog is limited to $`|b|20^{}`$ and the reconstructions (see Baker et al. 1998) within the ZOA are strongly influenced by 1.2 Jy IRAS Redshift Survey data and a mock galaxy distribution in the inner ZOA. We will therefore concentrate on reconstructions based on the 1.2 Jy IRAS Redshift Survey only.
In the following, the structures identified in the ZOA by (a) Webster et al. (1997) using Wiener filter plus spherical harmonics and linear theory and (b) Bistolas (1998), who applied a Wiener filter plus linear theory and constrained realizations, will be discussed and compared to observational data. Figure 2 in Webster et al. displays the reconstructed density fields on shells of 2000, 4000, 6000 and 8000 km s<sup>-1</sup>; Fig. 5.2 in Bistolas displays the density fields in the ZOA from 1500 to 8000 km s<sup>-1</sup> in steps of 500 km s<sup>-1</sup>.
The reconstructions by Webster et al. 1997 clearly show the recently identified nearby cluster at ($`33{}_{}{}^{},5{}_{}{}^{}15^{}`$,1500 km s<sup>-1</sup>; see Sect. 5.4.), whereas Bistolas reveals no clustering in the region of the Local Void out to 4000 km s<sup>-1</sup>. At the same longitudes, the clustering at 7500 km s<sup>-1</sup> is seen by Bistolas, but not by Webster et al. The Perseus-Pisces chain is strong in both reconstructions, and the second Perseus-Pisces arm – which folds back at $`\mathrm{}95^{}`$ – is clearly confirmed. Both reconstructions find the Perseus-Pisces complex to be very extended in space, i.e. from 3500 km s<sup>-1</sup> out to 9000 km s<sup>-1</sup>. Whereas the GA region is more prominent compared to Perseus-Pisces in the Webster et al. reconstructions, the signal of the Perseus-Pisces complex is considerably stronger than the GA in Bistolas, where it does not even reveal a well-defined central density peak. Both reconstructions find no evidence for the suspected cluster around PKS1343$``$601, but its signal could be hidden in the central (A3627) density peak due to the smoothing. While the Cygnus-Lyra complex ($`60{}_{}{}^{}90{}_{}{}^{},0{}_{}{}^{},4000`$ km s<sup>-1</sup>) discovered by Takata et al. (1996) stands out clearly in Bistolas, it is not evident in Webster et al. Both reconstructions find a strong signal for the Vela SCL ($`285{}_{}{}^{},6{}_{}{}^{},6000`$ km s<sup>-1</sup>), labelled as HYD in Webster et al. 1997. The Cen-Crux cluster identified by Woudt (1998) is evident in Bistolas though less distinct in Webster et al. A suspected connection at ($`\mathrm{},v)(345{}_{}{}^{},6000`$ km s<sup>-1</sup>; see middle panel of Fig. 12) is supported by both methods. The Ophiuchus cluster just becomes visible in the most distant reconstruction shells (8000 km s<sup>-1</sup>).
### 10.6. Conclusions
Not all reconstructions find the same features, and when they do, the prominence of the density peaks as well as their locations in space do vary considerably. At velocities of $`4000`$ km s<sup>-1</sup> most of the dominant structures happen to lie close to or within the ZOA while at larger distances, clusters and voids seem to be more homogeneously distributed over the whole sky. Out to 8000 km s<sup>-1</sup> none of the reconstructions predict any major structures which are not mapped or suggested from observational data. Thus, no major surprises seem to remain hidden in the ZOA. The various multi-wavelength explorations of the Milky Way will soon be able to verify this. Still, the combination of both the reconstructed potential fields and the observationally mapped galaxy distribution will lead to estimates of the cosmological parameters $`\mathrm{\Omega }_0`$ and $`b`$.
## 11. Discussion
In the last decade, enormous progress has been made in unveiling the extragalactic sky behind the Milky Way. At optical wavebands, the entire ZOA has been systematically surveyed. It has been shown that these surveys are complete for galaxies larger than $`D^o=1.^{}3`$ (corrected for absorption) down to extinction levels of $`A_B=3.^\mathrm{m}0`$. Combining these data with previous “whole-sky” maps reduces the “optical ZOA” by a factor of about 2 - 2.5, which allows an improved understanding of the velocity flow fields and the total gravitational attraction on the Local Group. Various previously unknown structures in the nearby Universe could be mapped in this way.
At higher extinction levels, other windows to the ZOA become more efficient in tracing the large-scale structures. Very promising in this respect are the current near-infrared surveys which find galaxies down to latitudes of $`|b|1.^{}5`$. Source confusion will remain a problem at low Galactic latitude which may be overcome by introducing novel statistical methods such as Artificial Neural Networks. The systematic H I surveys detect gas-rich spiral galaxies all the way across the Galactic Plane – slightly hampered only at very low latitudes ($`|b|<1.^{}0`$) because of the numerous continuum sources. These studies have already shown that no unknown dynamically important, H I-rich nearby galaxies whose gravitational influence could significantly impact the internal dynamics and the peculiar motion of the Local Group are hidden by the Milky Way. In addition, the deep ZOA H I surveys can be merged with the lower-sensitivity whole-sky H I-surveys currently in progress. The “Behind the Plane” survey resulted in a reduction from 16% to 7% of the “FIR ZOA” which soon should provide improved values of the dipole direction and convergence from IRAS data. In addition, new indications of possible hidden massive clusters behind the Miky Way are now forthcoming from the CIZA project – although again an “X-ray ZOA” will remain due to the absorption of X-ray radiation by the thick gas layer close to the Galactic Plane.
A difficult task is still awaiting us, i.e. to obtain a detailed understanding of the selection effects inherent to the various methods. Quantifying the selection effects is crucial for any optimal reconstruction method (e.g. Wiener) which attempts to merge the different data sets in an unbiased way. This is extremely important if we want to use these data for quantitative cosmography. Moreover, we need a better understanding of the effects of obscuration on the observed properties of galaxies identified through the dust layer (at all wavelengths), in addition to an accurate high-resolution, well-calibrated map of the Galactic extinction.
Despite the fact that our knowledge about the above issues is as yet limited, a lot can and has been learned from ZOA research. This is evident, for instance, from the detailed and varied investigations of the Great Attractor region. Mapping the GA and understanding the massive overdensity inferred from peculiar velocity fields had remained an enigma due the fact that the major and central part of this extended density enhancement was largely hidden by the obscuring veil of the Milky Way. The results from the various ZOA surveys now clearly imply that the Great Attractor is, in fact, a nearby “Great-Wall” like supercluster, starting at the nearby Pavo cluster below the GP, moving across the massive galaxy cluster A3627 toward the shallow overdensity in Vela at 6000 km s<sup>-1</sup>. The cluster A3627 is the dominant central component of this structure, similar to the Coma cluster in the (northern) Great Wall. Whether a second massive cluster around PKS1343$``$601 is part of the core of the GA remains uncertain.
##### Acknowledgments.
The enthusiastic collaborations of our colleagues in the exploration of the galaxy distribution behind the Milky Way is greatly appreciated. These are P.A. Woudt, C. Salem and A.P. Fairall with deep optical searches, C. Balkowski, V. Cayatte, A.P. Fairall, P.A. Henning with redshift follow-ups of optically identified galaxies, A. Schröder and G.A. Mamon in the exploration of DENIS images at low Galactic latitude, W.B. Burton, P.A. Henning, and A. Rivers in the northern ZOA HI-survey (DOGS) and the HIPASS ZOA team members L. Staveley-Smith , R.D. Ekers, A.J. Green, R.F. Haynes, P.A. Henning, S. Juraszek, M. J. Kesteven, B. Koribalski, R.M. Price, E. Sadler and A. Schröder in the southern ZOA survey.
Particular thanks go to P.A. Woudt for his valuable suggestions, to W. Saunders for preparing Fig. 8, to A. Schröder and G. Mamon for their comments on the NIR section, and to H. Ebeling for his input with regard to the X-ray section and Fig. 14.
This research has made use of the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, Caltech, under contract with the National Aeronautics and Space Administration, as well as the Lyon-Meudon Extragalactic Database (LEDA), supplied by the LEDA team at the Centre de Recherche Astronomique de Lyon, Observatoire de Lyon.
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Appendix A: Wiener Filtering
Here we give a brief review of the Wiener (1949) filter (WF) technique; the reader is referred to Lahav etal. (1994), Zaroubi et al. (1995) and Rybicki & Press (1992) for further details. Let us assume that we have a set of measurements, $`\{d_\alpha \}(\alpha =1,2,\mathrm{}N)`$ which are a linear convolution of the true underlying signal, $`s_\alpha `$, plus a contribution from statistical noise, $`ϵ_\beta `$, such that
$$d_\alpha =_{\alpha \beta }\left[s_\beta +ϵ_\beta \right],$$
$`(A1)`$
where $`_{\alpha \beta }`$ is the response or “point spread” function (summation convention assumed). Notice that we have assumed that the statistical noise is present in the underlying field and therefore is convolved by the response function.
The WF is the linear combination of the observed data which is closest to the true signal in a minimum variance sense. More explicitly, the WF estimate is given by $`s_\alpha (WF)=F_{\alpha \beta }d_\beta `$ where the filter is chosen to minimize $`|s_\alpha (WF)s_\alpha |^2`$. It is straightforward to show that the WF is given by
$$F_{\alpha \beta }=s_\alpha d_\gamma d_\gamma d_\beta ^{}^1,$$
$`(A2)`$
where
$$s_\alpha d_\beta ^{}=_{\beta \gamma }s_\alpha s_\gamma ^{},$$
$`(A3)`$
$$d_\alpha d_\beta ^{}=_{\alpha \gamma }_{\beta \delta }\left[s_\gamma s_\delta ^{}+ϵ_\gamma ϵ_\delta ^{}\right].$$
$`(A4)`$
In the above equations, we have assumed that the signal and noise are uncorrelated. From equation A4, it is clear that, in order to implement the WF, one must construct a prior which depends on the variance of the signal and noise.
The dependence of the WF on the prior can be made clear by defining signal and noise matrices given by $`S_{\alpha \beta }=s_\alpha s_\beta ^{}`$ and $`N_{\alpha \beta }=ϵ_\alpha ϵ_\beta ^{}`$. With this notation, we can rewrite equation A4 as
$$𝐬(WF)=𝐒\left[𝐒+𝐍\right]^1^1𝐝.$$
$`(A5)`$
Formulated in this way, we see that the purpose of the WF is to attenuate the contribution of low signal-to-noise ratio data and therefore regularize the inversion of the response function. The derivation of the WF given above follows from the sole requirement of minimum variance and requires only a model for the variance of the signal and noise. The WF can also be derived using the laws of conditional probability if the underlying distribution functions for the signal and noise are assumed to be Gaussian; in this more restrictive case, the WF estimate is, in addition to being the minimum variance estimate, also both the maximum a posterior estimate and the mean field. For Gaussian fields, the mean WF field can be supplemented with a realization of the expected scatter about the mean field to create a realization of the field; this is the heart of the “constrained realization” approach described in Hoffman & Ribak (1991). A generalization to non-Gaussian fields is given in Sheth (1995).
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# L-functions and Random Matrices
## 1 The GUE Conjecture
### 1.1 Introduction
In 1972 H. L. Montgomery announced a remarkable connection between the distribution of the zeros of the Riemann zeta-function and the distribution of eigenvalues of large random Hermitian matrices. Since then a number of startling developments have occurred making this connection more profound. In particular, random matrix theory has been found to be an extremely useful predictive tool in the theory of L-functions. In this article we will try to explain these recent developments and indicate some directions for future investigations.
### 1.2 The Riemann zeta-function
The Riemann zeta-function is defined by
$$\zeta (s)=\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n^s}(s=\sigma +it,\sigma >1).$$
(1)
It can also be expressed as a product over primes, the Euler product,
$$\zeta (s)=\underset{p}{}\left(1\frac{1}{p^s}\right)^1,$$
(2)
for $`\sigma >1`$. In 1859 Riemann proved that $`\zeta (s)`$ extends to a meromorphic function on the whole plane with its only singularity being a simple pole at $`s=1`$ with residue 1. He further proved that it has a functional equation relating the value of $`\zeta (s)`$ with the value of $`\zeta (1s)`$,
$$\zeta (s)=\chi (s)\zeta (1s)$$
(3)
where $`\chi (1s)=\chi (s)^1=2(2\pi )^s\mathrm{\Gamma }(s)\mathrm{cos}(\pi s/2)`$. He discovered that the distribution of prime numbers is governed by the zeros of $`\zeta `$. He was led to conjecture that all of the complex zeros $`\rho =\beta +i\gamma `$ of $`\zeta `$ have $`\beta =1/2`$. This assertion is the famous Riemann Hypothesis.
We know that $`0<\beta <1`$ for any complex zero of $`\zeta `$. Riemann estimated the zero counting function
$$N(T)=\mathrm{\#}\{\rho =\beta +i\gamma :0<\gamma T\}=\frac{T}{2\pi }\mathrm{log}\frac{T}{2\pi e}+O(\mathrm{log}T).$$
(4)
Thus, at a height $`T`$ the average spacing between zeros is asymptotic to $`2\pi /\mathrm{log}T`$. See \[T\] for additional background information about $`\zeta (s)`$.
### 1.3 Background
Montgomery \[M1\] was studying gaps between zeros of the Riemann zeta-function in an attempt to prove that the spacings between consecutive zeros can sometimes be less than 1/2 of the average spacing. Such a conclusion would have led to a good effective lower bound for the class number of an imaginary quadratic field. This estimate was not achieved, but through the course of his analysis Montgomery was led to conjecture that
$$\underset{T\mathrm{}}{lim}\frac{1}{N(T)}\underset{\genfrac{}{}{0pt}{}{0\gamma ,\gamma ^{}T}{\frac{2\pi \alpha }{\mathrm{log}T}\gamma \gamma ^{}\frac{2\pi \beta }{\mathrm{log}T}}}{}1=_\alpha ^\beta \left(1\left(\frac{\mathrm{sin}\pi x}{\pi x}\right)^2\right)𝑑x$$
(5)
where $`0<\alpha <\beta `$ are fixed.
When Montgomery told Freeman Dyson this formula, Dyson responded that the integrand was the pair-correlation function for eigenvalues of large random Hermitian matrices, or more specifically the Gaussian Unitary Ensemble, or GUE.
The GUE is the limit as $`N\mathrm{}`$ of the probability space consisting of $`N\times N`$ Hermitian matrices $`H`$ with a probability measure $`P(H)dH`$ that is invariant under conjugation by any unitary matrix $`U`$. Here $`dH=_{jk}d\mathrm{}H_{jk}_{j<k}d\mathrm{}H_{jk}`$. Mathematical physicists had studied various ensembles since the 1950s in connection with work of Wigner in nuclear physics. Mehta \[Me\] has given a thorough treatment of the development of the subject.
Montgomery went on to conjecture that the $`n`$-correlation function for zeros of $`\zeta `$ is the same as that for the GUE; this conjecture came to be known as Montgomery’s GUE conjecture.
Odlyzko and Schönhage developed an algorithm that allowed for the simultaneous calculation of many values of $`\zeta (1/2+it)`$ for $`t`$ near $`T`$ in average time $`T^ϵ`$. This algorithm allowed Odlyzko \[O\] to do extensive computations of the zeros of $`\zeta `$ at a height near zero number $`10^{20}`$; his computations of the pair correlation and nearest neighbor spacing for the zeros of $`\zeta `$ were amazingly close to those for the GUE. His famous pictures added much credibility to Montgomery’s conjecture.
### 1.4 Further evidence for GUE
In 1996 Rudnick and Sarnak \[RS\] made some interesting progress on the GUE conjecture. To explain their result, number the ordinates of the zeros of $`\zeta (s)`$: $`0<\gamma _1\gamma _2\mathrm{}`$. Introduce a scaling $`\stackrel{~}{\gamma }=\gamma \frac{\mathrm{log}\gamma }{2\pi }`$ so that the $`\stackrel{~}{\gamma }`$ have asymptotic mean spacing 1. Then Sarnak and Rudnick proved that
$$\underset{T\mathrm{}}{lim}\frac{1}{T}\underset{\genfrac{}{}{0pt}{}{\gamma _{j_1},\mathrm{},\gamma _{j_n}T}{j_mj_n}}{}f(\stackrel{~}{\gamma }_{j_1},\mathrm{},\stackrel{~}{\gamma }_{j_n})=_{P_n}W_{U,n}(\stackrel{}{x})f(\stackrel{}{x})𝑑\stackrel{}{x}$$
(6)
where $`W_{U,n}(\stackrel{}{x})=W_{U,n}(x_1,\mathrm{},x_n)`$ is the $`n`$-correlation function for the GUE (see section 5.5) and where $`f`$ is any function satisfying (1) $`f(\stackrel{}{x}+t(1,\mathrm{},1))=f(\stackrel{}{x})`$ for $`tR`$; (2) $`f`$ is smooth and symmetric in the variables and decays rapidly as $`x\mathrm{}`$ in the hyperplane $`P_n:=\{(x_1,\mathrm{},x_n):_{j=1}^nx_j=0\}`$; (3) the Fourier transform $`\widehat{f}(\stackrel{}{u})`$ of $`f`$ is supported in $`_{j=1}^n|u_j|<2`$. The condition (1) assures that $`f`$ is a function of the differences of the $`\gamma _j`$.
Thus, the GUE conjecture has been proven for all correlation functions, for a limited class of test functions
### 1.5 Other L-functions
The Riemann zeta-function is the prototype for some extraordinary objects known as arithmetic L-functions. An arithmetic L-function (or L-function for short) has many properties in common with the Riemann $`\zeta `$-function. It has a Dirichlet series
$$L(s)=\underset{n=1}{\overset{\mathrm{}}{}}a_nn^s;$$
(7)
it is meromorphic apart from a possible pole at $`s=1`$; it has a functional equation $`\gamma (s)L(s)=ϵ\overline{\gamma (1\overline{s})L(1\overline{s})},`$ where $`ϵ`$, the sign, satisfies $`|ϵ|=1`$. $`L(s)`$ has an Euler product in which the $`p`$-th factor is the reciprocal of a polynomial in $`p^s`$. Moreover, an arithmetic L-function is expected to satisfy the Riemann Hypothesis, that all complex zeros have real part 1/2. (Note that our L-functions are normalized so that the 1/2-line is the line of symmetry for the functional equation. See \[S\] for precise definitions.) An L-function is called primitive if it is not the product of two L-functions. Primitive L-functions arise from a variety of contexts: from Dirichlet characters, from Mellin transforms of certain cusp forms, from Galois representations, from algebraic varieties, etc. However, it is believed that each primitive arithmetic L-function is associated to a cuspidal automorphic representation of GL<sub>m</sub> over a number field.
Rudnick and Sarnak applied their methods to a fairly general L-function. They proved the analogue of (6) for any cuspidal automorphic L-function over $`Q`$, assuming the Riemann Hypothesis for the L-function, and with obvious changes to reflect the appropriate scaling of the zeros. In particular, the answer did not depend in any way on the distribution of the coefficients of the particular L-function.
The GUE Conjecture is now seen as a universal law governing the distribution of zero spacings for all arithmetic L-functions.
## 2 Families
The realization that the compact classical groups of matrices play a role in the theory of L-functions arose through seminal work of Katz and Sarnak \[KS1\], \[KS2\].
### 2.1 Function Field Analogues
Katz and Sarnak investigated the distribution of zeros of function field zeta-functions. Consider a curve $`f(x,y)=0`$ where $`f`$ is a polynomial with integer coefficients. The zeta-function for $`f`$ over a finite field $`F_q`$ can be obtained in a simple way from the generating function of the numbers $`N_n`$ of points on that curve in the finite field extensions $`F_{q^n},n=1,2,\mathrm{}.`$ It is known that this zeta-function is a rational function whose numerator is a polynomial with integer coefficients, degree 2$`g`$ where $`g`$ is the genus of the curve, and that it satisfies the Riemann Hypothesis that all zeros have modulus $`1/\sqrt{q}`$. One can order these zeros on the circle $`|z|=\sqrt{q}`$ in terms of their angles measured from the positive axis. One then considers statistics of the angles.
Katz and Sarnak proved that, after proper normalization, the $`n`$-correlation function of these angles, as $`g`$ and $`q`$ tend to infinity, is exactly the $`n`$-correlation of GUE. They also proved that the nearest neighbor, or consecutive spacing, statistic for these angles is the same as that for GUE.
The method of proof involved working with subgroups of U($`N`$), the group of $`N\times N`$ unitary matrices with the Haar measure (see sections 5.1 and 5.2). It had been shown by Dyson that the $`n`$-correlation and the spacing statistics are the same for GUE as for the appropriately scaled limit as $`N\mathrm{}`$ for $`U(N)`$ with Haar measure. (In the physics literature the limit of $`U(N)`$ is called the Circular Unitary Ensemble or CUE; see \[Me\].)
The function field zeta-functions that Katz and Sarnak were working with were characteristic polynomials of matrices from subgroups of U($`N`$). To deduce their result they applied a theorem of Deligne about the equidistribution of function field zeta-functions among the characteristic polynomials of conjugacy classes in these subgroups.
### 2.2 Introduction of Families
Katz and Sarnak investigated the robustness of their theorem. They conjectured that the conclusion remains true under the weaker hypothesis that $`q`$ is held fixed and $`g\mathrm{}`$. On the other hand, they provided examples of sequences of curves of increasing genus for which the spacing statistic is not the GUE consecutive spacing statistic.
They also considered special families such as (a) curves of the form $`y^2=f(x)`$ with all squarefree, monic $`f`$, and (b) quadratic twists of curves of the form $`y^2=x(x1)(xt)`$; that is $`\mathrm{\Delta }y^2=x(x1)(xt)`$ where $`\mathrm{\Delta }`$ is a squarefree, monic polynomial over $`F_q`$. When $`q`$ and the degree of $`\mathrm{\Delta }\mathrm{}`$ they discovered that the $`n`$-correlation and spacing statistics again matched up with CUE. However, the reason in each case is different depending on the geometric monodromy group. In case (a) they computed that the monodromy group was the symplectic group and in case (b) the orthogonal group.
We denote by USp($`N`$) (if $`N`$ is even) the unitary symplectic matrices, and by O($`N`$) the orthogonal matrices. These are subgroups of U($`N`$) and are equipped with their own Haar measure.
Katz and Sarnak proved that the $`n`$-correlation statistic and the spacing statistic for the limits of U($`N`$), USp($`N`$), and O($`N`$) are all the same as that for GUE.
By contrast, the Circular Orthogonal Ensemble (COE) and the Circular Symplectic Ensemble (CSE), which are well-known ensembles in mathematical physics (see \[Me\]) have the same underlying symmetry groups as O and Sp, but have a measure different from the Haar measure. They have different $`n`$-correlation and spacing statistics than those of GUE and CUE.
From now on we will use the letters O, Sp, U when referring to a statistic associated with the limits of O($`N`$), USp($`N`$), and U($`N`$). These statistics are computed for $`N\times N`$ matrices of the appropriate subgroup of U$`(N)`$ with their respective Haar measure and then $`N\mathrm{}`$ with an appropriate scaling limit in order to determine the statistics (see 5.4, 5.5, and 5.6).
Since orthogonal and symplectic matrices have eigenvalues which occur in complex conjugate pairs, it is clear that the eigenvalue 1 plays a special role for these matrices, whereas it does not for unitary matrices. Indeed, there are statistics which differentiate O, Sp, and U. In particular, the eigenvalue nearest to 1 is such an example. It turns out that this statistic is different for all three symmetry types. More generally the $`j`$-th eigenvalue nearest to 1 is a statistic that is dependent on the group (for a summary of the four statistics of interest to us, see section 5.3).
Another statistic that is different for all three symmetry types is “level-density”. The $`n`$-level density function is obtained from summing a test function at $`n`$-distinct eigenvalues. For comparison purposes, note that the $`n`$-correlation function is obtained by summing a test function of $`n`$-variables which is a function depending only on the differences of the arguments at $`n`$ distinct eigenvalues. See sections 5.5 and 5.6 for more on these statistics.
Katz and Sarnak showed that for their special families of function field zeta-functions these new statistics ($`j`$-th nearest eigenvalue to 1 and $`n`$-level density) match with the appropriate statistics from Sp and O, which they had computed.
In general, they found that if they could compute the geometric monodromy group associated with a family of function field zeta-functions, and if that monodromy group was U, Sp or O, then all four of the statistics we have been discussing for the function field zeta-functions could be proven to match with the appropriate statistic from U, Sp or O.
Thus, we say that the symmetry type for the family of curves of the form $`y^2=f(x)`$ the symmetry type is Sp; and for the family of quadratic twists of $`y^2=x(x1)(xt)`$ the symmetry type is O.
Katz and Sarnak also have examples of families of function field zeta-functions where the symmetry type is U.
### 2.3 L-Functions Over Number Fields
Katz and Sarnak speculated that their results for function field zeta-functions would have implications for L-functions over number fields. Two collections of L-functions which present themselves as natural analogues to the families above are (a) the collection of all Dirichlet L-functions $`L(s,\chi _d)`$ where $`\chi _d(n)`$ is a real primitive Dirichlet character with conductor $`|d|`$ (so that $`d`$ runs through the set of fundamental discriminants of quadratic number fields) and (b) the collection of L-functions associated with primitive Hecke newforms of a fixed weight. The analogy with the function field zeta-functions is as follows. The zeta-function for a member of the Sp family (a) above is obtained by counting solutions to the equation $`y^2=f(x)`$ over a finite field. These solutions are counted in a field with $`p`$ elements by the sum $`_{a=1}^p\chi _p(f(a))`$ where the real Dirichlet character $`\chi _p`$ modulo $`p`$ appears. This suggests that the family (a) could have a symmetry type Sp. Similarly, the zeta-function of a member of the O family (b) above is the reciprocal of the $`p`$-th factor in the Euler product for the L-function of the elliptic curve defined through twists of the equation $`y^2=x(x1)(xt)`$, the Legendre family of elliptic curves. These L-functions are known (by the solution of the Taniyama - Shimura conjecture) to be associated to primitive Hecke newforms of weight 2. Thus, the collection (b) could well be a family with symmetry type O.
The first evidence that (a) is a family with symmetry type Sp came from Michael Rubinstein’s thesis \[R\]. He computed the lowest lying zero of $`L(s,\chi _d)`$ and the data matched well with the eigenvalue nearest 1 for symplectic matrices. He examined theoretically the $`n`$-level density of zeros and showed (for test functions with restrictions on the support of their Fourier transforms) that the $`n`$-level density functions for the zeros of $`L(s,\chi _d)`$ are identical with the $`n`$-level density functions for Sp. Precisely, we index the ordinates of the zeros of $`L(s,\chi _d)`$ as
$$0\gamma _1^{(d)}\gamma _2^{(d)}\mathrm{}$$
(8)
and scale using
$$\stackrel{~}{\gamma }^{(d)}=\gamma ^{(d)}\frac{\mathrm{log}\gamma ^{(d)}}{2\pi }.$$
(9)
Let $`D^{}=_{|d|D}1`$. Then Rubinstein proved that
$$\frac{1}{D^{}}\underset{|d|D}{}\underset{\genfrac{}{}{0pt}{}{\gamma _{j_1},\mathrm{}\gamma _{j_n}}{j_mj_n}}{}f(\stackrel{~}{\gamma }_{j_1}^{(d)},\mathrm{},\stackrel{~}{\gamma }_{j_n}^{(d)})_{R^n}W_{Sp,n}(\stackrel{}{x})f(\stackrel{}{x})𝑑\stackrel{}{x},$$
(10)
as $`D\mathrm{}`$ where $`W_{Sp,n}`$ is as in section 5.5, provided that $`f`$ is a Schwarz function such that the support of $`\widehat{f}(u)`$ is contained in $`_{j=1}^n|u_j|<1`$. $`W_{Sp,n}`$ is called the $`n`$-level density function for the symplectic group.
Rubinstein also found evidence that another family has a symmetry type O. To describe this family let $`\mathrm{\Delta }(z)=_{n=1}^{\mathrm{}}\tau (n)\text{e}(nz)`$ where $`\tau `$ is Ramanujan’s tau-function. It is well-known that $`\mathrm{\Delta }`$ is a primitive Hecke newform of weight 12 and level 1. The family (c) is obtained from quadratic twists of the L-function associated with $`\mathrm{\Delta }`$ namely
$$L(\mathrm{\Delta },s,\chi _d)=\underset{n=1}{\overset{\mathrm{}}{}}\frac{\tau (n)n^{11/2}\chi _d(n)}{n^s}.$$
(11)
If $`d<0`$ then this L-function automatically vanishes at $`s=1/2`$ because the associated functional equation occurs with $`ϵ=1`$. Rubinstein computed the lowest zero of $`L(\mathrm{\Delta },s,\chi _d)`$ for $`d>0`$ and the lowest zero above the real axis for $`L(\mathrm{\Delta },s,\chi _d)`$ with $`d<0`$. At this point, we should mention that an orthogonal matrix has determinant +1 or $`1`$. So there are actually two symmetry types O<sup>+</sup> and O<sup>-</sup> (and the statistics of O are an average of the statistics of these). Katz and Sarnak had computed the statistics (neighbor spacing, correlations, density, eigenvalue nearest 1) for these two symmetry types as well. Rubinstein found that the lowest lying zero of L-functions in (c) with $`d>0`$ followed O<sup>+</sup> while the lowest lying zero above the real axis for L-functions with $`d<0`$ followed O<sup>-</sup>.
Rubinstein considered more generally the twisting of an arbitrary arithmetic L-function $`L(f,s)`$ associated with an automorphic form $`f`$ on $`GL_m`$ by quadratic characters $`\chi _d`$. The results here divide into three cases which have to do with signs of the functional equations. If the L-function $`L(f,s,\text{sym}^2)`$ associated with the symmetric square of $`f`$ is entire, the functional equations for $`L(f,s,\chi _d)`$ will always have $`ϵ=+1`$ and the average over $`d`$ is exactly as above with the $`n`$-level density function being $`W_{Sp,n}`$. In the case that the symmetric square $`L(f,s,\text{sym}^2)`$ has a pole (at $`s=1`$) then the sign of the functional equation for $`L(f,s,\chi _d)`$ is +1 for even characters $`\chi _d`$ and $`1`$ for odd characters $`\chi _d`$. Rubinstein averages over these two cases separately and discovers that the first case yields an $`n`$-level density function $`W_{O^+,n}`$ and in the second case an $`n`$-level density function $`W_{O^{},n}`$ (see section 5.5). In each of these cases a restriction is placed on the support of the Fourier transform of $`f`$.
### 2.4 The Diagonal Terms
Rubinstein’s work gave impressive confirmation of the theory, but the severe restriction on the support of the Fourier transform is worth investigating. In fact, this restriction occurs right at the place where something interesting is happening with the Fourier transfom of the density function. Returning for a moment to the case of the Riemann zeta-function, we can illustrate this idea.
Montgomery’s original theorem involved
$$F(\alpha ,T)=\frac{1}{N(T)}\underset{0<\gamma ,\gamma ^{}T}{}T^{i\alpha (\gamma \gamma ^{})}w(\gamma \gamma ^{}).$$
(12)
Here $`w`$ is a weight function that concentrates at the origin. (Montgomery used $`w(u)=4/(4+u^2)`$.) Assuming the Riemann Hypothesis he showed that
$$F(\alpha ,T)=T^{2\alpha }(1+o(1))+|\alpha |+o(1)$$
(13)
uniformly for $`1+ϵ<\alpha <1ϵ`$ for any $`ϵ>0`$.
By the definition of $`F`$,
$$\frac{1}{N(T)}\underset{0<\gamma ,\gamma ^{}T}{}r(\gamma \gamma ^{})w(\gamma \gamma ^{})=_{\mathrm{}}^{\mathrm{}}\widehat{r}(\alpha )F(\alpha ,T)𝑑\alpha .$$
(14)
Thus, Montgomery’s theorem gives information about the average behavior of differences between the zeros for test functions $`r`$ with the support of $`\widehat{r}`$ contained in $`(1,1)`$. Montgomery went on to conjecture (based on considerations of the behavior of prime pairs) that for $`|\alpha |>1`$ one has $`F(\alpha ,T)=1+o(1)`$; this assertion implies (1). Thus, $`F`$ is not differentiable at $`\alpha =1`$. In the proof of Montgomery’s theorem (via the explicit formula and the mean value theorem (23) for Dirichlet polynomials) for $`|\alpha |<1`$, the main term of $`F`$ arises from the “diagonal” contributions of the mean square of a Dirichlet polynomial (i.e. the terms $`m=n`$ in the integral $`_0^Ta_m\overline{a_n}(m/n)^{it}𝑑t`$). For $`\alpha >1`$, the off-diagonal terms (i.e. $`mn`$) contribute to the main-term.
A similar situation arises in the work of Rudnick and Sarnak and in the work of Rubinstein. All of the proofs of these theorems are valid only in the range where the diagonal terms dominate.
### 2.5 Beyond the Diagonal
So far we have seen that the theory of families is confirmed by numerical data as well as theoretical data up to the diagonal. Bogolmony and Keating gave a heuristic derivation of all of the GUE conjecture (i.e. all the $`n`$-level correlations) based on Hardy-Littlewood type conjectures for pairs of primes and pairs of almost primes; this work shows how the off-diagonal terms potentially contribute.
Özlük \[Oz\] proved an analogue for all primitive Dirichlet characters for the pair correlation theorem (12); he obtained a result for $`|\alpha |<2`$. This was the first example of going beyond the diagonal. See also \[OS\] and \[IS2\].
It is of great interest that Iwaniec, Luo, and Sarnak \[ILS\] have succeeded in going beyond the diagonal in several examples which represent three of the symmetry types. They work with the 1-level density functions assuming only that the Riemann Hypothesis holds for all of the L-functions in question. The 1-level density functions may be obtained from $`W_n`$ by taking $`n=1`$. They are
$$W(\text{O})(x)=1+\frac{1}{2}\delta _0(x),$$
(15)
$$W(\text{O}^+)(x)=1+\frac{\mathrm{sin}2\pi x}{2\pi x},$$
(16)
$$W(\text{O}^{})(x)=1\frac{\mathrm{sin}2\pi x}{2\pi x}+\delta _0(x),$$
(17)
$$W(\text{Sp})(x)=1\frac{\mathrm{sin}2\pi x}{2\pi x},$$
(18)
$$W(\text{U})(x)=1.$$
(19)
The families considered by Iwaniec, Luo, and Sarnak are related to modular forms. Let $`H_k^{}(N)`$ denote the set of holomorphic newforms $`f`$ of weight $`k`$ and level $`N`$. Let $`H_k^+(N)`$ denote the weight $`k`$ level $`N`$ newforms for which the associated L-function has a + in its functional equation, and $`H_k^{}(N)`$ is the subset of $`f`$ for which $`L(f,s)`$ has a $``$ in its functional equation. Let $`M^{}(K,N)`$ be the union of the $`H^{}(k,N)`$ for $`kK`$ and similarly define $`M^+`$ and $`M^{}`$. They consider the low lying zeros of $`L(f,s)`$ as $`f`$ varies through one of the $`M`$-sets. The average spacing for all the zeros of all the $`L(f,s)`$ with $`fH^{}(k,N)`$ up to a fixed height $`t_0`$ is asymptotic to $`2\pi /\mathrm{log}(k^2N)`$. Let $`\varphi `$ be a test function which is even and rapidly decaying. They proved that if the support of $`\widehat{\varphi }`$ is contained in $`(2,2)`$, then
$$\frac{1}{|M^{}(K,N)|}\underset{\genfrac{}{}{0pt}{}{fM^{}(K,N)}{\gamma _f}}{}\varphi \left(\frac{\gamma _f\mathrm{log}k^2N}{2\pi }\right)_{\mathrm{}}^{\mathrm{}}\varphi (x)W(\text{O})(x)𝑑x.$$
(20)
as $`KN\mathrm{}`$. Similar statements hold with $`M^{}`$ replaced by $`M^+`$ and $`M^{}`$ and O replaced by O<sup>+</sup> and O<sup>-</sup>.
It should be pointed out that the Fourier transforms of the density functions $`W(\text{O})(x)`$, $`W(\text{O}^+)(x)`$, and $`W(\text{O}^{})(x)`$ all agree in the diagonal range; so it is only when one goes beyond the diagonal that the distinguishing features of these three symmetry types becomes apparent.
Iwaniec, Luo, and Sarnak also consider the symmetric square L-functions of the $`fM^{}`$ and verify that the above statements hold with symmetry type Sp and the support of $`\widehat{\varphi }`$ in (-3/2,3/2). Also, the average zero spacing is $`2\pi /\mathrm{log}(k^2N^2)`$ so $`N`$ should be replaced by $`N^2`$ in the argument of $`\varphi `$ in the left hand side of (20).
Thus, the theoretical and numerical evidence that the zeros of families of L-functions depend on the symmetry type of the family is pretty convincing.
## 3 Moments
So far we have seen that the eigenvalues of matrices from unitary groups are excellent models for zeros of families of L-functions. Now we want to take the matrix models a significant step further and argue that the characteristic polynomials of these matrices on average reveal very important features of the value distribution of the L-functions in the family.
### 3.1 Moments of the Riemann zeta-function
We first look at the situation of the Riemann zeta-function and its moments.
To give some background, we cite the theorem of Hardy and Littlewood:
$$\frac{1}{T}_0^T|\zeta (\frac{1}{2}+it)|^2𝑑t\mathrm{log}T$$
(21)
and the theorem of Ingham:
$$\frac{1}{T}_0^T|\zeta (\frac{1}{2}+it)|^4𝑑t\frac{1}{2\pi ^2}\mathrm{log}^4T.$$
(22)
The asymptotics of no other moments (apart from the trivial 0-th moment) are known. In general it has been conjectured that
$$\frac{1}{T}_0^T|\zeta ({\scriptscriptstyle \frac{1}{2}}+it)|^{2k}𝑑tc_k\mathrm{log}^{k^2}T.$$
(23)
The basic tools for investigating mean-values in $`t`$-aspect are the mean value theorem for Dirichlet polynomials (due to Montgomery and Vaughan):
$$_0^T\left|\underset{nN}{}a_nn^{it}\right|^2𝑑t=\underset{nN}{}(T+O(n))|a_n|^2$$
(24)
and some sort of formula expressing the function in question in terms of Dirichlet polynomials (such as an approximate functional equation) such as
$$\zeta (s)^k=\underset{n\tau ^k}{}\frac{d_k(n)}{n^s}+\chi (s)^k\underset{n\tau ^k}{}\frac{d_k(n)}{n^{1s}}+E(s),$$
(25)
where $`E(s)`$ should be small on average, $`s=1/2+it`$, $`\tau =\sqrt{t/(2\pi )}`$, $`\chi (s)`$ is the factor from the functional equation (3), and where
$$\zeta (s)^k=\underset{n=1}{\overset{\mathrm{}}{}}\frac{d_k(n)}{n^s}(\sigma >1)$$
(26)
so that $`\zeta (s)^k`$ is the generating function for $`d_k(n)`$. Note that the mean-value theorem for Dirichlet polynomials detects only diagonal contributions.
Conrey and Ghosh \[CG2\] gave the moment conjecture a more precise form, namely that there should be a factorization
$$c_k=\frac{g_ka_k}{\mathrm{\Gamma }(1+k^2)}$$
(27)
where
$$a_k=\underset{p}{}\left(1\frac{1}{p}\right)^{k^2}\underset{j=0}{\overset{\mathrm{}}{}}\frac{d_k(p^j)^2}{p^j}$$
(28)
is an arithmetic factor and $`g_k`$, a geometric factor, should be an integer. Note that by the mean-value theorem for Dirichlet polynomials it is not difficult to show that
$$\frac{1}{T}_0^T\left|\underset{nx}{}\frac{d_k(n)}{N^{1/2+it}}\right|^2𝑑t\frac{a_k(\mathrm{log}x)^{k^2}}{\mathrm{\Gamma }(1+k^2)},$$
(29)
provided that $`x=o(T)`$. Thus, an interpretation of $`g_k`$ is
$$g_k=\underset{T\mathrm{}}{lim}\frac{_0^T|\zeta ^k({\scriptscriptstyle \frac{1}{2}}+it)|^2𝑑t}{_0^T\left|_{nT}\frac{d_k(n)}{n^{1/2+it}}\right|^2𝑑t}$$
(30)
assuming that the limit exists, so that $`g_k`$ represents the ‘number’ of Dirichlet polynomial approximations to $`\zeta (s)^k`$ of length $`T`$ required to measure the mean square of $`\zeta (s)^k`$.
In this notation, the result of Hardy and Littlewood is that $`g_1=1`$ and Ingham’s result is that $`g_2=2`$.
These results can be obtained essentially from the mean-value theorem for Dirichlet polynomials. To go beyond the fourth moment requires taking into account off-diagonal contributions. Goldston and Gonek \[GG\] describe a precise way to transform information about coefficient correlations $`_{nx}a(n)a(n+r)`$ into a formula for the mean square of a long Dirichlet polynomial $`_{nx}a(n)n^s`$ where $`x`$ is bigger than the length of integration.
Using Dirichlet polynomial techniques Conrey and Ghosh \[CG1\] conjectured that $`g_3=42`$ and Conrey and Gonek \[CGo\] conjectured that $`g_4=24024`$. Meanwhile, Keating and Snaith \[KeSn1\] computed the moments of characteristic polynomials of matrices in $`U(N)`$ and found that for any real $`x`$ and any complex number $`s`$,
$$M_{U,N}(s)=_{U(N)}|det(AI\mathrm{exp}(ix))|^{2s}𝑑A=\underset{j=1}{\overset{N}{}}\frac{\mathrm{\Gamma }(j)\mathrm{\Gamma }(j+2s)}{\mathrm{\Gamma }(j+s)^2},$$
(31)
where $`dA`$ denotes the Haar measure for the group $`U(N)`$ of $`N\times N`$ unitary matrices. To do this calculation, they made use of Weyl’s formula for the Haar measure (see section 5.3) and Selberg’s integral (see section 5.6). They also showed that
$$\underset{N\mathrm{}}{lim}\frac{M_N(s)}{N^{s^2}}=\frac{G(1+s)^2}{G(1+2s)},$$
(32)
where $`G(s)`$ is Barnes’ double Gamma-function which satisfies $`G(1)=1`$ and $`G(z+1)=\mathrm{\Gamma }(z)G(z)`$. Note that for $`s=k`$ an integer,
$$\frac{G(1+k)^2}{G(1+2k)}=\underset{j=0}{\overset{k1}{}}\frac{j!}{(j+n)!}.$$
(33)
For $`k=1,2,3`$ the above is 1/1!,2/4!, 42/9! in agreement with the theorems of Hardy and Littlewood, and Ingham and the conjecture of Conrey and Ghosh. Keating and Snaith argued that one should thus model the moments of the zeta-function from 0 to $`T`$ by moments of characteristic polynomials of unitary matrices of size $`N\mathrm{log}T`$. (More precisely, one should take $`N`$ to be the integer nearest to $`\mathrm{log}\frac{T}{2\pi }`$). They then conjectured that
$$g_k=k^2!\underset{j=0}{\overset{k1}{}}\frac{j!}{(j+k)!}$$
(34)
for integer $`k`$. The initial public announcements of the conjectures of Conrey and Gonek (that $`g_4=24024`$) and of Keating and Snaith ( $`g_k`$ for all real $`k1/2`$) occurred at the Vienna conference on the Riemann Hypothesis only moments after it was checked that the Keating and Snaith conjecture does indeed predict that $`g_4=24024`$.
### 3.2 Moments of L-functions at 1/2
Subsequently, Conrey and Farmer \[CF\] analyzed known results for moments of L-functions at 1/2 and made a general conjecture. (These moments had been considered by a number of authors; see especially \[GV\].) The conjecture has the shape
$$\frac{1}{X^{}}\underset{\genfrac{}{}{0pt}{}{f}{c(f)X}}{}L(f,1/2)^k\frac{g_ka_k}{\mathrm{\Gamma }(1+q(k))}(\mathrm{log}X)^{q(k)}$$
(35)
for some $`a_k`$, $`g_k`$, and $`q(k)`$ where $``$ is a family of $`f`$ parametrized by the conductor $`c(f)`$, and $`X^{}=_{c(f)X}1`$. The observations of Conrey and Farmer were that $`g_k`$ and $`q(k)`$ depend only on the symmetry type of the family, and that $`a_k`$ depends on the family itself, but is explicitly computable in any specific case. Thus, the conjecture is that $`q(k)=k^2`$ for a unitary family, $`q(k)=k(k+1)/2`$ for a symplectic family, and $`q(k)=k(k1)/2`$ for an orthogonal family. The values of $`g_k`$ were left unspecified, but were then predicted as before from random matrix theory by Keating and Snaith \[KeSn2\] and independently by Brezin and Hikami \[BH\] by computing moments of characteristic polynomials of matrices from O$`(N)`$ and from USp(2$`N)`$. Each is a quotient of products of Gamma-functions.
Thus, for the family of Dirichlet L-functions $`L(s,\chi _d)`$ with a real primitive Dirichlet character $`\chi _d`$ modulo $`d`$ we have the following results ($`D^{}=_{|d|D}1`$): Jutila \[J\] proved that
$$\frac{1}{D^{}}\underset{|d|D}{}L({\scriptscriptstyle \frac{1}{2}},\chi _d)a_1\mathrm{log}(D^{\frac{1}{2}})$$
(36)
and
$$\frac{1}{D^{}}\underset{|d|D}{}L^2({\scriptscriptstyle \frac{1}{2}},\chi _d)2\frac{a_2\mathrm{log}^3(D^{\frac{1}{2}})}{3!}.$$
(37)
Soundararajan \[So1\] showed that
$$\frac{1}{D^{}}\underset{|d|D}{}L^3({\scriptscriptstyle \frac{1}{2}},\chi _d)16\frac{a_3\mathrm{log}^6(D^{\frac{1}{2}})}{6!}$$
(38)
and conjectured that
$$\frac{1}{D^{}}\underset{|d|D}{}L^4({\scriptscriptstyle \frac{1}{2}},\chi _d)768\frac{a_4\mathrm{log}^{10}(D^{\frac{1}{2}})}{10!}.$$
(39)
The general conjecture coming from random matrix theory (see section 5.7), which agrees with the above, is:
$$\frac{1}{D^{}}\underset{|d|D}{}L^k({\scriptscriptstyle \frac{1}{2}},\chi _d)\underset{\mathrm{}=1}{\overset{k}{}}\frac{\mathrm{}!}{2\mathrm{}!}a_k\mathrm{log}^{k(k+1)/2}(D)$$
(40)
where
$$a_k=\underset{p}{}\frac{\left(1{\scriptscriptstyle \frac{1}{p}}\right)^{\frac{k(k+1)}{2}}}{\left(1+{\scriptscriptstyle \frac{1}{p}}\right)}\left(\frac{\left(1{\scriptscriptstyle \frac{1}{\sqrt{p}}}\right)^k+\left(1+{\scriptscriptstyle \frac{1}{\sqrt{p}}}\right)^k}{2}+\frac{1}{p}\right).$$
(41)
An example of an orthogonal family where several moments are known arises from $`_q`$ the set of primitive cusp forms of weight $`2`$ and level $`q`$ ($`q`$ prime). Then, from results of Duke \[D\], Duke, Friedlander, and Iwaniec \[DFI\], Iwaniec and Sarnak \[IS1\], and Kowalski, Michel, and VanderKam \[KMV1\] and \[KMV2\], we have
$$\frac{1}{|_q|}\underset{f_q}{}L(1/2,f)a_1$$
(42)
$$\frac{1}{|_q|}\underset{f_q}{}L^2(1/2,f)2a_2\mathrm{log}q^{\frac{1}{2}}$$
(43)
$$\frac{1}{|_q|}\underset{f_q}{}L^3(1/2,f)8a_3\frac{\mathrm{log}^3q^{\frac{1}{2}}}{3!}$$
(44)
$$\frac{1}{|_q|}\underset{f_q}{}L^4(1/2,f)128a_4\frac{\mathrm{log}^6q^{\frac{1}{2}}}{6!}$$
(45)
where $`a_1=\zeta (2),`$
$$a_2=\zeta (2)^2\underset{p}{}\left(1+{\scriptscriptstyle \frac{1}{p^2}}\right)$$
(46)
$$a_3=\zeta (2)^3\underset{p}{}\left(1{\scriptscriptstyle \frac{1}{p}}\right)\left(1+{\scriptscriptstyle \frac{1}{p}}+{\scriptscriptstyle \frac{4}{p^2}}+{\scriptscriptstyle \frac{1}{p^3}}+{\scriptscriptstyle \frac{1}{p^4}}\right)$$
(47)
$$a_4=\zeta (2)^5\underset{p}{}\left(1{\scriptscriptstyle \frac{1}{p}}\right)^3\left(1+{\scriptscriptstyle \frac{3}{p}}+{\scriptscriptstyle \frac{11}{p^2}}+{\scriptscriptstyle \frac{10}{p^3}}+{\scriptscriptstyle \frac{11}{p^4}}+{\scriptscriptstyle \frac{3}{p^5}}+{\scriptscriptstyle \frac{1}{p^6}}\right)$$
(48)
We have not found a simple expression for $`a_k`$, though it can be determined explicitly for each $`k`$. As before, a general conjecture is:
$$\frac{1}{|_q|}\underset{f_q}{}L^k(1/2,f)2^{k1}\underset{\mathrm{}=1}{\overset{k1}{}}\frac{\mathrm{}!}{2\mathrm{}!}a_k\mathrm{log}^{k(k1)/2}q.$$
(49)
We believe that in formulas for moments of L-functions over a family the power of the log of the conductor and the value of $`g_k`$ should only depend on the symmetry type of the family and that the value of $`a_k`$ will depend on the family but can always be determined explicitly.
## 4 Further directions
In this section we mention some questions where further research is desirable.
### 4.1 Full moment conjecture
What are the lower order terms in the moment formulae for $`|\zeta (1/2+it)|^{2k}`$ and for $`L(1/2)^k`$? These are known in a few instances (see \[In\], \[C\] for the second and fourth moments of $`\zeta (s)`$) but not in general. The difficulty is that random matrix theory does not “see” the contribution of the arithmetic factor $`a_k`$. Lower order terms will likely involve a mix of derivatives of $`a_k`$ and secondary terms from the moments of the characteristic polynomials of matrices. In general, a better understanding of how $`\zeta (s)`$ is modeled by a characteristic polynomial of a certain type of matrix is needed; how do the primes come into play? Perhaps we should think of $`\zeta (1/2+it)`$ as a partial Hadamard product over zeros multiplied by a partial Euler product. Perhaps these two parts behave independently, and the Hadamard product part can be modeled by random matrix theory.
### 4.2 Distribution of Values
Keating and Snaith compute explicit formulas for the $`s`$-th moment of the characteristic polynomials of matrices from O($`N`$), Sp($`N`$), and U($`N`$). Consequently the value distributions for these characteristic polynomials can be explicitly computed; they involve the Fourier transform of the $`s`$-th moment. Preliminary investigations indicate that there is a good fit between the random matrix formulae and numerical data. One particularly interesting feature of this investigation involves the understanding of zeros which occur exactly at 1/2. These seem to occur only for L-functions in an orthogonal family. If the sign of the functional equation is $`1`$, there is automatically a zero at 1/2. The interesting situation is when the sign is + and there is still a zero; for example if E is an elliptic curve defined over $`Q`$ then examination of the distribution of values of O<sup>+</sup> suggest that twists $`L(E,s,\chi _p)`$ of the L-function by quadratic characters seem to vanish for about $`X^{3/4}(\mathrm{log}X)^{5/8}`$ values of $`|p|<X`$ with sign +1 in the functional equation.
### 4.3 Extreme Values
How large is the maximum value of $`|\zeta (1/2+it)|`$ for $`T<t<2T`$? It is known that the Riemann Hypothesis implies that the maximum is at most $`\mathrm{exp}(c\mathrm{log}T/\mathrm{log}\mathrm{log}T)`$ for some $`c>0`$. It is also known that the maximum gets as big as $`\mathrm{exp}(c_1(\mathrm{log}T/\mathrm{log}\mathrm{log}T)^{1/2})`$ for a sequence of $`T\mathrm{}`$ for some $`c_1>0`$. It has been conjectured that the smaller bound (the one that is known to occur) is closer to the truth. However, the new conjectures about moments suggest that it may be the larger.
This question has a number of equivalent and analogous (for L-values) formulations. How big can $`S(T):=\frac{1}{\pi }\mathrm{arg}\zeta (1/2+iT)`$ be? Assuming the Riemann Hypothesis, it is known that $`S(T)\mathrm{log}T/\mathrm{log}\mathrm{log}T`$ but that infinitely often it is bigger than $`c(\mathrm{log}T/\mathrm{log}\mathrm{log}T)^{1/2}`$ for some $`c>0`$ Which is closer to the truth? What is the maximum size of the class number of an imaginary quadratic field (as a function of the discriminant) (see \[Sh\]for a discussion and numerical investigation.)? How big can the least quadratic non-residue of a given prime $`p`$ be ($`\mathrm{log}p`$ or $`(\mathrm{log}p)^{2ϵ}`$? See \[M2\] for a discussion of this question? What is the maximal order of vanishing of an L-function at 1/2? In terms of the conductor $`N`$, can it be as big as $`\mathrm{log}N/\mathrm{log}\mathrm{log}N`$ or is it at most the square root of that, or something entirely different? All of these questions are related, at least by analogy, and they may all have similar answers. It would be interesting and surprising if in each case it is the larger bound which is closer to the truth.
### 4.4 Zeros of $`\zeta ^{}(s)`$
Can one use random matrix theory to predict the horizontal distribution of the real parts of the zeros of $`\zeta ^{}`$? It is known that the Riemann Hypothesis is equivalent to the assertion that each non-real zero of $`\zeta ^{}(s)`$ has real part greater than or equal to 1/2. Moreover, if such a zero has real part 1/2, then it is also a zero of $`\zeta (s)`$ (and so a multiple zero of $`\zeta (s)`$). These assertions are the point of departure for Levinson’s work on zeros of the Riemann zeta-function on the critical line. It would be interesting to know the horizontal distribution of these zeros; in particular what proportion of them with ordinates between $`T`$ and $`2T`$ are within $`a/\mathrm{log}T`$ of the 1/2-line?
In a similar vein, the Riemann $`\xi `$-function is real on the 1/2-line and has all of its zeros there (assuming the Riemann Hypothesis). It is an entire function of order 1; because of its functional equation, $`\xi (1/2+i\sqrt{z})`$ is an entire function of order 1/2. It follows that the Riemann Hypothesis implies that all zeros of $`\xi ^{}(s)`$ are on the 1/2-line. Assuming this to be true, one can ask about the vertical distribution of zeros of $`\xi ^{}(s)`$, and more generally of $`\xi ^{(m)}(s)`$. It seems that the zeros of higher derivatives will become more and more regularly spaced; can these distributions be expressed in a simple way using random matrix theory?
### 4.5 Long Mollifiers, Local Integrals, and GUE
David Farmer \[F1\], \[F2\] has made two very interesting conjectures having to do with $`\zeta (s)`$. The first is a conjecture about the mean square of $`\zeta (1/2+it)`$ times an arbitrarily long mollifier. A mollifier is a Dirichlet polynomial with coefficients equal to the Möbius $`\mu `$-function times a smooth function. The length of the mollifier is the length of the Dirichlet polynomial. He has also conjectured that
$$\frac{1}{T}_0^T\frac{\zeta (s+u)\zeta (1s+v)}{\zeta (s+a)\zeta (1s+b)}𝑑t1+\frac{(ua)(vb)}{(u+v)(a+b)}(1T^{(u+v)}).$$
(50)
where $`a,b,u,v`$ are complex numbers with positive real part, and $`s=1/2+it`$. These two conjectures are essentially equivalent and imply certain parts of the GUE conjecture. It would be interesting to generalize these and relate them to the full GUE conjecture.
## 5 Appendices
### 5.1 The Classical Groups
* The unitary group U($`N`$) is the group of $`N\times N`$ matrices $`U`$ with entries in $`C`$ for which $`UU^{}=I`$ where $`U^{}`$ denotes the conjugate transpose of $`U`$, i.e. if $`U=(u_{i,j})`$, then $`U^{}=(\overline{u_{j,i}})`$.
* The orthogonal group O($`N`$) is the subgroup of U($`N`$) consisting of matrices with real entries.
* The special orthogonal group SO($`N`$). This is the subgroup of O($`N`$) consisting of matrices with determinant 1. SO(2$`N`$) leads to the symmetry type we have called O<sup>+</sup> and SO(2$`N`$+1) leads to the symmetry type we call O<sup>-</sup>.
* The symplectic group USp(2$`N`$) is the subgroup of U(2$`N`$) of matrices $`U`$ for which $`UZU^t=Z`$ where $`U^t`$ denotes the transpose of $`U`$ and
$$Z=\left(\begin{array}{cc}0& I_N\\ I_N& 0\end{array}\right)$$
(51)
### 5.2 The Weyl Integration Formula
The $`N\times N`$ unitary matrices can be parametrized by their $`N`$ eigenvalues on the unit circle. Any configuration of $`N`$ points on the unit circle corresponds to a conjugacy class of U($`N`$). If $`f(A)=f(\theta _1,\mathrm{}\theta _N)`$ is a symmetric function of $`N`$ variables, then Weyl’s formula \[W\] gives
$$_{U(N)}f(A)𝑑A=\frac{1}{N!}_{[0,1]^N}f(\theta )\underset{1j<kN}{}|e(\theta _j)e(\theta _k)|^2d\theta _1\mathrm{}d\theta _N$$
(52)
where $`dA`$ is the Haar measure. Similarly, on $`\text{Sp}(2N)`$ and SO($`2N`$) we have respectively
$$dA=\frac{2^{N^2}}{N!}\underset{j<k}{}(\mathrm{cos}\pi \theta _j\mathrm{cos}\pi \theta _k)^2\underset{j=1}{\overset{N}{}}\mathrm{sin}^2\pi \theta _j\underset{j=1}{\overset{N}{}}d\theta _j;$$
(53)
$$dA=\frac{2^{(N1)^2}}{N!}\underset{i<j}{}(\mathrm{cos}\pi \theta _j\mathrm{cos}\pi \theta _k)^2\underset{j=1}{\overset{N}{}}d\theta _j.$$
(54)
### 5.3 Four Statistics
Suppose we have a sequence $`𝒯`$ of $`N`$-tuples of numbers $`T_N=\{t_1,t_2,\mathrm{},t_N\}`$ where $`t_1t_2<\mathrm{}<t_N`$ such that for each set the average spacing $`t_{j+1}t_j`$ is asymptotically 1. We write $`t_{i,N}`$ in place of $`t_i`$ if we need to indicate that $`t_iT_N`$.
* The $`n`$-level density of $`𝒯`$ is $`W(\stackrel{}{x})=W(x_1,\mathrm{}x_n)`$ means that
$$\underset{N\mathrm{}}{lim}\underset{\genfrac{}{}{0pt}{}{(i_1,\mathrm{}i_n),i_jN}{i_ji_k}}{}f(t_{i_1},\mathrm{}t_{i_n})=_{R^n}f(\stackrel{}{x})W(\stackrel{}{x})𝑑\stackrel{}{x}.$$
(55)
for a Schwarz-class $`f`$. The sum is over $`n`$-tuples with distinct entries.
* The $`j`$-th lowest zero density is $`\nu _j(x)`$ means that for a test function $`f`$
$$\underset{N\mathrm{}}{lim}\frac{1}{N}\underset{nN}{}f(t_{j,n})=_0^{\mathrm{}}f(x)\nu _j(x)𝑑x.$$
(56)
* The consecutive spacing density is $`\mu (x)`$ means that for a test function $`f(x)`$ we have
$$\underset{N\mathrm{}}{lim}\frac{1}{N}\underset{iN1}{}f(t_{i+1}t_i)=_0^{\mathrm{}}f(x)\mu (x)𝑑x.$$
(57)
* The $`n`$-correlation density is $`V(x_1,\mathrm{},x_n)`$ means that for test functions $`f`$ that are symmetric in all of the variables, depend only on the differences of the variables (i.e. $`f(x_1+u,\mathrm{},x_n+u)=f(x_1,\mathrm{},x_n)`$ for all $`u`$), and are rapidly decaying on the hyperplane $`P_n:\{(x_1,\mathrm{},x_n:x_i=0\}`$, we have, as $`N\mathrm{}`$,
$$\frac{1}{N}\underset{\genfrac{}{}{0pt}{}{t_1,\mathrm{}t_nT_N}{i_ji_k}}{}f(t_1,\mathrm{}t_n)_{P_n}f(\stackrel{}{x})V(\stackrel{}{x})𝑑x_1\mathrm{}𝑑x_{n1}$$
(58)
as $`N\mathrm{}`$. The spacing and $`n`$-correlation densities are universal, i.e. the same for each of O, Sp, and U, whereas the $`n`$-level and $`j`$-th lowest zero densities depend on the symmetry type.
### 5.4 Gaudin’s Lemma
Associated to each $`N\times N`$ unitary matrix $`A`$ are its $`N`$ eigenvalues $`\text{e}(\theta _j)`$ where $`0\theta _1\mathrm{}\theta _N1`$. We integrate a function $`F(A)`$ over U($`N`$) by parametrizing the group by the $`\theta _i`$ and using Weyl’s formula to convert the integral into an $`N`$-fold integral over the $`\theta _i`$.
Often one wants to integrate with respect to Haar measure over U($`N`$) a function $`F(A)=\stackrel{~}{f}(A)=\stackrel{~}{f}(\theta _1,\mathrm{},\theta _N)`$ of $`N`$ variables that is “lifted ” from a function $`f`$ of $`n`$ variables:
$$\stackrel{~}{f}(\theta _1,\mathrm{},\theta _N)=\underset{\genfrac{}{}{0pt}{}{(i_1,\mathrm{},i_n)}{i_ji_k}}{}f(\theta _{i_1},\mathrm{},\theta _{i_n})$$
(59)
where the sum is over all possible $`n`$-tuples $`(i_1,\mathrm{},i_n)`$ of distinct integers between 1 and $`N`$. Gaudin’s lemma gives a simplification of this computation from an $`N`$-fold integral to an $`n`$-fold integral. The Haar measure (see section 5.2) at the matrix $`A`$ can be expressed as
$$dA=\frac{1}{N!}\underset{1j<kN}{}|e(\theta _j)e(\theta _k)|^2d\theta _1\mathrm{}d\theta _N.$$
(60)
The product here is the square of the absolute value of the $`N\times N`$ Vandermonde determinant with $`j,k`$ entry $`\text{e}(\theta _k)^{j1}=\text{e}((j1)\theta _k)`$ . It is also the $`N\times N`$ determinant of the matrix with $`j,k`$ entry $`J_N(\theta _j\theta _k)`$ where
$$J_N(\theta )=\underset{m=0}{\overset{N1}{}}\text{e}(m\theta )=\text{e}((N1)\theta /2)\frac{\mathrm{sin}\pi N\theta }{\mathrm{sin}\pi \theta }.$$
(61)
Thus,
$$_{U(N)}\stackrel{~}{f}(A)𝑑A=_{[0,1]^N}\stackrel{~}{f}(\theta _1,\mathrm{},\theta _N)\frac{1}{N!}\underset{N\times N}{det}J_N(\theta _j\theta _k)\underset{j=1}{\overset{N}{}}d\theta _j.$$
(62)
Then, Gaudin’s lemma asserts the equality
$$_{U(N)}\stackrel{~}{f}(A)𝑑A=_{[0,1]^n}f(\theta _1,\mathrm{},\theta _n)\frac{1}{n!}\underset{n\times n}{det}J_N(\theta _j\theta _k)\underset{j=1}{\overset{n}{}}d\theta _j.$$
(63)
This principle works for all of the subgroups of U($`N`$) under consideration here as well. (See \[KS\] section 5.1 for a general statement and proof of this important lemma.) We illustrate by computing the $`n`$-level density function for U($`N`$). Note that
$$\underset{N\mathrm{}}{lim}\frac{1}{N}J_N(\theta /N)=\text{e}(\theta /2)\frac{\mathrm{sin}\pi \theta }{\pi \theta }$$
(64)
from which it follows easily that
$$\underset{N\mathrm{}}{lim}\frac{1}{N^n}\underset{n\times n}{det}J_N(\theta _j\theta _k)=\underset{n\times n}{det}K_0(\theta _1,\mathrm{},\theta _n)$$
(65)
where $`K_ϵ`$ is defined in section 5.5.
Now let $`f(\stackrel{}{x})=f(x_1,\mathrm{},x_n)`$ be a test function. To compute the $`n`$-level density (compare with section 5.4) we need to evaluate
$$\underset{N\mathrm{}}{lim}_{U(N)}\underset{\genfrac{}{}{0pt}{}{(i_1,\mathrm{},i_n)}{i_ji_k}}{}f(\widehat{\theta }_{i_1},\mathrm{},\widehat{\theta }_{i_n})\underset{j<k}{}|\text{e}(\theta _j)\text{e}(\theta _j)|^2d\theta _1\mathrm{}d\theta _N$$
(66)
By Gaudin’s lemma and after using the new expression for the Haar measure and changing variables $`\theta _jx_j/N`$, the above is equal to
$`\underset{N\mathrm{}}{lim}`$ $`{\displaystyle \frac{1}{N^n}}{\displaystyle _{[0,N]^n}}f(\theta _1,\mathrm{},\theta _n){\displaystyle \frac{1}{n!}}\underset{n\times n}{det}J_N(\theta _j\theta _k){\displaystyle \underset{j=1}{\overset{n}{}}}d\theta _j`$ (68)
$`={\displaystyle _{R^n}}f(x_1,\mathrm{},x_n)\underset{n\times n}{det}K_0(x_1,\mathrm{},x_n)dx`$
so that $`W_{U,n}(x_1,\mathrm{},x_n)=det_{n\times n}K_0(x_1,\mathrm{},x_n)`$.
### 5.5 Formulas for the Density Functions
* The $`n`$–level density is $`W(x_1,\mathrm{},x_n)=det_{n\times n}K_ϵ(x_1,\mathrm{},x_n)`$ where $`K_ϵ(x_1,\mathrm{},x_n)`$ is the $`n\times n`$ matrix with entries
$$(K_ϵ(x_1,\mathrm{},x_n))_{i,j}=\frac{\mathrm{sin}\pi (x_ix_j)}{\pi (x_ix_j)}+ϵ\frac{\mathrm{sin}\pi (x_i+x_j)}{\pi (x_i+x_j)}$$
(69)
where $`ϵ=0`$ for U; $`ϵ=1`$ for Sp; $`ϵ=1`$ for O<sup>+</sup>. Also,
$$W_{O^{},n}(\stackrel{}{x})=\underset{n\times n}{det}(K_1(\stackrel{}{x}))+\underset{m=1}{\overset{n}{}}\delta (x_m)\underset{n1\times n1}{det}(K_1^{(m)}(\stackrel{}{x}))$$
(70)
where $`\delta `$ is the Dirac $`\delta `$-function and the superscript $`m`$ denotes that the $`m`$-th row and $`m`$-th column have been deleted from $`K_1(\stackrel{}{x})`$.
* The lowest zero density is $`\nu _1(x)`$ where
$$\nu _1(x)=\frac{d}{dx}\underset{j=0}{\overset{\mathrm{}}{}}(1\lambda _j(x))\text{U};$$
(71)
$$\nu _1(x)=\frac{d}{dx}\underset{j=0}{\overset{\mathrm{}}{}}(1\lambda _{2j+1}(2x))\text{Sp};$$
(72)
$$\nu _1(x)=\frac{d}{dx}\underset{j=0}{\overset{\mathrm{}}{}}(1\lambda _{2j}(2x))\text{O},$$
(73)
where $`1\lambda _0(x)\lambda _1(x)\mathrm{}`$ are the eigenvalues of
$$_{x/2}^{x/2}\frac{\mathrm{sin}\pi (tu)}{\pi (tu)}f(u)𝑑u=\lambda (x)f(t)$$
(74)
* The consecutive spacing density is
$$\mu (x)=\underset{j=0}{\overset{\mathrm{}}{}}(1\lambda _j(x)).$$
(75)
* The $`n`$–correlation density, $`V(x_1,\mathrm{},x_n)=W_{U,n}(x_1,\mathrm{},x_n)`$,
### 5.6 The Selberg Integral
There are many versions of Selberg’s integral see \[Me\]; one is as follows.
If $`\mathrm{}\alpha >0,\mathrm{}\beta >0,\mathrm{}\gamma >\mathrm{min}(\frac{1}{n},\frac{\mathrm{}\alpha }{n1},\frac{\mathrm{}\beta }{n1})`$, then
$`{\displaystyle _1^1}\mathrm{}{\displaystyle _1^1}{\displaystyle \underset{1i<jN}{}}|x_ix_j|^{2\gamma }{\displaystyle \underset{j=1}{\overset{n}{}}}(1x_j)^{\alpha 1}(1+x_j)^{\beta 1}dx_j`$ (76)
$`=2^{\gamma n(n1)+n(\alpha +\beta 1)}{\displaystyle \underset{j=0}{\overset{n1}{}}}{\displaystyle \frac{\mathrm{\Gamma }(1+\gamma +j\gamma )\mathrm{\Gamma }(\alpha +j\gamma )\mathrm{\Gamma }(\beta +j\gamma )}{\mathrm{\Gamma }(1+\gamma )\mathrm{\Gamma }(\alpha +\beta +\gamma (n+j1))}}.`$ (77)
### 5.7 Moments of Characteristic Polynomials
$`M_{U,N}(s)`$ $`=`$ $`{\displaystyle _{U(N)}}|det(AI\mathrm{exp}(ix))|^{2s}𝑑A`$ (78)
$`=`$ $`{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{\mathrm{\Gamma }(j)\mathrm{\Gamma }(j+2s)}{\mathrm{\Gamma }(j+s)^2}},`$ (79)
$`M_{Sp,2N}(s)`$ $`=`$ $`{\displaystyle _{Sp(2N)}}|det(AI)|^s𝑑A`$ (80)
$`=`$ $`2^{2Ns}{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{\mathrm{\Gamma }(1+N+j)\mathrm{\Gamma }(1/2+s+j+s)}{\mathrm{\Gamma }(1/2+j)\mathrm{\Gamma }(1+s+N+j)}},`$ (81)
$`M_{O,2N}(s)`$ $`=`$ $`{\displaystyle _{O(2N)}}|det(AI)|^s𝑑A`$ (82)
$`=`$ $`2^{Ns}{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{\mathrm{\Gamma }(N+j1)\mathrm{\Gamma }(s+j1/2)}{\mathrm{\Gamma }(j1/2)\mathrm{\Gamma }(s+j+N1)}}.`$ (83)
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# Charge and Orbital Ordering and Spin State Transition Driven by Structural Distortion in YBaCo2O5
\[
## Abstract
We have investigated electronic structures of antiferromagnetic YBaCo<sub>2</sub>O<sub>5</sub> using the local spin-density approximation (LSDA) + $`U`$ method. The charge and orbital ordered insulating ground state is correctly obtained with the strong on-site Coulomb interaction. Co<sup>2+</sup> and Co<sup>3+</sup> ions are found to be in the high spin (HS) and intermediate spin (IS) state, respectively. It is considered that the tetragonal to orthorhombic structural transition is responsible for the ordering phenomena and the spin states of Co ions. The large contribution of the orbital moment to the total magnetic moment indicates that the spin-orbit coupling is also important in YBaCo<sub>2</sub>O<sub>5</sub>.
\]
Recently, an interesting spin state transition of the Co<sup>2+</sup> ion in YBaCo<sub>2</sub>O<sub>5</sub> has been reported by Vogt et al. using the neutron powder diffraction (NPD) measurements. The transition is induced by the long-range orbital and charge ordering of Co<sup>2+</sup>/Co<sup>3+</sup> ions. The ordered oxygen-deficient double perovskite $`R`$BaCo<sub>2</sub>O<sub>5+δ</sub> ($`R`$ = rare-earths) has attracted much attention as a new spin-charge-orbital coupled system like manganites and also as a new Co-based colossal magnetoresistance (CMR) material. Indeed, giant magnetoresistance are observed for $`R`$=Gd and Eu, $`(\rho _0\rho _{H=7\mathrm{T}})/\rho _0`$ = 41% and 40% for GdBaCo<sub>2</sub>O<sub>5.4</sub> and EuBaCo<sub>2</sub>O<sub>5.4</sub>, respectively.
In the paramagnetic phase, YBaCo<sub>2</sub>O<sub>5</sub> is crystallized in the tetragonal structure of the space group $`P4/mmm`$. It consists of double CoO<sub>5</sub> square base pyramidal backbone layers along the $`c`$-axis in which Y and Ba layers intervene alternatively and oxygens are deficient exclusively from the Y layers. According to the valency consideration, Co<sup>2+</sup> and Co<sup>3+</sup> ions coexist similarly as the Mn<sup>3+</sup>/Mn<sup>4+</sup> covalency in hole-doped La<sub>1-x</sub>Sr<sub>x</sub>MnO<sub>3</sub>. Below $`T_\mathrm{N}330`$ K, YBaCo<sub>2</sub>O<sub>5</sub> undergoes a $`G`$-type antiferromagnetic (AFM) transition and the lattice changes slightly from the tetragonal to orthorhombic structure. At $`T_{\mathrm{CO}}220`$ K, a pronounced upturn is observed in the resistivity indicating that another transition takes place, i.e., the long-range charge and orbital ordering of Co<sup>2+</sup>/Co<sup>3+</sup> ions. The stripe type charge ordering is formed in the $`ab`$ plane. Further, the spin state of the Co<sup>2+</sup> ion changes from the low to high spin across $`T_{\mathrm{CO}}`$, which is evidenced by the increased magnetic moment per Co ion from $`2.10\mu _B`$ at 300 K to $`3.45\mu _B`$ at 25 K. More recently, for the isostructural HoBaCo<sub>2</sub>O<sub>5</sub>, essentially the same features are observed of $`T_\mathrm{N}340`$ K and $`T_{\mathrm{CO}}210`$ K.
The phenomenon of the spin state transition is observed usually in cobaltates such as LaCoO<sub>3</sub> and La<sub>1-x</sub>Sr<sub>x</sub>CoO<sub>3</sub>. In LaCoO<sub>3</sub>, the magnetic ground state of Co<sup>3+</sup> ion corresponds to the low spin (LS) state with $`t_{2g}^6e_g^0`$. Upon heating, a successive spin state transition to an intermediate spin (IS) state and then to a high spin (HS) state occurs. Note that, distinctly from the case in LaCoO<sub>3</sub>, the spin state transition in YBaCo<sub>2</sub>O<sub>5</sub> seems to occur for Co<sup>2+</sup> ion. In fact, this issue is still controversial. Based on the reduction of the magnetic susceptibility below 220K, Akahoshi and Ueda suggested that the AFM transition takes place at $`T220`$ K which is associated with a spin state transition of Co<sup>3+</sup> from the HS to LS state upon cooling. Thus, the nature of the spin state transition and the interplay of the spin state with the charge and orbital ordering are still unclear.
To reveal the mechanism of the spin state transition as well as the charge and orbital ordering, we have explored the electronic structure of the $`G`$-type AFM YBaCo<sub>2</sub>O<sub>5</sub> using the local-spin density approximation (LSDA) + $`U`$ scheme implemented in the linearized muffin-tin orbital band method.
We have employed two structural data of nearly tetragonal structure at 300 K (L1) and orthorhombic structure at 25 K (L2). In the L1 structure, all the Co sites have an equal average bond length of $`d`$(Co-O) = 1.97 Å. Whereas in the L2 structure, there are two different kinds of Co sites, CoI and CoII. The bond length at CoI sites is extended to $`d`$(CoI-O) = 2.03 Å and that at CoII sites becomes shortened to $`d`$(CoII-O) = 1.92 Å. The $`G`$-type AFM spin order is assumed in all our calculations.
In Fig. 1, we have compared the LSDA Co $`3d`$ partial density of states (PDOS) for the L1 and L2 structures. Most notable is the change of the exchange splitting which becomes larger for CoI site and smaller for CoII in L2 than that for Co in L1. This indicates that, due to the lattice distortion, $`3d`$ electrons in L2 becomes more localized for CoI and less localized for CoII in comparison to those for Co in L1. The calculated spin magnetic moments are $`\mu _S=1.68\mu _B`$ and $`0.95\mu _B`$ for each CoI and CoII site, respectively, in L2 and $`\mu _S=1.13\mu _B`$ for Co in L1. The sizes of the spin moments are consistent with the degree of $`3d`$ electron localization. Hence, the LSDA gives a qualitatively good information about the structural transition effects on Co $`3d`$ electron states. The LSDA, however, yields an incorrect metallic ground state which pertains even after the structural transition from the L1 to L2 structure. This is different from the experiment which shows unambiguously the semiconducting resistivity behavior below 220 K.
Using the LSDA, one cannot expect proper description of localized Co $`3d`$ electrons in YBaCo<sub>2</sub>O<sub>5</sub>. To resolve the above problem, we have applied the LSDA + $`U`$ method with parameter values of $`U=5.0`$ eV and $`J=0.89`$ eV. Although there is an arbitrariness of the $`U`$-value in our calculation, the LSDA + $`U`$ results are usually not much sensitive to the used $`U`$-value within $`\mathrm{\Delta }U\pm 1`$ eV. The spin-orbit interaction is also taken into account in the self-consistent variational loop, because Co $`3d`$ electrons are expected to retain atomic properties to some extent due to their localized nature.
Figures 2 and 3 present CoI and CoII $`3d`$ PDOS, respectively, in the L2 structure obtained by the LSDA + $`U`$ calculations. In the bottom panels, $`t_{2g}`$ and $`e_g`$ decompositions of $`3d`$ PDOS are also provided. It is amusing to note that the energy gap of $`E_g0.6`$ eV opens at the Fermi level $`E_\mathrm{F}`$ and so YBaCo<sub>2</sub>O<sub>5</sub> becomes an insulator as expected in consideration of the Co<sup>2+</sup>/Co<sup>3+</sup> charge ordering. For CoI $`3d`$ electrons, the majority-spin bands are fully occupied by three electrons in $`t_{2g}`$ and two electrons in $`e_g`$ bands. The minority-spin $`t_{2g}`$ bands are only partially occupied by two electrons, and one $`t_{2g}`$ and two $`e_g`$ bands are almost empty. Considering the pyramidal environment of Co, the occupied $`t_{2g}`$ states corresponds to $`d_{zx}`$ and $`d_{yz}`$ while the empty $`t_{2g}`$ to $`d_{xy}`$. Accordingly, the nominal valency and the $`3d`$ electron configuration at CoI site are assigned to be Co<sup>2+</sup> and $`3d^7`$ ($`t_{2g}^5e_g^2`$), respectively. Hence, Co<sup>2+</sup> ion is in the HS state with spin magnetic moment of $`\mu _S=3\mu _B`$ ($`S=3/2`$), which is consistent with the NPD data.
For CoII $`3d`$ electrons, the majority-spin bands are not fully occupied with a split-off empty $`e_g`$ state above $`E_\mathrm{F}`$ (see Fig. 3). For the minority-spin bands, the situation is similar to the case of CoI. It is thus possible to identify the $`3d`$ electron configuration of CoII as $`3d^6`$ ($`t_{2g}^5e_g^1`$) and the valency as Co<sup>3+</sup>. With one less electron than Co<sup>2+</sup>, the lower $`d_{3z^2r^2}`$ out of two $`e_g`$ states in the majority-spin bands is occupied and the upper $`d_{x^2y^2}`$ becomes empty. Hence, the spin state of Co<sup>3+</sup> ion is the IS state which is also in agreement with the experimental analysis of $`\mu _S=2\mu _B`$ ($`S=1`$).
The calculated charge occupancies of each Co $`3d`$ orbitals are shown in Table I. For Co<sup>2+</sup> ion, the majority-spin bands are almost completely occupied by $`t_{2g}=2.98`$ and $`e_g=1.99`$ electrons, while the minority-bands are only partially occupied by $`t_{2g}=2.02`$ and $`e_g=0.37`$ electrons. For Co<sup>3+</sup> ion, it is noticeable that the majority-spin $`e_g`$ states are partially occupied by $`e_g=1.54`$ electrons. Hence, the calculated charge occupancies are consistent with the nominal valencies of Co<sup>2+</sup> and Co<sup>3+</sup>, if one takes into account the band hybridization effects.
Vogt et al. have deduced magnetic moments from the NPD experiments as $`\mu _{\mathrm{exp}}=4.2\mu _B`$ and $`2.7\mu _B`$ for each Co<sup>2+</sup> and Co<sup>3+</sup> ion, respectively. In the analysis of the experiments, they counted only the spin moment contribution, assuming that the orbital moment is quenched. However, the orbital moment is only partially quenched in YBaCo<sub>2</sub>O<sub>5</sub>. In Table I, we have summarized the calculated magnetic moments using the LSDA + $`U`$ method. For Co<sup>2+</sup> ion, the spin and orbital magnetic moments are $`\mu _S=2.61\mu _B`$ and $`\mu _L=1.04\mu _B`$, respectively, and for Co<sup>3+</sup> ion, $`\mu _S=1.84\mu _B`$ and $`\mu _L=0.40\mu _B`$. The orbital moment of Co<sup>2+</sup> ion is as much as that of CoO, and the orbital moment of Co<sup>3+</sup> is comparable to that of NiO. The non-negligible orbital moment, which originates from the localized nature of Co $`3d`$ electrons, suggests that YBaCo<sub>2</sub>O<sub>5</sub> should fall in a class of strongly correlated electron system like CoO and NiO. The calculated total magnetic moments of $`\mu _{\mathrm{tot}}=3.65\mu _B`$ and $`2.24\mu _B`$ for Co<sup>2+</sup> and Co<sup>3+</sup> ion, respectively, are only slightly smaller than the experimental values. Evidently, this interpretation will also be valid for HoBaCo<sub>2</sub>O<sub>5</sub>. Suard et al. have improperly assigned the NPD measured $`\mu _{\mathrm{exp}}=3.7\mu _B`$ and $`2.7\mu _B`$ in HoBaCo<sub>2</sub>O<sub>5</sub> to Co<sup>3+</sup> and Co<sup>2+</sup> ion, respectively, However, their assumptions of spin-only moments and the HS states for both Co<sup>2+</sup> and Co<sup>3+</sup> ions are discarded by the present results. To test the calculated results, more experimental works like the x-ray scattering measurement are encouraged, in which separate determination of the spin and orbital moments are possible.
In Fig. 4, we have plotted the geometry of the orbital ordering which is obtained from the orbital dependent occupancy of the $`3d`$ minority-spin states at each Co site. At Co<sup>2+</sup> sites, the orbitals are aligned along $`a`$-axis, while at Co<sup>3+</sup> sites, the orbitals are along $`b`$-axis. This feature is understandable by considering that the bond length of $`d`$(Co<sup>2+</sup>-O) is larger in $`a`$-axis than in $`b`$-axis and vice versa for that of $`d`$(Co<sup>3+</sup>-O). As for the charge ordering configuration, Co<sup>2+</sup> and Co<sup>3+</sup> chains of a stripe type are formed in the $`ab`$ plane along $`b`$-axis, which are alternating in the $`a`$ and $`c`$ direction. This is in contrast to the charge ordering observed in the isostructural YBaMn<sub>2</sub>O<sub>5</sub>. In YBaMn<sub>2</sub>O<sub>5</sub>, the Mn<sup>2+</sup>/Mn<sup>3+</sup> orders in a checkerboard type. This difference gives rise to the different magnetic structures: $`G`$-type AFM phase for YBaCo<sub>2</sub>O<sub>5</sub> and $`G`$-type ferrimagnetic phase for YBaMn<sub>2</sub>O<sub>5</sub>. The theoretical result of the charge and orbital ordering geometry coincides with the experimentally proposed one. Thus, it can be inferred that the deformed bond lengths of $`d`$(Co-O) determine the charge and orbital ordering geometry.
The $`G`$-type AFM ordering in YBaCo<sub>2</sub>O<sub>5</sub> is consistent with the above charge and orbital ordering geometry. In the Co<sup>3+</sup> chains, the kinetic-exchange energy gain between the occupied $`d_{yz}`$ and $`d_{zx}`$ states and the empty $`d_{x^2y^2}`$ state for the AFM configuration of neighboring sites stabilizes the AFM spin ordering. In a similar way, the AFM ordering between neighboring Co<sup>3+</sup> and Co<sup>2+</sup> ions can be explained. The AFM ordering in the Co<sup>2+</sup> chains, however, is hard to understand in terms of the direct kinetic-exchange gain, because the overlap integral between two neighboring Co<sup>2+</sup> ions would be negligible as seen in Fig. 4. Instead, the AFM interaction in the Co<sup>2+</sup> chains is expected to be derived indirectly via the Co<sup>2+</sup>-Co<sup>3+</sup> and Co<sup>3+</sup>-Co<sup>3+</sup> AFM interactions.
As mentioned above, the structural transition plays a crucial role in determining the ground state properties of YBaCo<sub>2</sub>O<sub>5</sub>. Although the tetragonal to orthorhombic structural transition occurs simultaneously with the $`G`$-type AFM transition at $`T_\mathrm{N}330`$ K, the structural deformation is not significant above 220 K. Only near $`T_{\mathrm{CO}}220`$ K, the lattice splitting between $`a`$\- and $`b`$-axis becomes pronounced and the charge and orbital ordering emerges with the Co<sup>2+</sup> spin state transition from the low to high spin. Furthermore, it is known that the long-range charge ordering and the spin state transition are very sensitive to the oxygen stoichiometry. Therefore, the structural distortion is thought to be responsible for the orderings and the spin state transition by inducing different local environment at each Co ion site. This feature implies that the electron-lattice interaction is very important in this system. A detailed study on the electron-phonon interaction effects in YBaCo<sub>2</sub>O<sub>5</sub> is urgently demanded.
In conclusion, we have performed the LSDA + $`U`$ calculations for a new spin-charge-orbital-lattice coupled system YBaCo<sub>2</sub>O<sub>5</sub>. It is found that the Co<sup>2+</sup>/Co<sup>3+</sup> charge and orbital ordering and the Co<sup>2+</sup> HS state transition are closely correlated with the lattice distortion from the tetragonal to orthorhombic structure. The orbital moment has a substantially large contribution to the total magnetic moment. All of the effects of the Coulomb correlation, the spin-orbit coupling, and the electron-phonon interaction should be properly taken into account to understand physical properties of YBaCo<sub>2</sub>O<sub>5</sub>.
Acknowledgements$``$ This work was supported by the KOSEF (1999-2-114-002-5) and by the Brain Korea 21 Project.
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# Effects of Luminosity Functions Induced by Relativistic Beaming on Statistics of Cosmological Gamma-Ray Bursts
## 1 Introduction
The observation of high-redshift GRBs such as GRB970508, GRB971214, and GRB980703 (Metzger et al., 1997; Kulkarni et al., 1998; Djorgovski et al., 1998) confirmed the cosmological origin of GRBs (Paczy$`\stackrel{´}{\mathrm{n}}`$ski, 1986; Meegan et al., 1992). High redshifts ($`z>1`$) of GRBs would imply a much larger apparent energy budget than the conventional energy budget $`10^{5354}`$ergs unless some significant beaming is involved (Kumar, 1999; Kumar and Piran, 2000). The isotropic energies of observed GRBs with known redshifts range from $`10^{50}`$ergs to $`10^{54}`$ergs. A broad luminosity function can be a reasonable answer for this wide range, and can be naturally expected if we assume beaming of GRBs. The suggested correlation between GRBs and certain types of supernovae (Ib/c) implies that the total energy available for GRBs could be rather limited at least in the low energy GRBs $`L_{\mathrm{int}}10^{49}`$erg/sec (Wang and Wheeler 1998; Woosley et al. 1999, see also Kippen et al. 1998; Graziani et al. 1999). If we consider the supernova explosion as a universal GRB source, then the reported GRB-supernova correlation implies that the high redshift GRBs must be strongly beamed with the beam opening solid angle, which is approximately given by
$$d\mathrm{\Omega }(L_{\mathrm{int}}/L_{\mathrm{obs}})10^4(L_{\mathrm{int}}/10^{49}\mathrm{erg}/\mathrm{sec})(L_{\mathrm{obs}}/10^{53}\mathrm{erg}/\mathrm{sec})^1.$$
(1)
When we consider beaming, a broad luminosity range and a stringent constraint on the intrinsic energy are anticipated. Yi (1994) considered the effect of the luminosity function due to the cylindrical-beam with the inhomogeneous distribution of GRB sources on statistics of GRBs. He analytically showed that the luminosity function follows a power law with the index $`4/3`$ in the very narrow beam model when $`\alpha =1.0`$. His result has been confirmed by Mao and Yi (1994) and Chang and Yi (2000). Mao and Yi (1994) studied the statistical properties of luminosity functions derived using a conic beam with a uniform spatial distribution of GRB sources. They examined the cumulative probability distribution of the peak count rates. They found that the maximum redshift increases as a product of the Lorentz factor and the opening angle decreases. We attempt to generalize the luminosity function induced by the cylindrical-beam by introducing the conical shape in such a beam. If the progenitors of GRBs were produced in the late stage of massive stars (Woosley, 1993; Paczy$`\stackrel{´}{\mathrm{n}}`$ski, 1998; MacFadyen and Woosley, 1999), their spatial distribution may follow the star-formation rate (SFR) of massive stars (Totani, 1997; Wijers et al., 1998; Blain and Natarajan, 2000). We use the recently observed SFR data to derive the possible number density distribution of GRB sources (Steidel et al., 1999).
In this paper, we study how beaming-induced luminosity functions affect the statistics of the observed GRBs in the BATSE 4B catalog (Paciesas et al., 1999) considering the two different spatial distributions of GRB sources: an SFR-motivated distribution and the uniform distribution. We employ two different SFR-motivated distributions. We discuss the effect of the photon index $`\alpha `$ ($`=1.0`$ and 2.0) in the conlusion. We assume a flat universe with no cosmological constant, and adopt the Hubble constant $`H_{}=50\mathrm{k}\mathrm{m}/\mathrm{sec}/\mathrm{Mpc}`$ (cf. Yi 1994; Mao and Yi 1994). We analyze 775 observed bursts in the BATSE 4B catalog in order to improve the statistical significance. Assuming the bimodal distribution of durations of GRBs (Kouveliotou et al., 1993; Lamb et al., 1993; Mao, Narayan and Piran, 1994; Katz and Canel, 1996) we divide the total sample into two subgroups according to their duration and study the statistics of each subgroup.
The selection criteria applied to observational data are described in section 2. Luminosity functions due to beaming are discussed in section 3. We present results and discussions in section 4. We summarise our conclusion in section 5.
## 2 Classification of Observational Data
We adopt the bursts in the BATSE 4B catalog (Paciesas et al., 1999) and calculate their $`V/V_{\mathrm{max}}`$ for selected GRBs (Schmidt et al., 1988), using fluxes in channels 2 and 3. We choose bursts detected in 1024 ms trigger time scale with the peak count rates satisfying $`C_{\mathrm{max}}/C_{\mathrm{min}}1`$. Applying these criteria, we select 775 bursts among 915 bursts in the BATSE 4B $`C_{\mathrm{MAX}}/C_{\mathrm{MIN}}`$ table. The $`V/V_{\mathrm{max}}`$ value for this total sample is 0.3177 $`\pm `$ 0.0102. Then we divide the selected sample into two subgroups according to the burst durations ($`T_{90}`$), which is motivated by the reported bimodal structure in GRB durations (Kouveliotou et al., 1993; Lamb et al., 1993; Mao, Narayan and Piran, 1994; Katz and Canel, 1996). We define bursts with $`T_{90}>2.5\mathrm{sec}`$ as long bursts, and those with $`T_{90}<2.5\mathrm{sec}`$ as short bursts. Their estimated $`V/V_{\mathrm{max}}`$ values are 0.2901 $`\pm `$ 0.0113 (588 bursts) and 0.4178 $`\pm `$ 0.0239 (149 bursts), for long bursts and short bursts respectively. The number in parentheses indicates the number of bursts belonging to each subgroup. The $`V/V_{\mathrm{max}}`$ values obtained from the two subgroups show deviations from the values which are expected for a uniform distribution in a Euclidean space. In our subgroups, the long bursts have a larger deviation from the Euclidean value. It agrees with the claim by Tavani (1998). According to Tavani (1998), bursts with long durations ($`T_{90}>2.5\mathrm{sec}`$) and hard spectra ($`HR_{32}>3.0`$) show the largest deviation from the so-called Euclidean value.
## 3 Beaming and Luminosity Functions
We assume two kinds of number density distributions of GRB sources: (i) the SFR-motivated distribution $`n_{\mathrm{SFR}}(z)`$ and (ii) the uniform distribution $`n_\mathrm{u}(z)=n_{,\mathrm{u}}`$. We assume a two-sided Gaussian function as $`n_{\mathrm{SFR}}(z)`$ (Kim, Yang, and Yi, 1999), and fit $`n_{\mathrm{SFR}}(z)`$ to the extinction-corrected SFR data (Steidel et al., 1999). The observed SFR gradually decreases or even keeps flat over a wide redshift range beyond a certain redshift (cf. Madau et al. 1998). The SFR-motivated number density distribution is approximated by
$`n_{\mathrm{SFR}}(z)=`$ $`n_{,\mathrm{SFR}}\mathrm{exp}\left[{\displaystyle \frac{(zz_\mathrm{c})^2}{\mathrm{\Delta }z_{1}^{}{}_{}{}^{2}}}\right],`$ $`(z<z_\mathrm{c})`$
$`n_{,\mathrm{SFR}}\mathrm{exp}\left[{\displaystyle \frac{(zz_\mathrm{c})^2}{\mathrm{\Delta }z_{2}^{}{}_{}{}^{2}}}\right],`$ $`(z>z_\mathrm{c}),`$
where $`\mathrm{\Delta }z_1`$ and $`\mathrm{\Delta }z_2`$ are widths of the fitting functions for the left part and the right part and $`z_\mathrm{c}=1.5`$ is a critical redshift. All parameters are determined by fitting the observational data. We set $`\mathrm{\Delta }z_1=1`$ $`(z<z_\mathrm{c})`$ and we also effectively fix $`n_{\mathrm{SFR}}`$(z) to $`n_{,\mathrm{SFR}}`$ beyond $`z_\mathrm{c}`$. The proportional constants $`n_{,\mathrm{SFR}}`$ and $`n_{,\mathrm{u}}`$ are normalized to satisfy the number from the observational data.
### 3.1 Cylindrical Beaming
We consider the beaming-induced luminosity function following Yi (1993, 1994), where the relativistic beam has a constant cross-sectional area, i.e. a perfectly collimated cylindrical beam. Under this circumstance, a bulk motion of all the particles emitting the photons is parallel to the cone axis. In this case, the beam opening angle is irrelevant for the luminosity function. We assume that the beamed emission from ”standard bursts” which have the same intrinsic luminosity $`L_{\mathrm{int}}`$. The advantage of this model is that it physically defines a range of the luminosity, i.e. $`[L_{\mathrm{min}},L_{\mathrm{max}}]`$, for a given photon index $`\alpha `$, and the Lorentz factor $`\gamma `$. An apparent luminosity function $`\mathrm{\Phi }(L)`$ due to the uniformly distributed beam in space is given by
$$\mathrm{\Phi }(L)=\frac{1}{p\beta \gamma }L_{\mathrm{int}}^{1/p}L^{(p+1)/p}$$
(3)
where $`\beta =(1\gamma ^2)^{1/2}`$ and $`p=\alpha +2`$. We adopt $`\alpha =1.0`$ (Mallozzi et al., 1996) and 2.0 (Yi, 1994) and the Lorentz factor $`\gamma =100`$ (see Piran 1999). If we assume a bi-polar jet, the minimum and the maximum luminosities are naturally given as follows:
$$Ł_{\mathrm{min}}=L_{\mathrm{int}}B_{\mathrm{min}}^p,L_{\mathrm{max}}=L_{\mathrm{int}}B_{\mathrm{max}}^p,$$
(4)
where $`B_{\mathrm{min}}=\gamma ^1`$ and $`B_{\mathrm{max}}=[\gamma (1\beta )]^1`$ (Blandford and K$`\ddot{\mathrm{o}}`$nigle, 1979). The intrinsic luminosity $`L_{\mathrm{int}}`$ for each subgroup is determined by the $`V/V_{\mathrm{max}}`$ test. The obtained $`L_{\mathrm{int}}`$ values are $`10^{43}\mathrm{erg}/\mathrm{sec}`$ for the specific set of adopted beaming parameters. This is a much smaller value than that required in the standard candle model, which is $`L_{}10^{51}\mathrm{erg}/\mathrm{sec}`$.
We calculate two statistical quantities using the beaming-induced luminosity function explained in this subsection. First, we calculate the $`V/V_{\mathrm{max}}`$ as a function of threshold flux $`F_{\mathrm{th}}`$ (Che et al., 1999). The threshold flux is obtained when we put the maximum detectable redshift $`z_{\mathrm{max}}`$ in
$$F=\left(\frac{H_0^2L}{16\pi c^2}\right)\frac{(1+z)^{1\alpha }}{[(1+z)^{1/2}1]^2},$$
(5)
where $`L`$ is the luminosity, $`c`$ is the speed of light, $`H_0`$ is the Hubble constant, and $`\alpha `$ is the photon index. After we obtain a set of $`F_{\mathrm{th}}`$, we calculate the $`V/V_{\mathrm{max}}`$ for each threshold flux. The $`V/V_{\mathrm{max}}`$ curve converges to 0.5 as $`F_{\mathrm{th}}`$ increases sufficiently large. We convert flux to the number of bursts and plot $`V/V_{\mathrm{max}}`$ as a function of the number of bursts for all bursts and long and short bursts in Fig. 1 (a) for $`\alpha =1.0`$ (thin lines) and $`\alpha =2.0`$ (thick lines). The $`V/V_{\mathrm{max}}`$ vs. number of bursts relation essentially gives the same information that one obtains from the conventional logN $``$ logP plot. Our $`V/V_{\mathrm{max}}`$ test does not include the threshold effect due to the detection efficiency.
Second, we calculate the fraction of GRBs. We define the fraction of bursts located at a redshift larger than $`\mathrm{z}^{}`$ as
$$f_{>\mathrm{z}^{}}=\frac{\mathrm{Number}\mathrm{of}\mathrm{bursts}\mathrm{with}\mathrm{a}\mathrm{redshift}\mathrm{larger}\mathrm{than}\mathrm{z}^{}}{\mathrm{Total}\mathrm{number}\mathrm{of}\mathrm{bursts}}.$$
(6)
In Fig. 1 (b), we plot the fraction of bursts as a function of redshift for each burst subsample for the cylindrical-beam case with different number density distributions.
### 3.2 Conic Beaming
As a more realistic example, we adopt conically-beamed emission from GRBs, which was studied by Mao and Yi (1994), and extended by Chang and Yi (2000). We consider the luminosity function produced using three different opening angles, i.e. $`\mathrm{\Delta }\theta =0^{}.1,1^{}.0`$, and $`3^{}.0`$. We calculate $`V/V_{\mathrm{max}}`$ and the $`f_{>\mathrm{z}^{}}`$ for each subgroup with the derived luminosity function. The bulk Lorentz factor is set to be 100. According to Mao and Yi (1994), the probability that we observe the bright bursts rapidly increases as the opening angle increases, while relatively dim bursts are not as detectable as brighter ones. Therefore, as the opening angle increases, the derived luminosity function becomes similar to that of the standard candle case, in which all bursts have the same maximum luminosity given by this luminosity function. When the opening angle is small $`(\mathrm{\Delta }\theta 1/\gamma )`$, the luminosity function gives the same result as the cylindrical beaming case. Using this luminosity function, we calculate the same statistics as we did in previous subsection. Results are shown in Figs. 2-5 for $`\alpha =1.0`$ (thin lines) and $`\alpha =2.0`$ (thick lines) in the cases of each sample.
## 4 Results and Discussions
The left panels in Figs. 1-5 (a) show the $`V/V_{\mathrm{max}}`$ as a function of the number of bursts. We show results obtained from different luminosity distributions and number density distributions. The thin lines in Fig. 1-4 represent the results obtained with $`\alpha =1.0`$ and thick lines $`\alpha =2.0`$. In Figs. 1-4, we assume that the SFR-motivated number density distribution, which is flat beyond a redshift larger than the critical redshift, $`z_\mathrm{c}`$. Finally, we show our results obtained with the SFR-motivated number density distribution, which gradually decreases and $`\alpha =1.0`$ in Fig. 5 (a). In all cases, we assume that n(z) rapidly increases if $`z<z_\mathrm{c}`$. The dotted lines represent the observed $`V/V_{\mathrm{max}}`$ curves with the $`\pm 3\sigma `$ bound. The solid lines indicate the theoretical results derived from our luminosity functions with the SFR-motivated number density distribution of GRB sources. We plot the results obtained for the luminosity functions with the uniform distribution of burst sources as dashed lines, for comparaison. Apparently all the luminosity functions we have studied satisfy the observed $`V/V_{\mathrm{max}}`$ curve at the 3$`\sigma `$ significance level. Although the number density distribution of GRB sources is believed to correlate with the SFR, we cannot rule out the uniform distribution of burst sources in this analysis (see also Krumholz et al. 1998).
The right panels in Figs. 1-5 (b) show $`f_{>\mathrm{z}^{}}`$ in terms of $`z^{}`$. The y-axis corresponds to the fraction of bursts which have redshifts larger than a certain redshift $`z^{}`$. Solid lines result from the SFR-motivated number density distribution for the spatial distribution of GRBs and dashed lines the uniform distribution. Filled squares in Figs. 1, 3, and 5 (b) locate the known redshifts of 15 observed GRBs so far. The redshift values are quoted from http://cossc.gsfc.nasa.gov/batse/counterparts/GRB\_table.html. We compare the observed redshift values only with the long subsample because all GRBs with the known redshifts belong to the long bursts. It is interesting to note that, if we assume $`\alpha =1.0`$, a large fraction of GRBs are distributed at high redshifts, provided that GRBs are assumed to follow the SFR, regardless of specific beaming models. However, in the case of $`\alpha =2.0`$, $`z_{\mathrm{max}}`$ is significantly reduced, though it still remains quite large. The fraction for $`f_{>\mathrm{z}^{}}`$ at $`z^{}=3.42`$ for a cylindrical-beam case is $`75\%`$ for long and total samples when $`\alpha =1.0`$. This is too large a value given the fact that the observed GRBs with such high redshifts account for $`10\%`$ of the total bursts. This fraction become $`<50\%`$ when $`\alpha =2.0`$. It is also interesting to note that, though their $`V/V_{\mathrm{max}}`$ value is close to a Euclidean value, most of the short bursts are distributed at high redshifts in the case of $`\alpha =1.0`$. When we assume that the long bursts are uniformly distributed in space, the $`z_{\mathrm{max}}`$ estimate for the broad beaming cannot explain the largest observed redshift in both cases of $`\alpha =1.0`$ and 2.0.
Our study is based on an assumption that all GRBs have the same intrinsic luminosity and they are all beamed. Because the $`V/V_{\mathrm{max}}`$ test is not sufficiently sensitive to the observed data, different luminosity functions are essentially indistinguishable. If there is an intrinsic luminosity function and the degree of beaming of GRBs is moderate, we may obtain results similar to what we have presented in this study. If all GRBs are highly beamed and their sources follow a SFR-like distribution, the maximum detectable redshift becomes very large, even in the case that the observed $`V/V_{\mathrm{max}}`$ is sufficiently close to the Euclidean value, i.e., $`0.5`$. Therefore, it is difficult to rule out a possibility that the apparent Euclidean value of $`V/V_{\mathrm{max}}`$ may be due to the luminosity function, for instance, induced by beaming. In other words, the so-called Euclidean value may have nothing to do with the Euclidean distribution. If the total burst sample is a mixture of broad (or nearly non-beamed) beams and narrow beams, the exact fraction of the strongly beamed bursts may affect the burst rate estimate significantly. This is beyond the scope of this paper.
For a given Lorentz factor, the range of the derived beamed luminosity distribution becomes broader as the beam opening angle decreases. We obtain $`\mathrm{log}L_{\mathrm{max}}/L_{\mathrm{min}}13`$ in the cylindrical-beam case. However, in the case of broad beaming, $`\mathrm{log}L_{\mathrm{max}}/L_{\mathrm{min}}`$ becomes substantially smaller. This is only an effective range because the cut-off for the minimum luminosity given by the conic beam is rather loosely constrained. It is obvious that the luminosity range decreases when the opening angle increases for a fixed Lorentz factor. This is simply explained as follows: Assuming that we observe a beamed emission with an opening angle $`\mathrm{\Delta }\theta `$, if we define $`\psi `$ as an angle between the line of sight and the symmetry axis of the conic beam, we can observe the beamed emission only within a certain range of $`\psi `$, $`\mathrm{\Delta }\psi `$. For instance, when $`1/\gamma \mathrm{\Delta }\theta `$, $`\mathrm{\Delta }\psi `$ is approximately $`\mathrm{\Delta }\theta `$. For a given Lorentz factor, $`\mathrm{\Delta }\psi `$ increases as the opening angle of cone becomes larger. The apparent luminosity range from a broad beam would decrease gradually with $`\psi `$, and in the extreme case, the isotropic radiation would have constant luminosity regardless of $`\psi `$. However, when $`\mathrm{\Delta }\theta `$ is very small, the apparent luminosity rapidly decreases with $`\psi `$, and the derived luminosity range is narrower than a broad-beam case.
## 5 Conclusion
We have demonstrated that the beaming-induced luminosity function may explain the statistics of the observed GRBs. In the case of the cylindrically beamed luminosity function, we determined the intrinsic luminosity $`L_{\mathrm{int}}`$ by the $`V/V_{\mathrm{max}}`$ test for each subsample. In all cases, the intrinsic luminosity is much smaller than that of non-beamed luminosity functions, i.e. $`L10^{5253}\mathrm{erg}/\mathrm{sec}`$. The maximum luminosity obtained from our luminosity functions, in both the cylindrical-beam and the conic-beam cases, is $`L_{\mathrm{max}}10^{5051}\mathrm{erg}/\mathrm{sec}`$. This is compatible with the required value for isotropic radiation. The ratio between the maximum luminosity for the long bursts and that for the short bursts $`L_{\mathrm{max},\mathrm{long}}/L_{\mathrm{max},\mathrm{short}}`$ is $`7`$ for the narrow beam and $`3`$ for the broad beam, respectively. This implies that these two subgroups may have two different intrinsic luminosities, and hence probably two different origin. The obtained $`z_{\mathrm{max}}`$ in the cylindrical-beam case are $`10,14,3`$ ($`\alpha =1.0`$) and $`4.5,5.5,1.6`$ ($`\alpha =2.0`$) for the total selected sample, and long and short subgroup, respectively. The obtained $`z_{\mathrm{max}}`$ in the broader beam case (i.e. $`\mathrm{\Delta }\theta =3^{}.0`$) are $`4,6,2`$ ($`\alpha =1.0`$) and $`2.2,2.6,1.2`$ ($`\alpha =2.0`$) for the total selected sample, and long and short subgroup, respectively. The maximum known redshift we consider, $`z=3.42`$, is estimated for GRB971214. However, even a larger value of $`z<3.9`$ is suggested for GRB980329 (http://cossc.gsfc.nasa.gov/batse/counterparts/GRB$`\mathrm{\_}`$table.html). Without any detailed information on the redshift distribution of GRBs, the maximum redshift gives a rather strong constraint on the ratio between the Lorentz factor and the opening angle $`\mathrm{\Delta }\theta `$. For the conic beam case, the luminosity function derived for the narrow opening angle seems to fit the $`V/V_{\mathrm{max}}`$ curve better.
From this simple analysis we have carried out, we have shown that the beaming-induced luminosity function may account for the basic statistical properties of the observed GRBs. The apparent Euclidean distribution of GRBs may be an indication of the presence of the GRB luminosity function.
Although it is the beyond the scope of this work, it is interesting to point out that there could be a potentially observable correlation between spectral hardness and luminosity since the most Doppler-boosted burst would probe the hardest part of the spectrum. Any wide scatters in intrinsic spectral shapes would make the detection of this correlation complicated.
We thank H. Kim and K. Kwak for useful discussions. IY is supported in part by the KRF grant No. 1998-001-D00365. We appreciate the referee, Ralph A. M. J. Wijers, for his suggestions and some useful information.
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# 1 Introduction
## 1 Introduction
The principal result we shall present in this paper is a physical system which at the outset is not related to gravity but which nevertheless requires curved spacetime for its very existence. This situation is best illustrated with the example of a magnetic monopole in the framework of both electrodynamical and gravitational topologically massive theories in $`3`$-dimensions. We find that the essential new feature introduced by topological mass is to open up the Dirac string of a monopole into a cone. The intuitive example of this phenomenon takes place for Maxwell-Chern-Simons (MCS) theory in a Riemannian $`3`$-manifold with Euclidean signature which shows that solutions of the MCS field equations naturally lead us into de Sitter (dS) space with conical deficit. Three dimensional flat spacetimes with, or without conical deficit do not allow such a solution.
In $`4n+3`$ dimensions there exists the Chern-Simons action through which we can introduce topological mass into Maxwell’s electrodynamics and Einstein’s gravity . In the simplest case of $`n=0`$ pure Einstein gravity has no propogating degrees of freedom and no Newtonian limit . On the other hand, three dimensional gravity with a pure Chern-Simons action is equivalent to the Yang-Mills gauge theory of the conformal group and therefore is finite and exactly solvable . There is, however, a very interesting non-trivial theory of gravitation in $`2+1`$ dimensions which has been proposed by Deser, Jackiw and Templeton (DJT) where the gravitational Chern-Simons action is added to the Hilbert action. This is the theory of topologically massive gravity (TMG). It is a dynamical theory of gravity unique to three dimensions and the geometry of its exact solutions is non-trivial. Mathematically the DJT field equations pose an interesting challenge in that they are qualitatively different from the Einstein field equations while posessing their elegance and consistency. Clement has made the most thorough investigation of the solutions of DJT field equations for TMG as well as TME which uncovered many interesting effects due to topological mass. Self-dual solutions of TME coupled to Einsteinian gravity were discussed by Fernando and Mansouri and by Dereli and Obukhov who gave the general analysis.
This class of fields we shall consider falls outside the domain of solutions considered earlier - and illustrate in its purest form some of the new interesting effects that take place in the presence of topological mass. Earlier we presented the spinor formulation of TMG in terms of real $`2`$-component spinors which provides a very useful formalism analogous to the Newman-Penrose formalism of general relativity. We shall extend this formalism to include topologically massive electrodynamics and gravity. This formalism is helpful for constructing physically meaningful exact solutions of the coupled DJT-MCS field equations. We shall use it to derive the exact solution for a topologically massive magnetic monopole.
## 2 Dirac Monopole
It will be useful to start our considerations with a brief review of the Dirac monopole and its extension to TME in order to explain the essential idea we shall use throughout this paper. Maxwell’s electrodynamics is given by the action principle
$$I_M=\frac{1}{2}(F\frac{1}{2}dA)^{}F$$
(1)
in Minkowski spacetime and leads to the Maxwell field equations
$`F`$ $`=`$ $`dA,dF=0`$ (2)
$`d^{}F`$ $`=`$ $`0`$ (3)
the first one of which is the second Bianchi identity. Dirac pointed out that for a magnetic monopole eqs.(2) must fail at least in one point on every Gaussian surface enclosing the monopole. The set of all such points forms the Dirac string. The Maxwell potential that satisfies these requirements is a generalization of the $`1`$-form obtained for the polar angle $`d\varphi `$ on the plane which is closed but not exact. The $`U(1)`$ potential $`1`$-form and the field $`2`$-form for the Dirac monopole are given by
$`A`$ $`=`$ $`g\left[1cos\theta \right]d\varphi ,`$ (4)
$`F`$ $`=`$ $`g\mathrm{sin}\theta d\theta d\varphi ,`$ (5)
the latter of which is the familiar element of area on $`S^2`$. The semi-infinite Dirac string is at $`\theta =0`$ and the surface integral
$$F=4\pi g$$
(6)
determines the monopole magnetic charge.
We shall now consider the Euclidean Maxwell-Chern-Simons topologically massive electrodynamics in order to illustrate the essential new idea brought in by making the Dirac monopole topologically massive. With the inclusion of the electromagnetic Chern-Simons term the action is given by
$$I_{MCS}=\frac{1}{2}\left\{(F\frac{1}{2}dA)^{}F\nu dAA\right\}$$
(7)
which yields the MCS field equation
$$d^{}F=\nu F$$
(8)
and the Bianchi identity (2) where $`\nu `$ is a coupling constant, the electromagnetic topological mass. In order to satisfy these field equations with a $`U(1)`$ potential $`1`$-form satisfying the properties of the Dirac monopole (4) we must introduce a deficit in the polar angle $`\theta `$ that led to the Dirac string. That is, topological mass has the effect of turning the string into a cone. Thus we should consider a field $`2`$-form of the type
$$F=g\mathrm{sin}(b\theta )d\theta d\varphi $$
(9)
where $`b`$ is a constant deficit parameter which will be related to topological mass. Now it is clear that the potential $`1`$-form (4) must be modified but still lead to eq.(9) as the field. This suggests that we consider a Riemannian manifold with the co-frame consisting of a modified form of the left-invariant $`1`$-forms of Bianchi Type $`IX`$ parametrized in terms of Euler angles
$`\sigma ^0`$ $`=`$ $`d\psi +\mathrm{cos}(b\theta )d\varphi `$
$`\sigma ^1`$ $`=`$ $`\mathrm{sin}(b\psi )d\theta +\mathrm{cos}(b\psi )\mathrm{sin}(b\theta )d\varphi `$ (10)
$`\sigma ^2`$ $`=`$ $`\mathrm{cos}(b\psi )d\theta +\mathrm{sin}(b\psi )\mathrm{sin}(b\theta )d\varphi `$
satisfying the Maurer-Cartan equations
$$d\sigma ^i=\frac{1}{2}C_{jk}^i\sigma ^j\sigma ^k$$
(11)
with non-vanishing structure constants
$$C_{\mathrm{\hspace{0.33em}\hspace{0.33em}12}}^0=C_{\mathrm{\hspace{0.33em}\hspace{0.33em}21}}^0=b,C_{\mathrm{\hspace{0.33em}\hspace{0.33em}20}}^1=C_{\mathrm{\hspace{0.33em}\hspace{0.33em}02}}^1=b,C_{\mathrm{\hspace{0.33em}\hspace{0.33em}01}}^2=C_{\mathrm{\hspace{0.33em}\hspace{0.33em}10}}^2=b.$$
(12)
Then a $`U(1)`$ potential $`1`$-form of the type
$$A=g\sigma ^0$$
(13)
will have all the desired properties and lead to the field $`2`$-form (9). The clue to the satisfaction of the field equation (8) for TME lies in the fact that with the co-frame (10) the Cartan-Killing metric $`ds^2=\eta _{ik}\sigma ^i\sigma ^k`$ with $`\eta _{ik}=diag.(1,1,1)`$ becomes
$$ds^2=d\theta ^2+d\varphi ^2+d\psi ^2+2\mathrm{cos}(b\theta )d\psi d\varphi $$
(14)
which is simply de Sitter space with the polar angle suffering a defect. The duality relations for the basis (10) immediately leads to the result that for the potential $`1`$-form (13) eqs.(8) of TME will be satisfied identically in dS provided
$$b=\nu ,$$
(15)
the deficit in the Eulerian polar angle is identified with topological mass. From the curvature of (14) we find
$$\lambda =\frac{\nu ^2}{4}$$
(16)
relating the cosmological constant to electromagnetic topological mass. For the case of Lorentzian signature, c.f. section 6, this would be anti-de Sitter spacetime. The field $`2`$-form (9) is shown in figure 1 where because we are in the Euclidean sector the deficit in the polar angle caused by topological mass can be given an explicit illustration. The maximal analytical extension of dS space is given in the chart
$`\mathrm{cos}\left({\displaystyle \frac{\nu }{2}}\theta \right)\mathrm{cos}\left[{\displaystyle \frac{\nu }{2}}(\varphi +\psi )\right]`$ $`=`$ $`\mathrm{tanh}(\nu \alpha ),`$
$`\mathrm{sin}\left({\displaystyle \frac{\nu }{2}}\theta \right)`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}(\nu \beta )}{\mathrm{cosh}(\nu \alpha )}},`$ (17)
$`{\displaystyle \frac{1}{2}}(\varphi \psi )`$ $`=`$ $`\gamma `$
whereby the metric (14) becomes
$$d\stackrel{~}{s}^2=\mathrm{\Omega }^2\left(d\alpha ^2+d\beta ^2+\mathrm{sin}^2(\nu \beta )d\gamma ^2\right)$$
(18)
with the conformal factor $`\mathrm{\Omega }=\frac{1}{2}\mathrm{cosh}(\nu \alpha )`$. The incomplete Einstein static cylinder is manifest in eq.(18).
In the discussion of TME monopole we started out with MCS field equations (2) and (8) which are written in a general background. The expectation was that these field equations will admit a solution in flat background spacetime for a physical system which is electrodynamic in nature and a priory completely unrelated to gravity. This proved to be impossible. With the missing cone in the field $`2`$-form (9) eqs.(2) and (8) could only be satisfied in curved space (14) with a corresponding conical deficit.
Thus we arrive at a remarkable conclusion that a physical system of electrodynamic type requires curved spacetime for its existence.
## 3 Spinor formalism in $`(2+1)`$-dimensions
We shall now turn our attention to Lorentz signature and introduce the Newman-Penrose version of TME equations. This type of study for topologically massive gravity was given by Hall, Morgan and Perjes and its $`2`$-component spinor description with differential forms was constructed in which will henceforth be referred to as $`𝐈`$. Here we shall extend this work by first presenting the spinor formulation of TME and in section 5 couple it to TMG.
We begin by recalling some basic relations from $`𝐈`$. At each point of a three dimensional space-time with the metric of Lorentz signature we can introduce a pair of real 2-component spinors
$$\zeta _{(a)}^A=\{o^A,\iota ^A\}A=1,2a=0,1$$
(19)
which will define the basis. Spinor indices will be raised and lowered by the Levi-Civita symbol $`ϵ_{AB}`$ from the right and the normalization of the spin frame is given by $`o_A\iota ^A=1`$ with all others vanishing identically. It is evident that such a spin frame will imply the triad of real basis vectors which can be connected to basis spinors through the Infeld-van der Waerden symmetric quantities $`\sigma _{AB}^i`$. We recall that the co-frame is given by
$$\sigma _i^{AB}dx^i=\left(\begin{array}{cc}n& \frac{1}{\sqrt{2}}m\\ \frac{1}{\sqrt{2}}m& l\end{array}\right)$$
(20)
and the space-time metric is
$$ds^{\mathrm{\hspace{0.33em}2}}=ln+nlmm$$
(21)
where $`l,n`$ are null and $`m`$ is space-like. The Newman-Penrose intrinsic derivative operators in the direction of the legs $`l^i,n^i,m^i`$ of the triad are given by $`D,\mathrm{\Delta },\delta `$ respectively so that
$$d=l\mathrm{\Delta }+nDm\delta $$
(22)
is resolution of the exterior derivative along the legs of the triad. We recall that taking the exterior derivative of the basis 1-forms and expressing the result in terms of the basis 2-forms yields the spin coefficients through the solution of a linear algebraic system. The result is given by eqs.(I. 23)
$$\begin{array}{ccc}dl& =\hfill & ϵln+(\alpha \tau )lm\kappa nm\hfill \\ dn& =\hfill & ϵ^{}ln\kappa ^{}lm(\alpha +\tau ^{})nm\hfill \\ dm& =\hfill & (\tau ^{}\tau )ln\sigma ^{}lm\sigma nm\hfill \end{array}$$
(23)
from which the spin coefficients are obtained through the solution of a linear algebraic system. Here prime denotes the symmetry operation resulting from the interchange of $`l`$ and $`n`$ leaving $`m`$ fixed. We note that $`\alpha ^{}=\alpha `$.
Earlier we had not introduced the spinor equivalent of the basis 2-forms which are wedge products of the Infeld-van der Waerden matrices of basis 1-forms (20)
$$L^{AXBY}=\sigma ^{AX}\sigma ^{BY}$$
(24)
with the spinor equivalent
$$L^{AXBY}=L^{AB}ϵ^{XY}+L^{XY}ϵ^{AB}$$
(25)
due to skew symmetry in the pair of indices $`AX`$ and $`BY`$ and the basic spinor relation (I. 11). The 2-component spinors $`L^{AB}`$ and $`L^{XY}`$ are real and symmetric $`L^{AB}=L^{(AB)}`$ $`L^{XY}=L^{(XY)}`$ and we have
$`L^{00}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}nm`$
$`L^{01}`$ $`=`$ $`{\displaystyle \frac{1}{2}}ln`$ (26)
$`L^{11}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}lm`$
for the basis $`2`$-forms. Using the definition of Hodge star operator (I. 57) and the completeness relation (I. 8) we find that
$${}_{}{}^{}(ln)=m,^{}(lm)=l,^{}(nm)=n$$
(27)
determine the duals of the basis 2-forms in (2+1)-dimensions. All other duality relations can be obtained from $`{}_{}{}^{}=1`$.
## 4 Topologically massive electrodynamics
In general relativity the spinor approach has turned out to be very useful for the investigation of physically interesting solutions of the Einstein and Maxwell equations. This should be the case for three dimensional spacetimes as well. So we shall now derive the TME equations in the spinor formalism. We recall that the field $`2`$-form is given by
$$F=\frac{1}{2}F_{ik}dx^idx^k=\frac{1}{2}F_{AXBY}\sigma ^{AX}\sigma ^{BY}$$
(28)
in terms of the basis 2-forms (26). Using the same considerations that led to (25) we decompose the electromagnetic spinor
$$F_{AXBY}=\phi _{AB}ϵ_{XY}+\phi _{XY}ϵ_{AB}$$
(29)
where $`\phi _{AB}`$ and $`\phi _{XY}`$ are symmetric second rank real 2-spinors. Taking into account this relation in (28) together with (25) and (26) we find that
$$F=\phi _0mn\phi _1ln+\phi _2lm$$
(30)
where
$`\phi _0:`$ $`=`$ $`\sqrt{2}\phi _{00}=F_{ik}l^im^k`$
$`\phi _1:`$ $`=`$ $`2\phi _{01}=F_{ik}l^in^k`$ (31)
$`\phi _2:`$ $`=`$ $`\sqrt{2}\phi _{11}=F_{ik}m^in^k.`$
are the three real triad scalars of the electromagnetic field. Under the action of the prime operation the triad scalars undergo the transformation $`\phi _1\phi _1,\phi _0\phi _2`$. We also need the dual of the field 2-form (30) which is given by
$$^{}F=\phi _0n+\phi _1m\phi _2l$$
(32)
due to the relations (27). Using eqs.(30) and (32) together with (22) and (23) the field eqs.(8) assume the form
$`(\delta \alpha \tau ^{}\nu )\phi _0(D\sigma )\phi _1\kappa \phi _2`$ $`=`$ $`0,`$
$`(\mathrm{\Delta }\sigma ^{})\phi _1(\delta +\alpha \tau +\nu )\phi _2+\kappa ^{}\phi _0`$ $`=`$ $`0,`$ (33)
$`(\mathrm{\Delta }+ϵ^{})\phi _0(D+ϵ)\phi _2(\tau ^{}\tau +\nu )\phi _1`$ $`=`$ $`0`$
and the Bianchi identity (2) is given by
$$(\delta \tau \tau ^{})\phi _1(D+ϵ\sigma )\phi _2(\mathrm{\Delta }+ϵ^{}\sigma ^{})\phi _0=0.$$
(34)
These are the Newman-Penrose version of TME equations. We note that under the prime operation the first two equations in (33) go over into each other, provided that the sign of the TME coupling constant is also changed simultaneously, whereas the last equation in (33), as well as the Bianchi identity (34) remain invariant.
Finally, we note that the Maxwell stress tensor
$$T_{ik}=F_i^jF_{kj}+\frac{1}{4}g_{ik}F^{mn}F_{mn}$$
(35)
can be expressed in terms of the electromagnetic triad scalars as follows
$`T_{ik}=\phi _0^2n_in_k+\phi _1^2l_{(i}n_{k)}+\phi _2^2l_il_k+(\phi _0\phi _2+{\displaystyle \frac{1}{2}}\phi _1)m_im_k`$
$`2\phi _0\phi _1n_{(i}m_{k)}2\phi _1\phi _2l_{(i}m_{k)}`$ (36)
where round parantheses denote symmetization. The Chern-Simons term makes no contribution to the energy-momentum tensor.
## 5 Topologically massive gravity with sources
Deser, Jackiw and Templeton’s theory of topologically massive gravity overcomes the dynamically trivial character of Einsteinian gravity in three dimensions. In our previous work I we wrote the DJT field equations in terms of differential forms with triad scalar coefficients. Here we shall extend this formalism to topologically massive gravity with sources, in particular TME. It is a property of $`3`$ dimensions that symmetric second rank tensors can be written as $`2\times 2`$ matrices of $`2`$-forms which enables us to write the field equations in compact form. For this purpose we shall start by constructing the energy-momentum $`2`$-form in three dimensions.
### 5.1 Energy-momentum 2-form
The spinor equivalent of the symmetric energy-momentum tensor $`T_{ik}`$ admits the decomposition
$$T_{ABXY}=\frac{1}{2}\left(T_{ABXY}+T_{ABYX}\right)+\frac{1}{2}\left(T_{BAYX}T_{ABYX}\right)$$
(37)
where the first paranthesis is symmetric in pair of indices $`A`$, $`B`$ and $`X`$, $`Y`$, while the second one is skew in the same pair of indices. Applying the basic spinor relation (I. 11) to the second group of indices we obtain
$$T_{AXBY}=S_{ABXY}+\frac{1}{3}Tϵ_{AB}ϵ_{XY}$$
(38)
where
$$T=T_{AX}^{AX}=T_i^i$$
and
$$S_{ABXY}=T_{(AB)XY}=T_{AB(XY)}=T_{(AB)(XY)}$$
is the trace-free part of the energy-momentum tensor. Next, we shall use the trace-free part of the energy-momentum tensor to construct the spinor valued energy-momentum 2-form
$$T_A^B=S_{AXY}^B\mathrm{\Sigma }^{XM}\mathrm{\Sigma }_M^Y$$
(39)
where
$$\mathrm{\Sigma }_A^B=\frac{1}{\sqrt{2}}\sigma _{Ai}^Bdx^i.$$
(40)
are obtained by by lowering a spinor index in eqs.(20). We can now express the components of the energy-momentum 2-form in terms of triad scalars. We have
$`2T_0^{\mathrm{\hspace{0.33em}\hspace{0.33em}0}}`$ $`=`$ $`\left(T_{02}+{\displaystyle \frac{1}{2}}T\right)lnT_{12}lm+T_{01}nm`$
$`\sqrt{2}T_0^{\mathrm{\hspace{0.33em}\hspace{0.33em}1}}`$ $`=`$ $`T_{01}ln+{\displaystyle \frac{1}{2}}T_{02}lmT_{00}nm`$ (41)
$`\sqrt{2}T_1^{\mathrm{\hspace{0.33em}\hspace{0.33em}0}}`$ $`=`$ $`T_{12}lnT_{22}lm+{\displaystyle \frac{1}{2}}T_{02}nm`$
where we have introduced the definitions
$$\begin{array}{cccccccccc}T_{00}\hfill & :=\hfill & S_{0000},\hfill & T_{01}\hfill & :=\hfill & \sqrt{2}S_{0010},\hfill & T_{02}\hfill & :=\hfill & 2S_{0011},\hfill & \\ T_{11}\hfill & :=\hfill & S_{0101},\hfill & T_{12}\hfill & :=\hfill & \sqrt{2}S_{0111},\hfill & T_{22}\hfill & :=\hfill & S_{1111}\hfill & \end{array}$$
(42)
which consist of the triad scalars
$$\begin{array}{ccccccccc}T_{00}\hfill & =\hfill & T_{ik}l^il^k,\hfill & T_{01}\hfill & =\hfill & T_{ik}l^im^k,\hfill & T_{02}\hfill & =\hfill & T_{ik}m^im^k,\hfill \\ T_{11}\hfill & =\hfill & T_{ik}l^in^k\frac{1}{3}T,\hfill & T_{12}\hfill & =\hfill & T_{ik}n^im^k,\hfill & T_{22}\hfill & =\hfill & T_{ik}n^in^k\hfill \end{array}$$
(43)
of the energy-momentum tensor. We note that under the prime operation the index 1 remains unchanged while $`02`$.
### 5.2 DJT field equations with sources
The DJT field equations of TMG with sources are given by
$$G^{ik}+\frac{1}{\mu }C^{ik}=\lambda g^{ik}\text{æ}T^{ik}$$
(44)
where $`G^{ik}`$ is the Einstein tensor and $`C^{ik}`$ is Cotton’s conformal tensor of three-dimensional manifolds. The constants $`\mu `$ and æ are the DJT topological and Einstein matter coupling constants with $`\lambda `$ standing for the cosmological constant. The sign of the matter coupling constant is taken to be negative, in contrast to four-dimensional gravity, to choose the physical non ghost-like excitation modes. The above definition of the matrix of energy-momentum 2-form along with the curvature and Cotton 2-forms described by eqs.(I. 41) and (I. 55) enable us to write the field equations (44) in the form
$$R_A^B+\frac{1}{\mu }C_A^B+\left(\lambda \frac{1}{2}\text{æ}T\right)\mathrm{\Sigma }_A^M\mathrm{\Sigma }_M^B=\text{æ}T_A^B$$
(45)
that consist of matrices of differential $`2`$-forms with triad scalar entries. The expression for the DJT field equations follows from the substitution of the results in eqs.(41), (I. 42) and (I. 60-62) into eq.(45). Thus we arrive at the following set of DJT field equations
$`D\mathrm{\Phi }_{12}\mathrm{\Delta }\mathrm{\Phi }_{10}3(\tau \tau ^{})\mathrm{\Phi }_{11}ϵ^{}\mathrm{\Phi }_{01}+ϵ\mathrm{\Phi }_{12}+\kappa \mathrm{\Phi }_{22}\kappa ^{}\mathrm{\Phi }_{00}`$
$`=\mu \mathrm{\Phi }_{02}{\displaystyle \frac{1}{2}}\mu (\lambda +9\mathrm{\Lambda }+\text{æ}T_{02})`$ (46)
$`\delta \mathrm{\Phi }_{01}D\mathrm{\Phi }_{02}2\kappa \mathrm{\Phi }_{12}(\alpha +2\tau ^{})\mathrm{\Phi }_{01}+\sigma ^{}\mathrm{\Phi }_{00}+3\sigma \mathrm{\Phi }_{11}+{\displaystyle \frac{9}{4}}D\mathrm{\Lambda }`$
$`=\mu (\mathrm{\Phi }_{01}{\displaystyle \frac{\text{æ}}{2}}T_{01})`$
$`\delta \mathrm{\Phi }_{12}\mathrm{\Delta }\mathrm{\Phi }_{02}2\kappa ^{}\mathrm{\Phi }_{01}+(\alpha 2\tau )\mathrm{\Phi }_{12}+\sigma \mathrm{\Phi }_{22}+3\sigma ^{}\mathrm{\Phi }_{11}+{\displaystyle \frac{9}{4}}\mathrm{\Delta }\mathrm{\Lambda }`$ (47)
$`=\mu (\mathrm{\Phi }_{12}{\displaystyle \frac{\text{æ}}{2}}T_{12})`$
$`D\mathrm{\Phi }_{11}\mathrm{\Delta }\mathrm{\Phi }_{00}+(\tau ^{}2\tau )\mathrm{\Phi }_{01}2ϵ^{}\mathrm{\Phi }_{00}+\kappa \mathrm{\Phi }_{12}{\displaystyle \frac{3}{4}}D\mathrm{\Lambda }`$
$`=\mu (\mathrm{\Phi }_{01}{\displaystyle \frac{\text{æ}}{2}}T_{01})`$
$`\mathrm{\Delta }\mathrm{\Phi }_{11}D\mathrm{\Phi }_{22}+(\tau 2\tau ^{})\mathrm{\Phi }_{12}2ϵ\mathrm{\Phi }_{22}+\kappa ^{}\mathrm{\Phi }_{01}{\displaystyle \frac{3}{4}}\mathrm{\Delta }\mathrm{\Lambda }`$ (48)
$`=\mu (\mathrm{\Phi }_{12}{\displaystyle \frac{\text{æ}}{2}}T_{12})`$
$`\mathrm{\Delta }\mathrm{\Phi }_{01}\delta \mathrm{\Phi }_{11}+\kappa ^{}\mathrm{\Phi }_{00}+3\tau \mathrm{\Phi }_{11}\sigma \mathrm{\Phi }_{12}+(ϵ^{}\sigma ^{})\mathrm{\Phi }_{01}+{\displaystyle \frac{3}{4}}\delta \mathrm{\Lambda }`$
$`={\displaystyle \frac{1}{2}}\mu \left[(\lambda +\mathrm{\Phi }_{02})\text{æ}(T_{11}+{\displaystyle \frac{1}{3}}T)\right]`$
$`D\mathrm{\Phi }_{12}\delta \mathrm{\Phi }_{11}+\kappa \mathrm{\Phi }_{22}+3\tau ^{}\mathrm{\Phi }_{11}\sigma ^{}\mathrm{\Phi }_{01}+(ϵ\sigma )\mathrm{\Phi }_{12}+{\displaystyle \frac{3}{4}}\delta \mathrm{\Lambda }`$ (49)
$`={\displaystyle \frac{1}{2}}\mu \left[(\lambda +\mathrm{\Phi }_{02})\text{æ}(T_{11}+{\displaystyle \frac{1}{3}}T)\right]`$
$`D\mathrm{\Phi }_{01}\delta \mathrm{\Phi }_{00}+3\kappa \mathrm{\Phi }_{11}(ϵ+2\sigma )\mathrm{\Phi }_{01}+(\tau ^{}+2\alpha )\mathrm{\Phi }_{00}`$
$`=\mu (\mathrm{\Phi }_{00}{\displaystyle \frac{\text{æ}}{2}}T_{00})`$
$`\mathrm{\Delta }\mathrm{\Phi }_{12}\delta \mathrm{\Phi }_{22}+3\kappa ^{}\mathrm{\Phi }_{11}(ϵ^{}+2\sigma ^{})\mathrm{\Phi }_{12}+(\tau 2\alpha )\mathrm{\Phi }_{22}`$ (50)
$`=\mu (\mathrm{\Phi }_{22}{\displaystyle \frac{\text{æ}}{2}}T_{22}).`$
Finally the DJT field equations (44) imply the trace relation
$$\lambda +\frac{1}{6}R=\frac{1}{3}\text{æ}T$$
(51)
since the Cotton tensor is traceless. We note that the first equation in the above set remains invariant and the rest equations in each pair go into other under the prime operation provided that the sign of the DJT coupling constant is also changed simultaneously.
## 6 Exact solutions
We shall now extend the discussion of the topologically massive Dirac monopole given in section 2 by presenting exact solutions of the system of coupled DJT-MCS field equations that describe a self-gravitating magnetic monopole. For this purpose, we shall use the Newman-Penrose version of the TME (33) and TMG (46) -(50) field equations. We start with a homogeneous space which is given by the left-invariant 1-forms
$`\sigma ^0`$ $`=`$ $`d\psi +\mathrm{cosh}(b\theta )d\varphi `$
$`\sigma ^1`$ $`=`$ $`\mathrm{sin}(b\psi )d\theta +\mathrm{cos}(b\psi )\mathrm{sinh}(b\theta )d\varphi `$ (52)
$`\sigma ^2`$ $`=`$ $`\mathrm{cos}(b\psi )d\theta +\mathrm{sin}(b\psi )\mathrm{sinh}(b\theta )d\varphi `$
of modified Bianchi Type $`VIII`$. There will be no confusion as the earlier definition of $`\sigma ^i`$ in eqs.(10) will not be used in the rest of this paper. The $`1`$-forms (52) satisfy the Maurer-Cartan equations (11) with structure constants
$$C_{\mathrm{\hspace{0.33em}\hspace{0.33em}12}}^0=C_{\mathrm{\hspace{0.33em}\hspace{0.33em}21}}^0=b,C_{\mathrm{\hspace{0.33em}\hspace{0.33em}20}}^1=C_{\mathrm{\hspace{0.33em}\hspace{0.33em}02}}^1=b,C_{\mathrm{\hspace{0.33em}\hspace{0.33em}01}}^2=C_{\mathrm{\hspace{0.33em}\hspace{0.33em}10}}^2=b.$$
(53)
We define the co-frame
$$\omega ^0=\lambda _0\sigma ^0,\omega ^1=\lambda _1\sigma ^1,\omega ^2=\lambda _2\sigma ^2,$$
(54)
where $`\lambda _0,\lambda _1`$ and $`\lambda _2`$ are constant scale factors . The triad basis 1-forms will be defined by
$$l=\frac{1}{\sqrt{2}}\left(\omega ^0\omega ^1\right),n=\frac{1}{\sqrt{2}}\left(\omega ^0+\omega ^1\right),m=\omega ^2.$$
(55)
Then the Ricci rotation coefficients
$`ϵ`$ $`=`$ $`ϵ^{}=\sigma =\sigma ^{}=0,`$
$`\tau `$ $`=`$ $`\tau ^{}={\displaystyle \frac{b\lambda _2}{2\lambda _0\lambda _1}},`$
$`\kappa `$ $`=`$ $`\kappa ^{}={\displaystyle \frac{b}{2\lambda _0\lambda _1\lambda _2}}\left(\lambda _0^2\lambda _1^2\right),`$ (56)
$`\alpha `$ $`=`$ $`{\displaystyle \frac{b}{2\lambda _0\lambda _1\lambda _2}}\left(\lambda _0^2+\lambda _1^2\lambda _2^2\right)`$
are obtained by taking the exterior derivative of eqs.(55) and comparing the result with eqs.(23). The Ricci identities (I. 45-49) now yield the expression for the scalar of curvature
$$R=\frac{b^2}{2\lambda _0^2\lambda _1^2\lambda _2^2}(\lambda _0+\lambda _1+\lambda _2)(\lambda _1+\lambda _2\lambda _0)(\lambda _0\lambda _1+\lambda _2)(\lambda _2\lambda _0\lambda _1)$$
(57)
which holds for both Bianchi Types $`VIII`$ and $`IX`$. It is related to Menger curvature $`K`$ by
$$R\frac{b^2}{2}K^2$$
(58)
for three points on a space curve in a fictitious flat $`3`$-dimensional Euclidean space where $`\lambda _0,\lambda _1`$ and $`\lambda _2`$ denote the distances between these points . If we consider the limit $`\lambda _i0`$ keeping one point fixed, then Menger curvature reduces to the definition of curvature in the Serre-Frenet formulas. This identification offers the possibility of classifying homogeneous solutions of TMG in terms of of space curves. Namely, for vacuum solutions where $`\lambda _0=\lambda _1+\lambda _2`$ the space curve is a straight line.
Now turning our attention to the potential 1-form, we note that we can take
$$A=g\sigma ^0$$
(59)
as in eq.(13) but on modified Bianchi Type $`VIII`$ left-invariant $`1`$-forms (52). Then the field $`2`$-form is given by
$$F=\frac{1}{\sqrt{2}}\frac{gb}{\lambda _1\lambda _2}(lmnm)$$
(60)
and we get
$$\phi _0=\phi _2=\frac{1}{\sqrt{2}}\frac{gb}{\lambda _1\lambda _2},\phi _1=0$$
(61)
for the electromagnetric triad scalars.
Starting with the TMG frame (6), (54) and TME potential (59) all triad scalars reduce to constants and the field equations of TME (33) and TMG (46) -(50) consist of a set of polynomials for constant scale factors $`\lambda _i`$ in terms of constants in the theory, namely the topological masses $`\mu ,\nu `$ and the cosmological constant æ.
The TME eqs.(33) are satisfied identically provided that the deficit angle in eqs.(6) is determined by
$$b=\nu \frac{\lambda _1\lambda _2}{\lambda _0}$$
(62)
in terms of the electromagnetic topological mass. Next, we consider the triad components of the Maxwell stress tensor using eqs.(36) and (43) along with eq.(62). For the nonvanishing components of the energy-momentum tensor we have
$$T_{00}=T_{02}=T_{22}=3T_{11}=T=\frac{g^2\nu ^2}{2\lambda _0^2}$$
(63)
which enter in the right-hand-side of the DJT field equations (46) -(50).
### 6.1 Two equal scale factors
First we shall consider the case of vanishing cosmological constant $`\lambda =0`$ and show that it forces the equality $`\lambda _1=\lambda _2`$ between scale factors. The resulting solution is a generalization of the Vuorio solution . Writing the trace equation (51) in terms of the Ricci rotation coefficients we have
$$2\alpha \tau +\kappa ^2\tau ^2=\text{æ}T$$
(64)
or using eqs.(6) we write it in the explicit form
$$\frac{1}{4}b^2K^2=\text{æ}T$$
(65)
and once again see that $`K=0`$ leads to vacuum solutions which is not of interest in this paper. Henceforth we shall take $`K0`$. The remaining DJT field equations reduce to
$$3\kappa \mathrm{\Phi }_{11}+(2\alpha \tau )\mathrm{\Phi }_{00}=\mu \mathrm{\Phi }_{00}+\frac{1}{2}\text{æ}T$$
(66)
$$3\tau ^{}\mathrm{\Phi }_{11}+\kappa \mathrm{\Phi }_{22}=\frac{1}{2}\mu \mathrm{\Phi }_{02},$$
(67)
and all other DJT field equations are identically satisfied. The explicit form of these equations is obtained by the substitution of the rotation coefficients (6) and sources (63). When we compare the resulting expressions eqs.(66) and (67) to eq.(64) we arrive at a polynomial constraint
$$\frac{b^3}{\lambda _0\lambda _1\lambda _2}\left(\frac{\lambda _1^2+\lambda _2^2\lambda _0^2}{\lambda _2^2\lambda _1^23\lambda _0^2}\right)(\lambda _2^2\lambda _1^2)K^2=0$$
(68)
which involves only the scale factors $`\lambda _i`$. We emphasize that this constraint holds only in the case of vanishing cosmological constant. As it is readily seen from eq.(65) only the roots $`\lambda _1=\pm \lambda _2`$ give rise to non-vacuum solutions of the DJT field eqations. For the sake of certainty, we shall take $`\lambda _1=\lambda _2`$, then the simultaneous solution of eqs.(65)- (67) has the form
$$\lambda _0^2=2\text{æ}g^2\frac{\nu +2\mu }{2\mu +3\nu }\lambda _1^2=\lambda _2^2=2\text{æ}g^2\frac{\mu +\nu }{2\mu +3\nu }$$
(69)
which is the generalization of the Vuorio solution for a TME-TMG magnetic monopole. Indeed, in the limiting case $`g0`$ when the monopole charge vanishes, the denominator in eqs.(69) vanishes as well, so that the ratio
$$\frac{2\text{æ}g^2}{2\mu +3\nu }=\lambda _1^3=\frac{\lambda _0^3}{8}$$
(70)
is constant and the solution (69) reduces to the Vuorio solution.
Finally, we note that the above solution of the DJT field equations with a topologically massive monopole can be readily generalized to include a cosmological constant. We find
$$\lambda _0^2=\frac{2\text{æ}g^2\nu ^2}{2\mu +3\nu }\frac{\nu +2\mu }{\nu ^2+4\lambda }\lambda _1^2=\lambda _2^2=\frac{2\text{æ}g^2}{2\mu +3\nu }\frac{\nu ^3+\mu \nu ^24\lambda \mu }{\nu ^2+4\lambda }$$
(71)
and the final metric is given by
$`ds^2`$ $`=`$ $`{\displaystyle \frac{2\text{æ}g^2}{(2\mu +3\nu )(\nu ^2+4\lambda )}}\{\nu ^2(2\mu +\nu )\left(\sigma ^0\right)^2`$ (72)
$`(\nu ^3+\mu \nu ^24\lambda \mu )[\left(\sigma ^1\right)^2+\left(\sigma ^2\right)^2]\}`$
which generalizes the solution to the case of a topologically massive magnetic monopole.
### 6.2 Tri-axial solution
When all three scale factors are different DJT-MCS field equations admit a solution which is obtained by considerations very similar to those given above. However, in this case polynomials involving the scale factors are much more complicated. For the tri-axial solution eq.(62) relating the deficit in the angle $`\theta `$ to TME mass $`\nu `$ remains the same. But the result (16) is now modified to
$$\frac{\lambda }{\nu ^2/4}=(\lambda _0^2\lambda _1^2\lambda _2^2)^2\frac{\lambda _1^{\mathrm{\hspace{0.33em}2}}\lambda _2^{\mathrm{\hspace{0.33em}2}}K^2}{\lambda _0^{\mathrm{\hspace{0.33em}2}}A}$$
(73)
where
$$A\lambda _0^{\mathrm{\hspace{0.33em}4}}3\lambda _1^{\mathrm{\hspace{0.33em}4}}3\lambda _2^{\mathrm{\hspace{0.33em}4}}+2\lambda _0^{\mathrm{\hspace{0.33em}2}}\lambda _1^{\mathrm{\hspace{0.33em}2}}+2\lambda _0^{\mathrm{\hspace{0.33em}2}}\lambda _2^{\mathrm{\hspace{0.33em}2}}2\lambda _1^{\mathrm{\hspace{0.33em}2}}\lambda _2^{\mathrm{\hspace{0.33em}2}}$$
(74)
is another interesting polynomial. The ratio of the two topological masses is given by
$$\frac{\mu }{\nu }=\frac{A}{2\lambda _0^{\mathrm{\hspace{0.33em}2}}(\lambda _0^{\mathrm{\hspace{0.33em}2}}\lambda _1^{\mathrm{\hspace{0.33em}2}}\lambda _2^{\mathrm{\hspace{0.33em}2}})}$$
(75)
and finally
$$g^2\text{æ}=\frac{4}{A}(\lambda _0^{\mathrm{\hspace{0.33em}2}}\lambda _1^{\mathrm{\hspace{0.33em}2}})(\lambda _0^{\mathrm{\hspace{0.33em}2}}\lambda _2^{\mathrm{\hspace{0.33em}2}})\lambda _1^{\mathrm{\hspace{0.33em}2}}\lambda _2^{\mathrm{\hspace{0.33em}2}}K^2$$
(76)
gives the relationship between the matter coupling constant and magnetic charge to the scale factors in the metric.
## 7 Conclusion
We have presented the exact solution for a self-gravitating magnetic monopole in topologically massive gravity and electrodynamics. Topological mass has the effect of turning the Dirac string into a cone as well as imparting conical deficit to the homogeneous Bianchi Type $`VIII`$ space. Just as in the case of topologically massive electrodynamical monopole without gravity, we have a physical system which requires curved spacetime for its existence. Furthermore, the conical deficit due to the topologically massive field can be accomodated only in curved spacetimes with a matching conical deficit. This is evidently a general phenomenon which we must expect when we consider monopole-type solutions of coupled MCS and DJT field equations. It is interesting to note that for pure Einstein gravity in $`3`$-dimensions coupled to Maxwell-Chern-Simons field all three scale factors must coincide and there exists no self-gravitating monopole solution.
## 8 Acknowledgement
We thank Stanley Deser for many interesting discussions.
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# Dynamics of multiply charged ions in intense laser fields
## I INTRODUCTION
Various techniques to generate ultrashort pulses such as chirped-pulse-amplification (CPA) have been developed and perfected over the last years such that nowadays powerful laser pulses are available over a wide range of frequencies and pulse lengths up to intensities of $`10^{21}W/cm^2`$ . For intensities of $`10^{16}W/cm^2`$ to $`10^{19}W/cm^2`$ the electric field strength of the laser pulse is already comparable up to far stronger than the atomic unit field strength ($`5.14\times 10^9V/cm`$) that is experienced by an electron on the first Bohr orbit of hydrogen and in fact those fields are already accessible in quite numerous laboratories worldwide in rather small table-top set-ups.
Especially for laser intensities till at about $`10^{16}W/cm^2`$ there has been a lot of activity over the last two decades in intense nonrelativistic laser interactions with matter including atoms , molecules , clusters and solids . Many highly nonlinear optical phenomena such as high-order harmonic generation , above-threshold ionization (ATI) and for higher frequencies and intensities stabilization , have attracted much attention for both experimentalists and theorists. Since optical laser fields of intensity $`10^{21}W/cm^2`$ have become available in a rather short time, both theoretical and experimental activities moved quite quickly into the regime of fully relativistic dynamics up to around $`10^{18}W/cm^2`$ to $`10^{22}W/cm^2`$ with experimental emphasis on free electron dynamics, QED effects and nuclear reaction processes . There has been rather little activity however for neutral atoms interacting with optical laser fields of intensity $`10^{16}W/cm^2`$ to $`10^{17}W/cm^2`$ merely for the reason, that there is almost instantaneous ionization for those intensities and the free electron dynamics is quite well understood for a long time . For higher frequencies ionization may be retained at those intensities however will also occur eventually for larger laser intensities due to the magnetic drift with a likelyhood for magnetic recollisions being usually smaller than that in the laser polarisation direction .
For an ionic system one may have the unique possibility to apply relativistic near-optical laser field pulses and still allow for bound dynamics. The physical processes occuring will not merely scale with the ionic charge but be fundamentally different because of relativity and later QED governing the dynamics. Nowadays ions may be processed from essentially all existing atoms with abitrary charge state, absolute purity and quite high density by sending them through foils while due to lasers high charge states have also been achieved as well however with limits in the width of the charge distribution and the absolute charge . The mean electric field sensed by a ground state electron of hydrogenic uranium corresponds to $`1.8\times 10^{16}V/cm`$. Thus laser fields that are comparable in strength need have an intensity of order $`10^{29}W/cm^2`$ and are thus far beyond reach nowadays. For those intensities the dynamics will be absolutely remote from that of hydrogen with say a laser field of $`10^{14}W/cm^2`$ and we do not aim to make the corresponding comparison in this paper. We consider it however more interesting to carry out a comparision of the nonrelativistic laser atom interaction with the weakly relativistic dynamics of an ion of charge around 10 with laser fields of intensity $`10^{16}W/cm^2`$ to $`10^{17}W/cm^2`$ (see figure 1). Here relativity induces merely corrections and it is possible to identify the leading deviations to nonrelativistic dynamics as that of the laser magnetic field component, the break-down of the dipole approximation, spin-orbit coupling, zitterbewegung amd the relativistic mass shift. In earlier letters we have shown that the magnetic field component of the laser pulse induces an enhanced angular momentum and thus a reduced electron expectation value in the vicinity of the multiply charged ion and the existence of spin signatures in bound electron dynamics and radiation .
In terms of applications multiply charged ions in relativistic laser fields appear also attractive for the generation of energic electrons and high harmonics. The kinetic energy of a laser-accelerated electron as characterised by the parameter $`U_p`$ does certainly increase with rising laser intensity. For harmonic generation especially it is necessary that the highly energetic electron interacts periodically with the nucleus, urging us to enhance the charge of the ionic core, i.e. the ionization potential $`I_p`$ correspondingly. As long as tunneling and recollisions can be assured with an appropriate relation of $`U_p`$ and $`I_p`$ the simultaneous increase of both parameters is obviously attractive as the cut-off energy of high harmonic generation is given by the sum of $`I_p`$ and the maximal kinetic energy at the time of recollision being of order $`U_p`$.
In this paper, we study the dynamics of multiply charged ions in intense lasers fields in the weakly relativistic parameter regime where terms up to $`1/c^2`$ are still of significance and higher terms become negligible. For this purpose we solve numerically the two-dimensional time-dependent expansion of the Dirac equation which is exact up to all orders in $`1/c^2`$ and which was first derived by Foldy and Wouthuysen . Different weakly relativistic effects arise and we can investigate each by solving the dynamics and comparing the situations in which we include or neglect the corresponding part of the Hamiltoninan associated with the effect of interest. As a function of the ionic charge relative to the laser field intensity the electric and magnetic field components of the laser field will be shown to be either strong enough to lead to ionization in the polarization and propagation directions or give rise to interesting structures in the near vicinity of the ionic core. We point out the role of the spin in inducing stronger binding of the electronic wavepacket and multiphoton spectral line splitting due to spin orbit coupling. This is compared with the situation via the Pauli equation where spin induced forces are neglected. We further show that the relativistic mass shift induces a significant shift of especially the highly excited eigenenergies of the ion with respect to the conventional Stark shift and point out the consequences for the spectral lines. For relatively high intensities special emphasis is placed on the tunneling regime in the weakly relativistic regime. For multiply charged ions we find that both tunneling and recollisions are still possible and indicate harmonic generation in the keV regime. We note however that the plateau of the harmonic spectrum, as well known in the nonrelativistic regime, is tilted increasingly in the relativistic regime with rising charge. This deviation to conventional high harmonic spectra is not too surprising as an electric wavepacket attempting to recollide due to a phase shift of the electric field in the laser polarization direction may still partly or totally miss the ionic core due the magnetic drift in the laser propagation direction, especially for long recollision times. In the above threshold spectra we find high electron energies in the keV regime which are separated by the photon energy of the applied laser field.
This paper is arranged as follows: In section II, we derive the Hamiltonian of interest along the lines of Foldy and Wouthuysen, then describe the numerical methods for solving the corresponding dynamical equation and present the details for computing the observable quantities of interest. In subsection III(A) we then present effects up to first order in $`1/c`$ as those induced by the magnetic field component of the laser field. In subsection III(B) we investigate the relativistic corrections to the Stark shift for laser-driven ions followed by a study of spin induced forces and the consequent splitting of spectral features in III(C). Subsequently, relativistic high order harmonic generation is discussed in III(D), as well as the photoelectron spectra in subsection III(E). Finally, a conclusion is drawn.
## II Preliminary Considerations
In this more technical section we present subsequently the Hamiltonian describing our system of interest, explain how we solve the corresponding dynamical equation numerically via the split-step mechanism and finally quantify the observables to be investigated in the following section.
### A System and Hamiltonian
We are interested throughout this article in the dynamics of multiply charged ions of charge state up to $`Z=12`$ with moderate laser intensities ($`10^{16}10^{17}W/cm^2`$) and near optical frequencies for the KrF (248nm; 0.1838 a.u.) laser system and the doubled Nd:glass laser frequency (527nm; 0.0866 a.u.). The maximum velocity of electron wavepackets reaches the order $`v=0.1c`$, where $`c=137.036`$ denotes the speed of light in atomic units. We are therefore entitled to consider the Dirac equation up to second order in $`v/c`$ and have confirmed that for our parameters high order corrections do not play a role. One may derive this Hamiltonian via various unitary transformations along the lines of the Foldy-Wouthuysen (FW) expansion of the Dirac equation or alternatively, with the same result, find the first order relativistic corrections to the Pauli equation . As opposed to nonrelativistic treatments we need to include at least two dimensions in the calculation as the magnetic field component of the laser pulse may induce a significant drift in the laser propagation direction. There is a spin-induced force in the magnetic field direction , i.e. in the remaining third dimension, but the influence is small for the observables and the parameters of interest here.
Our working Hamiltonian $`𝐇_{\mathrm{𝐅𝐖}}`$ involves a series of relativistic corrections to the usual nonrelativistic Schrödinger Hamiltonian. One well-known correction term in $`v/c`$ to the Schrödinger equation is the additional term in the Pauli equation representing the coupling of the laser magnetic field to the spin degree of freedom of the electron wavepacket. The second order terms include the spin orbit coupling, i.e. the feed-back of the oscillating spin to the electron motion, as well as the leading relativistic mass shift term and Zitterbewegung. Terms of order $`O(1/c^3)`$ do not play a role for the parameters employed though would certainly be necessary to include for the fully relativistic regime involving more intense laser fields and higher charged ions. The main advantage of using this equation in comparison to the full Dirac equation is the possible isolation of the influence of each physical mechanism arising. In addition this equation does not limit us numerically to use high frequency lasers as necessary so far with the full Dirac equation .
For the circumstances with laser parameters described above, the FW Hamiltonian ($`2\times 2`$ matrix) of a bound electron in a strong laser field can be written (in atomic units ) as
$$\begin{array}{ccc}𝐇_{\mathrm{𝐅𝐖}}\hfill & =& 𝐇_\mathrm{𝟎}+𝐇_𝐩+𝐇_{\mathrm{𝐤𝐢𝐧}}+𝐇_𝐃+𝐇_{\mathrm{𝐬𝐨}}\hfill \\ 𝐇_\mathrm{𝟎}\hfill & =& (𝐩+𝐀(z,t)/c)^2/2+V(x,z)\hfill \\ 𝐇_𝐏\hfill & =& \sigma 𝐁(z,t)/2c\hfill \\ 𝐇_{\mathrm{𝐤𝐢𝐧}}\hfill & =& 𝐩^4/8c^2\hfill \\ 𝐇_𝐃\hfill & =& div𝐄^{}(x,z,t)/8c^2\hfill \\ 𝐇_{\mathrm{𝐬𝐨}}\hfill & =& i\sigma curl𝐄^{}/8c^2+\sigma 𝐄^{}\times 𝐩/4c^2.\hfill \end{array}$$
(1)
Here $`𝐇_\mathrm{𝟎}`$ denotes the standard nonrelativistic Hamiltonian in Schrödinger form, where $`𝐩=(p_x,0,p_z)=(i/x,0,i/z)`$ is the two-dimensional momentum operator and $`𝐀(z,t)`$ is the time-spatial dependent vector potential of the laser field $`𝐄(z,t)`$, which is linearly polarized along the $`x`$-axis and propagates in $`z`$-direction. For the vector potential we include the magnetic field component and do not apply the dipole approximation, urging us to perform a two-dimensional numerical integration in the $`xz`$ plane. Since $`𝐇_\mathrm{𝟎}`$ all other Hamiltonians are $`2\times 2`$ matrices, they need be considered as multiplied by the unity matrix I even if not shown explicitely throughout this paper. We consider multiply charged ions in the single active electron approximation which are preionized by several or more than ten electrons and thus are easily availabe today via lasers or with highest accuracy via shooting the atoms through thin foiles . Those are well described by the soft-core potential to model the Coulomb field experienced by the active electron of a multiply charged ion, i.e.
$$V(x,z)=k/\sqrt{q_e+x^2+z^2}.$$
(2)
The parameter $`k`$ is a function of the effective number of positive charges $`Z`$ as sensed by the electron. $`q_e`$ compensates for the effect of possible inner electrons and reduced distances of the electronic wavepacket to the ionic core in two- rather three-dimensional calculations. $`k`$ may be adapted such that we obtain the correct ionization energy for the system of interest with effective charge of the ionic core $`Z`$ and charge of the ion $`Z1`$. The static field of the ionic core is expressed by the gradient of the potential $`V(x,z)`$ and $`𝐄^{}(x,z,t)`$ stands for the sum of this field plus the laser field $`𝐄(z,t)`$. The following term $`𝐇_𝐏`$ in Eq.(1) indicates the coupling of the laser magnetic field $`𝐁`$ to the electronic spin as described by the Pauli matrix $`\sigma `$. The sum $`𝐇_\mathrm{𝟎}+𝐇_𝐏`$ leads to the Hamiltonian in the well known Pauli equation. Further in Eq. (1) $`𝐇_{\mathrm{𝐤𝐢𝐧}}`$ denotes the leading term for the relativistic mass increase, and $`𝐇_𝐃`$ is the well-known Darwin term. Finally the last term in the Hamiltonian $`𝐇_{\mathrm{𝐬𝐨}}`$ stands for the spin-orbit coupling. Considering our central potential $`V(x,z)`$ the first term of $`𝐇_{\mathrm{𝐬𝐨}}`$ in Eq. (1) disappears because $`\times (V(x,z))=0`$ and the contribution due to the laser field is of order $`1/c^3`$. Thus, the spin-orbit coupling term becomes
$$𝐇_{\mathrm{𝐬𝐨}}=\sigma 𝐄^{}\times 𝐩/4c^2=\sigma 𝐄\times 𝐩/4c^2+f(x,z)\sigma 𝐋$$
(3)
with $`f(x,z)=k(q_e+x^2+z^2)^{3/2}/4c^2`$ and where $`𝐋=𝐫\times 𝐩=(0,zp_xxp_z,0)`$ is the orbital angular momentum, of which only the component along the y-direction is non-zero. The origin of spin-orbit coupling can alternatively be viewed also as being due to the interaction between the magnetic moment of the electron and the magnetic field $`𝐁^{}`$ due to the motion of the positively charged core as sensed by the electron in its own rest frame.
### B Dynamics and Numerical Approach
We investigate the dynamics of multiply charged ions exposed to an intense laser field through solving the following dynamical equation, involving the previously derived Hamiltonian $`𝐇_{\mathrm{𝐅𝐖}}`$ in Eq. (1) (for the convenience, we use the usual atomic units throughout this paper ):
$$i\frac{}{t}\left(\begin{array}{c}\mathrm{\Psi }_{up}(x,z,t)\\ \mathrm{\Psi }_{down}(x,z,t)\end{array}\right)=𝐇_{\mathrm{𝐅𝐖}}\left(\begin{array}{c}\mathrm{\Psi }_{up}(x,z,t)\\ \mathrm{\Psi }_{down}(x,z,t)\end{array}\right).$$
(4)
The wave function has two components corresponding to spin-up and spin-down polarization of the electron and $`𝐇_{\mathrm{𝐅𝐖}}`$ is consequently a $`2\times 2`$ matrix operator. The coupling to negative energy states as included in conventional Dirac theory is negligible in second order in $`v/c`$. The laser field is assumed to be linearly polarized along the $`x`$-axis so that the vector potential $`𝐀(z,t)`$ of the laser field may be of the form
$$𝐀(z,t)=(A_x(z,t),0,0)$$
(5)
where the $`z`$-dependence of the vector potential reflects the propagation of the laser pulse in the $`z`$ direction. The spatial dependence of the vector potential indicates that the magnetic component of the laser field $`B=\times 𝐀(z,t)/c0`$ is included and we do not (and can not) carry out the dipole approximation. We choose the vector potential $`A_x(z,t)`$ to be
$$A_x(z,t)=\{\begin{array}{cc}\frac{cE_0}{\omega t_{on}}\left[(tz/c)sin(\omega t\omega z/c)+\frac{1}{\omega }cos(\omega t\omega z/c)\right]\hfill & \hfill 0<tz/ct_{on}\\ \frac{cE_0}{\omega }sin(\omega t\omega z/c)\hfill & \hfill t_{on}<tz/c<t_p\end{array}$$
(6)
which is associated with a linearly polarized laser field with electric field $`E_x`$ and magnetic field $`B_y`$ components
$$E_x(z,t)=B_y(z,t)=\{\begin{array}{cc}E_0\frac{tz/c}{t_{on}}cos(\omega t\omega z/c)\hfill & \hfill 0<tz/ct_{on}\\ E_0cos(\omega t\omega z/c)\hfill & \hfill t_{on}<tz/c<t_p\end{array}$$
(7)
being oriented in the $`x`$ and $`y`$ direction, respectively and propagating both in phase in the $`z`$ direction. Here, $`E_0`$ and $`\omega `$ are the maximal amplitudes of both fields and the angular frequency of the laser field, respectively. Further $`t_{on}`$ is associated with the linear rising time of the laser pulse, i.e. 0 is the beginning of the turn-on and $`t_{on}`$ the end of it. We note that the vector potential Eq. (6) may not be continuous at the end of the turn-on phase, however it is when $`t_{on}`$ is chosen such that $`\omega (t_{on}z/c)=(m+0.25)2\pi `$ with m being an arbitrary integer. The measurable electric field strength however is always continuous. After the turn-on phase the pulse is assumed to have a constant amplitude till time $`t_p`$. Obviously a realistic pulse will turn off afterwards smoothly, however for all observables of interest here this phase is of no interest and numerical calculations usually terminate at $`t_p`$.
Since the laser field is linearly polarized, the interaction term involves a term of the form $`𝐩𝐀(𝐳,𝐭)/\mathrm{c}=\mathrm{p}_\mathrm{x}\mathrm{A}_\mathrm{x}(\mathrm{z},\mathrm{t})/\mathrm{c}`$. which means no coupling between momentum and coordinate space. This is because the term couples only the $`x`$ component of the momentum with a function which is dependent on $`z`$ but not on $`x`$. Thus, we can still apply the usual split-operator algorithm to solve the two-dimensional time-dependent matrix equation (4) via
$$\begin{array}{ccc}\hfill \left(\begin{array}{c}\mathrm{\Psi }_{up}(x,z,t+\mathrm{\Delta }t)\\ \mathrm{\Psi }_{down}(x,z,t+\mathrm{\Delta }t)\end{array}\right)& =& exp[i\mathrm{\Delta }t(p^2/4p^4/(16c^2))𝐈]\hfill \\ & & \times exp[i\mathrm{\Delta }t((𝐩𝐀/c+A^2/c^2)𝐈+𝐇_𝐃)]\times exp[i\mathrm{\Delta }t(𝐇_𝐏+𝐇_{\mathrm{𝐬𝐨}}]\hfill \\ & & \times exp[i\mathrm{\Delta }t(p^2/4p^4/(16c^2))𝐈]\times \left(\begin{array}{c}\mathrm{\Psi }_{up}(x,z,t)\\ \mathrm{\Psi }_{down}(x,z,t)\end{array}\right).\hfill \end{array}$$
(8)
Here, $`\mathrm{\Delta }t`$ denotes the time step and the unit matrix operator is described with I. All exponential operators except the term $`exp[i\mathrm{\Delta }t(𝐇_𝐏+𝐇_{\mathrm{𝐬𝐨}}]`$ are diagonal and we apply the split evolution operator on the wave function with the help of Fourier transforming between the coordinate representation and the momentum representation. Consequently all derivative operator can be transformed into multiplications with constants. For non-diagonal exponential operators, we usually need to diagonize them. Fortunately, the non-diagonal operator $`exp[i\mathrm{\Delta }t(H_P+H_{so}]`$ involved here only depends on the specific Pauli matrix $`\sigma `$, so that we can carry out the Taylor expansion for this exponential operator up to the order of interest here. Regarding the interaction term $`p_xA_x/c`$ special care is needed as mentionned above. Here we do the Fourier transformation only for the x-coordinate because of the z-coordinate dependence of the vector potential $`A_x(z,t)`$. This however is sufficient and the exponential function of operators of interaction terms on the wave function ends up being merely a sequence of Fourier transformations and c-number multiplications. Because of the splitting of the total Hamiltonian in the exponent we introduce an error following the Baker-Hausdorff formula because the split terms do generally not commute. The error of this algorithm is of order $`(\mathrm{\Delta }t)^3`$ between every successive time step. Thus, a small time step can ensure to get accurate results; and we have not experienced any problem with numerical convergence in the regime of interest here.
From the numerical point of view we first solve for the eigenstates of the bound electron in the ionic core potential by using the so-called ”spectral” method . In fact, we choose a testing wave function without any symmetries to propagate on the two-dimensional potential. This way we obtain all symmetric and asymmetric eigenstates of the system. As an example, figure 2 shows the energy-level structure of a model hydrogen-like ion $`Mg^{11+}`$ of which the single active electron experiences the nucleus with charges Z=12. Choosing k=80.32 and $`q_e=1.0`$ we get the corresponding ground-state energy -72au. Note that it is not our purpose to present an exact quantitative model of the ionic level structure of an ionic system. We model single electron ions with a softcore rather than a Coulomb potential and can adapt the ground state energy by choosing $`k`$ and $`q_e`$. With the assumption of a smooting around the ionic core, we deviate from the exact Coulombic potential and the system may also be considered as an active electron plus an ionic core where lower shells are filled and inactive. In the intensity range where level structures are important we thus can only make qualitative statements for the dynamics if we want to associate it with a particular realistic ion. In the tunneling regime however where structure becomes less important and only the correct ground state is significant we can be quantitative as well.
The ground-state wave packet of our model ion relates to a symmetric s-state wavepacket, and also the second excited-state wave packet (see figure 2b). The first, third and fourth excited-states are asymmetric (especially not s-states) in space. This system will be applied to evaluate the relativistic Stark shift in section III(B). After obtaining the eigenstates $`\mathrm{\Phi }_n(x,z)`$, we may use the above split-step operator in Eq.(8) to investigate the evolution of multiply charged ions under the irradiation of external intense lasers. When the laser pulse is turned on, the system is assumed to evolve from the spin-polarized ground state, that is, the initial condition $`\mathrm{\Psi }_{up}(x,z,t=0)=\mathrm{\Phi }_1(x,z)`$ and $`\mathrm{\Psi }_{down}(x,z,t=0)=0`$.
### C Observables of interest
With the knowledge of the time-dependent wavefunction we are in the position to calculate the spatial probability distribution $`|\mathrm{\Psi }(x,z,t)|^2=|\mathrm{\Psi }_{up}(x,z,t)|^2+|\mathrm{\Psi }_{down}(x,z,t)|^2`$ and the expectation value of any observable $`𝒪`$ associated with the system via
$$𝒪(t)=_{x_{min}}^{x_{max}}dx_{z_{min}}^{z_{max}}dz(\mathrm{\Psi }_{up}(x,z,t),\mathrm{\Psi }_{down}(x,z,t))𝒪(x,z,t)\left(\begin{array}{c}\mathrm{\Psi }_{up}(x,z,t)\\ \mathrm{\Psi }_{down}(x,z,t)\end{array}\right).$$
(9)
Here the integral turns into a sum in our case as the wavefunction is given on an equidistant grid from $`x_{min}`$ to $`x_{max}`$ in the polarization direction and from $`z_{min}`$ to $`z_{max}`$ in the propagation direction.
Particular interest in this article is placed on the spatial average values $`x(t)𝐈`$ and $`z(t)𝐈`$ of the wavepacket and because of their relevance for the radiation spectrum on the accelerations in the polarization direction $`a_x(t)=\stackrel{..}{x}(t)𝐈`$ and in the propagation direction $`a_z(t)=\stackrel{..}{z}(t)𝐈`$ via Eq. (9).
The radiation spectrum is generally given by a rather complex function of the accelerations and velocities in all spatial directions (see e.g. Eq. 10.7 in the first review in ). For simplicity we here restrict ourselves to an observation direction perpendicular to the 2d plane of motion, i.e. the $`y`$ direction. Also we are mostly interested in the dominating part of the radiation spectrum in the weakly relativistic regime which is polarized in the polarization direction of the laser field. In the far field spectrum this is proportional to the squared Fourier transform of simply the acceleration $`a_x(t)=\stackrel{..}{x}(t)𝐈`$. Also for the purpose of studying weakly relativistic signatures of the spectra it is interessing to study the less intense spectrum which is polarised in the laser propagation direction and governed by $`a_z(t)=\stackrel{..}{z}(t)𝐈`$.
Under the assumption of including relativistic corrections up to second order in $`v/c`$ it is not sufficient for $`a_x(t)`$ to consider the gradient of the potential only as in the nonrelativistic regime. We rather find in the weakly relativistic limit
$$\begin{array}{ccc}\stackrel{..}{x}& =& \sqrt{1\stackrel{2}{\stackrel{.}{x}}/c^2\stackrel{2}{\stackrel{.}{z}}/c^2}\times \left(\frac{V(x,z)}{x}+\frac{\stackrel{.}{x}}{c}\left(\frac{\stackrel{.}{x}}{c}\frac{V(x,z)}{x}+\frac{\stackrel{.}{z}}{c}\frac{V(x,z)}{z}\right)\right)\hfill \\ & =& \left(1+\frac{3}{2c^2}\frac{^2}{x^2}+\frac{3}{2c^2}\frac{^2}{z^2}+O(1/c^4)\right)\times \left(\frac{V(x,z)}{x}\right)\hfill \\ & & \\ \stackrel{..}{z}& =& \left(1+\frac{3}{2c^2}\frac{^2}{x^2}+\frac{3}{2c^2}\frac{^2}{z^2}+O(1/c^4)\right)\times \left(\frac{V(x,z)}{z}\right)\hfill \end{array}$$
(10)
where for the transformation to the second part of the equation we substituted the time derivatives via $`\mathrm{d}x/dt=dH_{FW}/dp_x`$ and $`\mathrm{d}z/dt=dH_{FW}/dp_z`$ and will neglect the higher order terms in $`O(1/c^4)`$. We note that the first term in the right-hand side is just the non-relativistic limit, and the following two terms correspond to weakly relativistic corrections which will be retained for the calculation of spectra. For the evaluation of the spectrum of interest the operators in Eq. (10) need just be inserted in Eq. (9) and then to be Fourier transformed.
We now turn to the technical aspects for the evaluation of the photoelectron spectra. Since the expected ionization is equivalent to a proportion of the electronic wave function leaving the vicinity of the ionic core and propagating outwards towards the (unphysical) boundaries of the numerical grid, we must avoid reflections of the wave function at those boundaries. This is achieved by absorbing all parts of the wave function approaching the boundaries by a $`cos^{\frac{1}{8}}`$ mask function , which results in a decrease of the norm of the wave function within the spatial box. While quite often only the part in the vicinity of the ionic core interests, it is here the opposite as we care only for the part of the wavefunction approaching the absorbing boundaries. Assume that we have calculated the wavepacket dynamics up to time $`t_f`$ in our finite box. Then we will have information about the photoelectron spectrum from what has been absorped till this time at the boundaries at all intermediate times $`t_\alpha `$ and from the ionized part of the wavefunction still in the box. We begin by treating the absorbed electron flux, say $`\mathrm{\Psi }_{flux}(x,z,t_\alpha )`$ at time $`t_\alpha `$ after entering the area of the boundaries. As this part will certainly not be influenced by the ionic core we will assume this to propagate to the end of interaction as a free particle (for details see ).
At the end of the interaction, at $`t_f`$, we also have to include the ionized part of the final spatial wave function still in the box. In order to obtain this from the full wave function it will be modified as follows, for example for the spin-up part of the wave function
$$\mathrm{\Psi }_{out}(x,z,t_f)=\{\begin{array}{cc}0\hfill & \left|x\right|<X_I\\ sin^2(\pi (\left|x\right|X_I)/2X_0)\mathrm{\Psi }_{up}(x,z,t_f)\hfill & X_I\left|x\right|X_I+X_0\\ \mathrm{\Psi }_{up}(x,z,t_f)\hfill & \left|x\right|>X_I+X_0\end{array}$$
(11)
Here $`X_I`$ represents a range which may be related to the ionic radius within which the electron may be considered bound and we will ignore this part of the wavepacket as we are only interested in the ionized part of the wave packet. We call $`X_0`$ the sliding range for the final wavefunction, where with increasing distance from the ionic core the wavefunction will contribute more to the photoelectron spectrum. Everything beyond $`X_I+X_0`$ can be considered as fully ionized and will fully contribute. The momentum wave function of the ionized electron can then be obtained by fast Fourier transforming both the spatial wave function $`\mathrm{\Psi }_{out}(x,z,t_f)`$ and the freely propagated $`\mathrm{\Psi }_{flux}(x,z,t_\alpha )`$ with respect to x. We note that the main part of the electron distribution is ejected along the polarization direction in our range of parameters, so that we focus here on the photoelectron spectrum in this direction even though it is analogous in the z-direction. The momentum wave function describing the emitted electron distribution in x direction can be thus expressed as
$$\begin{array}{ccc}\mathrm{\Psi }_p(p_x,z,t_f)& =& FFT_x\left[\mathrm{\Psi }_{out}(x,z,t_f)\right]\hfill \\ & & +_{t_\alpha }e^{ip_x^2(t_ft_\alpha )/2}FFT_x\left[\mathrm{\Psi }_{flux}(x,z,t_\alpha )\right].\hfill \end{array}$$
(12)
Here the symbol ’$`FFT_x`$’ means fast Fourier transform of the wave function with respect to the x-coordinate, and $`p_x`$ denotes the electron momentum along the x-axis. The summation in the above equation is for all times $`t_\alpha `$ during which the absorped part of the wave packet is detected. The kinetic energy spectrum of the photoelectron can then be obtained by integrating the momentum wave function over the z-coordinate, i.e. $`P(ϵ_x,t_f)(1+ϵ_x/c^2)/\sqrt{2ϵ_x}\left|\mathrm{\Psi }_p(p_x,z,t_f)𝑑z\right|^2`$. The prefactor has been derived via $`\mathrm{d}p_x/\mathrm{d}e_x`$ where the kinetic energy of the photoelectrons has been approximated via $`ϵ_xp_x^2/2p_x^4/8c^2`$ in atomic units, including only weakly relativistic corrections in second order. The same procedure is carried out for the spin-down wave function and to obtain the total photoelectron spectrum, we sum over both polarisations.
## III Results and Discussions
In this section we present and discuss the role of the magnetic laser component (A) and of the relativistic mass effects (B) on dynamics and radiation of laser driven ions. In (C) spin effects are investigated regarding the Pauli $`\sigma 𝐁`$ term that merely induces spin oscillations and the second order $`\sigma 𝐋`$ term which is responsible for spin induced forces. In (D) harmonic generation is studied in the weakly relativistic regime and in (E) highly energetic above threshold ionization.
### A Magnetic field effects
For neutral atoms interacting with moderately intense laser pulses, the magnetic component of the laser field can usually be ignored in the evaluation of the wavepacket dynamics. Once however the velocity of the electron wavepacket becomes nonnegligible as compared to the velocity of light the Lorentz force $`(𝐯/c)\times 𝐁`$ may not be ignored as compared to the electric field force of the laser field. This generally is the case when the relativity parameter $`r=eE/(m\omega )`$ is not neglible as compared to $`c`$, i.e. at least $`1\%`$ depending on the observable of interest.
The magnetic field component of a linearly polarized laser field induces a drift in the laser propagation direction which may be strong enough to induce instantaneous ionization even for high laser frequencies . There is a regime however where the laser field and correspondingly the drift is weak enough such that the attraction of singly charged ions can still compensate for it and induce recollisions along the propagation direction . If one does not want to be restricted in the laser field intensity, a more elegant solution to overcome the drift and avoid strong ionization would be to employ multiply charged ions.
In the following we focus on the magnetic field (B-field) induced time evolution of the electronic wavepacket in the vicinity of and away from the ionic core, and since second order relativistic effects turn out to be small for the employed laser and ion parameters, the Hamiltonian $`H_{FW}`$ need only include the term $`H_0`$. Thus, we numerically solve the two-dimensional time-dependent Schrödinger equation but beyond the usual dipole approximation and including the laser magnetic field, and the vector potential $`A=A(z,t)`$ consequently depends on both t and z.
In fig.3 we have displayed the dynamics of single electron ions with charge $`Z=3`$ in a) and $`Z=4`$ in b) in a laser field with a wavelength of 248nm and an intensity $`1.2\times 10^{17}W/cm^2`$ during the pulse duration of 10-cycle constant amplitude. We have merely displayed the dynamics of the center of mass of the wavepacket $`(<x(t)>,<z(t)>)`$ in order to compare with our classical intuition. For $`Z=1`$ and $`Z=2`$ we would have almost complete ionization for the laser intensity applied and the figure would ressemble that of a classical free electron drifting with its usual “zick-zack” in the laser propagation direction . For $`Z=3`$ we see in Fig. 3a) that the ionic core starts to seriously compete with the drift imposed by the laser field, i.e. there is substantial ionization in the laser polarization and propagation direction however also a significant part can still be kept close to the ionic core. Interesting to see is that parts of the wavepacket may move opposite to the laser propagation direction similar to the magnetic recollisions advocated earlier . We have also considered the center of mass motion restricted to a smaller area around the vicinity of the nucleus and found dynamics ressembling the “figure of 8 motion” known in the frame where the drift is transformed away. Simulations of the full wave packet have confirmed the dynamics towards and away from the nucleus in the laser propagation direction and the significant amount of ionization.
A different situation occurs in Fig 3b) for $`Z=4`$ when the ion charge is high enough to avoid substantial ionization and to allow for weakly relativistic bound dynamics. The electronic wave packet will oscillate around the ionic core with a particularly high velocity in its close vicinity. At the times when the electronic wave packet approaches the ionic core, the magnetic field component also turns out to be large. Therefore an extremely strong Lorentz push of the wave packet around the ionic core in the laser propagation direction arises and results in a low electron expectation value close to the ionic core . The ”open diamonds” represent the case without magnetic field, i.e. the situation within the dipole approximation. Here the electronic wave packet does merely carry out one-dimensional oscillations along the laser polarization direction, i.e. along the x-axis, as expected. However, when the magnetic component of the laser field is included in the calculation, the electronic wave packet is pushed towards the laser propagation direction (z-axis) as visible from the full circles. This magnetically induced pressure on the electronic wave packet is significantly increased around the ionic core as seen in the figure 3b). It results in a low electron probability in the vicinity of the ionic core. This can be interpreted as a hole around the nucleus at (x=0,z=0) as indicated by the arrow in figure 3b). We need emphasize that we do not put forward holes from knocking out complete inner shell electrons from multi-electron systems. We rather discuss the strong reduction of the wave packet expectation value of single active electrons near the nucleus. For a multi-electron system, however, inner-shell electrons are likely to be more affected by the magnetically and ionic core induced Lorentz push investigated here, simply because the attraction from the nucleus is even larger then.
It may be obvious that the electron velocity near the nucleus is high for multiply charged ions. However, we need explain why the magnetic field, also necessary for the Lorentz force, is significant at the time of the return of the electronic wave packet to the nucleus. When the electronic wave packet returns to the nucleus with essentially maximum velocity, the amplitude of the magnetic component of the laser field should become nearly zero. This is because velocity and force are generally phase shifted by $`\pi /2`$. However, due to the existence of the multiply charged ionic core and the increasing field intensity during the turn-on of the laser field, an important phase lag can be developed during the laser pulse turn-on, which is very critical for the hole formation around the ionic core. During the laser pulse turn-on the electronic wave packet does not move like a quasi-free particle because of the tight binding of the ionic core. When the electronic wave packet is dragged outwards by the laser force, the strong ionic core tries to attract it, which results in an important phase lag between the wave packet motion and the laser field. Usually, when the wave packet returns to the nucleus, this effect is again compensated because this time the ionic core accelerates the electric wave packet towards it. However, because of the increasing field in the turn-on phase of the laser pulse, the laser force is stronger then, the Coulomb force comparably weaker and the phase lag consequently only partially compensated. Once the pulse is at full intensity, the phase lag remains almost unchanged till the field maximal amplitude begins to decrease again.
In the following we investigate the electronic relaxation dynamics after the laser pulse arising from the reduced circular motion around the ionic core. For this purpose, we allow the “hollow” ion to evolve freely for a time period of thirty laser cycles after the laser-ion interaction and during this time observe its relaxation dynamics and radiation. We note that the hollowing out of the ion due to the magnetic field will be reversed step by step in space during the relaxation process. In figure 4.a we have shown the relaxation radiation accompanying this process for the same potential and laser field parameters as in figure 3b). Apparently, the reverse process to magnetically induced hollowing out of the ion gives rise to x-ray emission with distinct peaks up to $`230eV`$ photon energy. The main peak corresponds to the transition between the Stark-shifted first excited state $`|1e>`$ and the ground state $`|g>`$, as indicated in figure 4.b. Further the 122.67eV peak is related to the transition $`|4e>|g>`$. Finally we note, that the hollowing out of the ion may be optimized, by adjusting the interaction parameters such that the Lorentz force $`𝐯/c\times 𝐁`$ becomes maximal near the nucleus.
### B Relativistic Stark shift
In this subsection, we investigate the role of relativistic corrections to the Stark shift as visible in the positions of the energy levels of intense laser driven ions. For this situation we choose the Hamiltonian $`H_{FW}=H_0+H_P+H_{kin}`$. We find that the Pauli term has no notable effect on the shift of energy levels but we include it for the sake of inclusion of all first order relativistic effects. The leading second order term in the weakly relativistic regime is generally $`H_{kin}`$ which may be associated with a relativistic mass-increase. We note that the Darwin term is insignificant in its effect compared to the mass increase and in fact also to spin-orbit coupling. Since in this subsection we are not interested in line splitting but in the total shift of spectral components, we do not include $`H_{so}`$ here. The effect of the mass increase in intense laser driven atoms has been discussed before and associated with an enhancement of stabilization and with Doppler shifts of harmonics . Here we focus however on the multiphoton regime and how the ionic bound energy levels are shifted in addition to the conventional Stark shift.
We choose the model Z=12 ion as described in section II(B) as target. The laser parameters involved are the intensity $`7\times 10^{16}W/cm^2`$, the wavelength 527nm, and a pulse including a 5.25-cycle turn-on phase and 100-cycle with constant amplitude. The whole radiation spectrum is displayed in figure 5, in a) without including $`H_{kin}`$ in the full Hamiltonian and in (b) with including it in the solution of the dynamical equation. The relativistic signatures are small on this scale, however, it is interesting to note that X-rays via hundreds of photons of the applied laser frequency may be generated due to resonances in the multiphoton regime.
We note further numerous resonant structures due to multiphoton transitions among the ionic bound levels in addition to smaller peaks displaced by up to few photon energies of the applied laser field to the resonant lines. In order to have a clearer picture of the structure and the deviation due to the relativistic mass shift, we have depicted the enlargements of the first three dominant resonant spectral structures in figures 6, 7 and 8, respectively for the multiphoton resonances of the ground state $`|g>`$ to the first excited state $`1e>`$, to the second excited state $`|2e>`$ and to the fourth excited state $`4e>`$. The third excited state $`|3e>`$ resembles a d-state and we find that the resonance on $`|3e>|g>`$ involves a coupling three orders of magnitude smaller than those of $`|1e>|g>`$ and $`|4e>|g>`$. We emphasize that our calculations are beyond the dipole approximation so that the usual selection rules do not apply. In these three figures the solid line and dashed line represent the cases in which relativistic corrections due to the mass shift are ignored and included, respectively.
In figure 6, we display the dominant resonant line which is located at $`87.88\omega `$ following the calculation without the relativistic mass shift correction. This corresponds to the energy difference between the ground state and the first excited state including the conventional nonrelativistic Stark shift. When the leading relativistic correction to the electron kinetic energy is taken into account, we find from the dashed line in Fig 6 the corresponding resonance shifts to $`87.45\omega `$, with $`\omega `$ denoting the laser frequency. The additional redshift of the resonance transition of $`|1e>|g>`$ is as high as $`0.43\omega `$ for the parameters chosen here. We refer to the correction as the ”relativistic Stark shift”. Furthermore, we note clear peaks displaced by two laser photons to both sides of the main resonanant lines; i.e. the laser stimulates the absorption or emission of two photons prior the emission of the high frequency photon from the resonant state. The effect of the relativistic Stark shift is clearly visible from the comparison of the dashed and solid line. In addition, small lines are visible which are displaced at about one-photon away from the resonances. Those would be forbidden under the dipole approximation, however appear because of quadric and high order contributions included in our calculation. Those are higher order terms and small in the weakly relativistic regime of interest here. Thus, higher order corrections to the position of this small peak as by the relativistic Stark shift are even smaller and are not visible for this peak displayed by one photon around the $`|1e>|g>`$ transition. However, we note here already that it will be visible for higher excited states, for example for the $`|4e>|g>`$ transition in figure 8.
For the resonance between $`|2e>|g>`$, the result is plotted in figure 7. It is similar to figure 6, but the relativistic Stark shift is only $`0.39\omega `$. As addressed in section II(B), the second excited state $`|2e>`$ is symmetric in space and ressembles an s-state similar to the ground state $`|g>`$, so that the resonant transition $`|2e>|g>`$ is considered. Thus only high order terms beyond the dipole term contribute. Considering the scale of the vertical axis, we find that the resonant signal is three orders of magnitude smaller than in the case for the $`|1e>|g>`$ transition. In addition the relativistic Stark redshift turns out smaller than in the case of the $`|1e>|g>`$ resonance, although the $`|2e>`$ state is not so tightly bound as the state $`|1e>`$. Moreover, the one-photon spaced peaks are no longer visible for the $`|2e>|g>`$ transition. If we continue to explore the $`|4e>|g>`$ resonance, we find that the relativistic Stark redshift becomes relatively large again, as indicated in Fig. 8. Since the state $`|4e>`$ is asymmetric and resembles a p-state as $`|1e>`$, the dipole transition $`|4e>|g>`$ is permitted. Further the state $`|4e>`$ is less strongly bound by the ionic core and consequently much easier influenced by the external field. Therefore, the relativistic Stark shift attains $`0.82\omega `$ and is thus larger than in the cases for the transitions $`|1e>|g>`$ and $`|2e>|g>`$. Finally the peaks spaced one-photon away from the resonance appear again as already in figure 6, however the relativistic Stark redshift becomes visible here also for these small peaks.
Figure 9 shows how the energy levels of the multiply charged ion of interest move under the radiation of an intense laser field. In the left part, the column (a) represents the situation without laser field and the central part (b) indicates the inclusion of the usual Stark shift but without relativistic corrections. The column (c) includes the ”relativistic Stark shift” in addition to the usual Stark shift, which arises from the leading relativistic correction to the kinetic energy of the electron. The detailed position of the states $`|1e>,|2e>`$ and $`|4e>`$ is enlarged in the right part of figure 9. The series of transitions in columns (b) and (c) are corresponding to the solid lines and dotted lines in figures 6, 7 and 8, respectively. This figure gives us a clear picture of how the energy levels shift under an external intense laser field. The relativistic Stark shift, though small when compared to the absolute value of the energy levels, can clearly be identified in the radiation spectra, and should in principle be measurable in experiments.
### C Spin effects
We now turn to the role of the spin degree of freedom of the electron in intense laser-ion interaction. In first order in $`v/c`$ the Pauli equation already includes the coupling of the magnetic laser field to the spin as noted earlier . We will not discuss this here and are more interested in the interaction of the spin degree of freedom with an operator describing the electron rather than the laser field. This is given by spin-orbit coupling arising first in second order in $`v/c`$.
With the laser intensity and ion charge increasing, such that second order effects become important, spin signatures in the dynamics are expected to be visible as noted first recently . For free electrons, similar spin-induced forces have also been noted almost purely analyticaly, especially most interestingly a small one in the direction of the magnetic field of the laser field. Here, we concentrate on the effect of the spin degree of freedom on bound electron dynamics and radiation in intense laser fields. In this case the used Hamiltonian $`H_{FW}`$ includes all second order terms in $`1/c^2`$ as shown in equation (1), of which both the Pauli term $`H_P`$ and the spin-orbit coupling term contribute to spin effects. As noted in subsection A the Lorentz force may induce a significant angular motion around the ionic core with considerable orbital angular momentum $`𝐋`$ with respect to the origin set by the nucleus. We show that the resulting enhanced spin orbit coupling gives rise to observable effects in the electron dynamics and radiation. In particular we note a significant splitting of the non-symmetric bound states due to this additional interaction which leads to well separated doublets and four-line structures in the radiation spectra. In general terms we understand those also as indications that the influence of the spin and other relativistic effects are both principally observable in experiments and nonnegligible in theoretical treatments for relatively low laser intensities.
We are still interested in the weakly relativistic regime of optical laser intensities of up to at about $`10^{17}`$Wcm<sup>-2</sup>, which have been implemented in several laboratories worldwide and which still allows for laser - bound electron dynamics. We employ the 527nm (double Nd:glass) laser with an intensity of $`7\times 10^{16}W/cm^2`$ to interact with Z=12 ions, which are intially spin-up polarized. The laser pulse has a 5.25-cycle turn-on and 100-cycle constant amplitude duration. In figure 10 we address the spin polarization itself and have displayed the expectation value of the electronic wavefunction in “spin-down” configuration as a function of the interaction time. We compare results from the full second order Hamiltonian $`𝐇_{\mathrm{𝐅𝐖}}`$ including spin-orbit coupling with those where spin-orbit coupling $`𝐇_{\mathrm{𝐬𝐨}}`$ has been ignored. Both situations involve a spin flipping with twice the laser frequency. With spin-orbit coupling however the total flipping amplitude is higher because of a second oscillation due to the magnetic field of the frame of reference transformed nucleus $`𝐁^{}`$. Finally the figure shows an effective polarization due to spin-orbit coupling in the turn-on phase, while without spin orbit coupling the electron periodically returns to the initial polarization in complete spin-up configuration.
The most significant qualitative features appear in the radiation spectrum in terms of strongly laser-enhanced line splitting. In figure 11 we have displayed the radiation spectrum of light emitted perpendicular to the plane spanned by the laser polarization and propagation directions and being polarized in $`x`$-direction. The upper row describes the situation governed by the Pauli Hamiltonian while the lower involves the full second order Hamiltoninan $`𝐇_{\mathrm{𝐅𝐖}}`$ in Eq.(1). Figures 11i(a), 11i(b) and 11i(c) show the spectral segments corresponding to the resonances of the first excited state $`|1e>`$ to the ground state $`|g>`$, the third excited state $`|3e>`$ with the ground state $`|g>`$, and the third excited state $`|3e>`$ to the first excited state $`|1e>`$. Comparing the upper and lower row we note shifts and splittings of the spectral components into a doublet in (a) and (b) and a four-line structure in (c). In addition, the relativistic Stark redshift is also observable, as discussed at length in III(B). We found also that the Darwin term due to $`𝐇_𝐃`$ has no notable effects in this situation. We confirmed that the spectral features displayed are generally well separated for ions with different charge states, should a possible experiment not allow for a pure sample of the ion of choice.
We explain the doublets and four-line structure with the splitting of the asymmetrical excited states $`|1e>`$ and $`|3e>`$ due to the additional spin orbit interaction as depicted in Fig. 11ii , while the symmetrical states, possibly $`s`$-states, remain unchanged. The splitting becomes larger with increasing laser intensity or charge of the ionic core. All transitions give rise to spectral features because the common selection rules do not apply in the parameter regime beyond the dipole approximation as investigated here. The bound states in Fig. 11ii drawn with thick lines indicate that those states of the split doublet are most populated explaining the relative heights of the spectral lines in Fig. 11i. The amount of population in the excited states is well above $`1\%`$ (e.g. $`1.25\%`$ in the first excited state) such that the total radiation should not be insignificant. We note that the amount of the splitting is $`\frac{\mathrm{\Delta }E}{\omega _L}0.48`$ for the state $`|1e>`$ and $`0.04`$ for the state $`|3e>`$ (here, $`\omega _L=0.0866`$ a.u. is the applied laser frequency) so that the enhanced spin-orbit splitting should be easily measurable in experiments. Comparing those values with the amount of spin-orbit splitting without the presence of the laser field we have evaluated numerically via the same techniques $`\frac{\mathrm{\Delta }E_0}{\omega _L}0.042`$ for $`|1e>`$ and $`0.004`$ for the state $`|3e>`$. Thus, we find that the total enhancement factor of the spin-orbit splitting due to the intense laser field for our set of parameters amounts to approximately $`10`$. We note that those values should increase considerable for more intense laser fields and for higher charged ions. However, we also emphasize that for less charged ions, in particular hydrogen, spin orbit coupling has still little significance.
### D Generation of Coherent keV X-rays
From nonrelativistic laser atom interaction it is well known that parts of the electronic wavepacket, for appropriate parameters, may tunnel through the Coulomb barrier, then propagate essentially freely in the laser field and when returning to the ionic core in the oscillating field may release their energy in form of radiation . The corresponding radiation spectrum in terms of harmonics of the applied field is in fact quite attractive because quite high frequency coherent light is achievable. Instead of an exponential decay of the spectrum with respect to the harmonic order we here experience a plateau and a cut-off energy of the emitted harmonics at $`I_p+3.17U_p`$. Here $`I_p`$ is the ionization potential of the ions and $`3.17U_p`$ the maximal kinetic energy of the electron at the time of the return to the nucleus. This law is valid in the tunnel regime where the Keldysh parameter need fulfill $`\gamma _K=\sqrt{\frac{I_p}{2U_p}}1`$.
The kinetic energy of electrons increases significantly with the laser intensity. Consequently higher harmonics are expected with an increasing laser intensity if the ionic potential is raising correspondingly, as shown recently already in the nonrelativistic regime with an output of more than 2000 low-frequency harmonics . We here concentrate on the weakly relativitic regime and show that the harmonic order increases considerably by enhancing the charge of the ion and simultaneously increasing the laser field. Here we took special care on adapting the ratio of $`U_p`$ and $`I_p`$ such that we still remain in the tunneling regime.
In our calculations, we employ the $`KrF`$ laser (wavelength $`248nm`$) to interact with multiply charged ions of charge Z=3 and Z=4. There are two advantages to use short-wavelength lasers: (1) the generated harmonics have a high efficiency and are well separated even near the cut-off because of $`\omega _{2n+1}\omega _{2n1}=4\pi /\lambda `$ ($`\lambda `$ being the fundamental laser wavelength); (2) the numerical resolution is better, as the box size is proportional to the square of laser wavelength and the CPU time to the cube. Moreover, the chosen laser wavelength is available in experiments. The pulse envelope is assumed to have a 10-cycle linear turn-on, followed by 10-cycles with constant maximal amplitude. For Z=3 ions (potential parameters k=6.48, $`q_e=1.0`$), we use the laser intensity of $`2.5\times 10^{16}W/cm^2`$. The result is shown in figure 12, in which we find as highest harmonic the 131st order. To our best knowledge, the maximum harmonic order ever obtained in experiments for this short driving wavelength is the 37th . There the authors expressed that they believe that the $`He^+`$ ions may have also contributed to the harmonic emission. Our numerical simulation showns that this harmonic spectrum can be extended to the 131st order if the $`Z=3`$ ions ($`Li^{2+}`$) are employed instead. This harmonic has a wavelength $`1.9nm`$, and the efficiency indicated is still of the level $`10^{10}10^{11}`$ relative to the radiation at the fundamental frequency. Please note in figure 12 that we have only shown the radiation polarized in the laser polarization direction as the radiation polarized in the laser propagation direction is relatively small for the parameters chosen here.
Coherent X-rays even in the keV regime can be obtained when $`Z=4`$ ions (potential parameters k=10.7, $`q_e=1.0`$) are used. In figure 13, we display the harmonic spectrum of $`Z=4`$ ions driven by a laser pulse with the same parameters as in the previous figure but with an intensity of $`10^{17}W/cm^2`$. The spectrum polarized in the laser polarization direction in a) has been evaluated via the Fourier transform of $`a_x(t)`$, and the smaller contribution polarized in the laser propagation direction in b) via $`a_z(t)`$ in Eq. (10). The cut-off is enlarged as an insert in the right corner of figure 13 a). We emphasize that coherent radiation of the 427th harmonic is clearly observable in spite of the low efficiency of $`10^{13}`$ relative to the radiation at the fundamental frequency. The photon energy of this harmonic exceeds even 2keV and shall be useful for many applications, as e.g. for time-resolved X-ray diffraction. For the parameters chosen we find no significant deviation to the cut-off rule, even though the factor 3.17 is expected to alter with more than weakly relativitic free electron dynamics changing the kinetic energy at the recollision time. Most remarkably, and unfortunately, the plateau is tilting with increasing charge reducing the efficiency of the highest harmonics substantially. Thus our results indicate that a part of the electronic wave packet can still tunnel through the barrier formed by a deep ionic potential but the amount that tunnels out and returns reduces drastically with the ion charge. Thus with further increasing charge higher harmonic energies are possible but special care is devoted to the small efficiency.
The calculations for figure 13 have been carried involving all first order terms in $`v/c`$, i.e. including the role of the magnetic field, and the leading terms in $`(v/c)^2`$. We have chosen parameters in the weakly relativistic regime such that higher terms will not make a difference in the results presented. We note that the spectrum in figure 13 a) in fact would look hardly different without all the terms in $`v/c`$, however the one in figure 13 b) would hardly ressemble the correct one depicted. This is not surprising as the laser induced dynamics in the laser propagation direction is fully induced by the magnetic field component of the laser field.
Regarding experimental possibilities for high harmonics from multiply charged ions in the coherent keV X-rays, there are two ways at hand. There may be the double pulse experiment where the first pulse is used to strip several or many electrons of the neutral atoms and then a moderately delayed pulse will interact with the obtained multiply charged ions and generate the coherent keV X-rays as described above. However, here the free electrons ionized by the first pulse may influence the phase matching of harmonic emissions and there may also be a mixture of many charge states. The alternative way is to shoot atoms through thin foils which is clearly a cleaner way to generate ions with a pure charge state . Extremely high densities of multiply charged ions have been generated recently , which raises hope for a reasonable efficiency for ion charges well above those investigated here.
### E High-order Above-threshold-ionization
When laser field strengths are employed which are high with respect to the binding fields of the ionic core, a large part of the electron wavepacket will escape the vicinity of the ion and not return. Those electrons may in fact be quite energetic and will be of interest in this section.
We begin by calculating the photoelectron spectrum for $`Z=3`$ ions with the method described in section II(C). Here, we use a laser pulse with a 3-cycle turn-on, and 10-cycles with constant maximal amplitude. The applied maximal intensity is equal to $`2.5\times 10^{16}W/cm^2`$, and the laser wavelength is 248nm (corresponding to a photon energy near $`5eV`$). The resulting photoelectron spectrum is displayed in figure 14. The spectral range from 860eV to 900eV is enlarged as an insert in the right-up corner of this figure. We note very energic electrons with energies up to 2keV with the usual photon-energy spacing characteristics of photoelectron peaks. Figure 15 shows the case for $`Z=4`$ ions with an appropriately higher laser intensity of $`1.2\times 10^{17}W/cm^2`$. As expected more energetic (above 5keV) photoelectrons can be deteced. A rather regular spacing of the peaks in this above-threshold-ionization spectrum is still notable in this keV energy regime.
Thus next to extremely high harmonics we find also very energetic electrons due to above-threshold ionization. In more general terms those results and the ones in the previous section show that very high-order nonlinear effects are governing the interaction of multiply charged ions with very intense laser fields.
## IV Conclusion
We believe to have shown that the physics of multiply charged ions in very intense laser fields is even richer than that for neutral atoms with moderately intense laser fields. There is not merely the effect of scaling the known effects for neutral atoms to the new intensity regime but the upcoming of many relativistic influences imposes a fundamentally different dynamics. We have seen that the magnetic field component of the laser field completely modifies the dynamics and may even induce a partially circular motion around the nucleus with the effect of a reduced expectation value just in the near vicinity of the ionic core. In the regime where second order terms in $`v/c`$ are important, the relativistic mass shift modifies clearly the positions of the spectral components in the multiphoton regime. Amplification was found on resonances involving of the order of 100 photons with interesting non-dipole spectral features displaced few harmonics away from the atomic resonances. The spin, usually of no importance in the nonrelativistic regime, starts to oscillate and via spin-orbit coupling modifies substantially dynamics and radiation. The laser enhanced splitting of resonant spectral lines is a clear relativistic quantum signature.
Those aspects are associated with fundamentally new influences in the weakly relativistic regime with respect to the nonrelativistic case and were most conveniently described by applying the expansion of the Dirac equation up to the second order in $`v/c`$. This allowed us to relate each effect to different parts of the Hamiltonian and to carry out our numerical investigations even for the low frequencies available in most present day high power laser systems.
We stressed also that the combination of highly charged ions and high power lasers can be useful for applications. High-order above threshold ionization was shown to give rise to photoelectrons in the multi-keV regime already for laser intensities around $`10^{17}W/cm^2`$. More interestingly parts of those highly energetic electron wavepackets may return to the nucleus and we indicated also coherent high harmonics in the keV regime. Even though for high harmonic generation towards the coherent hard X-ray regime the quantitative aspect appears most attractive, we pointed out also qualitative changes as the problematic titling of the plateau of the harmonic spectrum with increasing ion charge.
The authors acknowledge funding from the German Research Foundation (Nachwuchsgruppe within Sonderforschungsbereich 276). SXH would like to thank W. Becker for fruitful discussions and acknowledges present funding from the Alexander von Humboldt Foundation. CHK acknowledges helpful discussions with M. Casu, J. Ullrich and M. W. Walser.
Present address: Max-Born-Institut for nonlinear Optics and Short Pulse Spectroscopy, Rudower Chaussee 6, 12489 Berlin, Germany. Present email: suxinghu@mbi-berlin.de
<sup>∗∗</sup> Email address: keitel@physik.uni-freiburg.de
Fig. 1: S. X. Hu and C. H. Keitel, “Dynamics of …”
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# 1 Introduction
## 1 Introduction
During the last decade an increase interest has been shown for the genuine multiparticle correlations in multihadron final states of hadronic, e<sup>+</sup>e<sup>-</sup> and other reactions . Recently OPAL, in its study of e<sup>+</sup>e<sup>-</sup> annihilations on the Z<sup>0</sup> mass, has established the existence of strong genuine multihadron correlations up to the fifth order . In hadron-hadron, like proton-proton, collisions the correlations of more than three particles have also been observed . In contrast to this situation, in heavy ion collisions, at low energies and/or in reactions of light nuclei, genuine correlations are found to have non-zero values only up to the third order . Furthermore it has been found out that in general these correlations become weaker as the reaction average multiplicity increases. In nucleus-nucleus collisions at high energies, of tens and hundreds GeV per nucleon, the two-particle correlations are the only one that survive .
This correlation dependence on the average multiplicity is very similar to the one observed in the investigations of multiparticle dynamical fluctuations, i.e. variation of many particle bunches in restricted phase space regions . In these studies, known as intermittency analyses, the average multiplicity dependence has been proposed to be the consequence of a mixing of several independent emission sources . As a result, the dynamical fluctuations in nucleus-nucleus collisions are already well accounted for by two-particle correlations , whereas in hadron-hadron interactions and in e<sup>+</sup>e<sup>-</sup> annihilations higher order genuine correlations do exist.
An analogous situation is also observed in Bose-Einstein correlations (BEC) where identical bosons are correlated when they emerge from the interaction in nearby phase space. A genuine three-pion BEC has been detected in hadron-hadron reactions and found to be even more pronounced in e<sup>+</sup>e<sup>-</sup> annihilations . In contrast to these, no genuine three-pion BEC were found in nucleus-nucleus collisions where the three-body correlations were well reproduced in terms of two-particle BEC . Since the intermittency phenomenon and BEC seem to be closely related , the dependence of many sources on the strength of the BEC cannot be excluded. The superposition of emitters may also be a reason for the suppression of BEC of hadrons produced from W-boson pairs in e<sup>+</sup>e<sup>-</sup> annihilations at LEP2 energies where the overlapping of hadrons affect the accuracy of the W mass measurements .
All this, as well as the obvious intrinsic interest in the genuine correlations which carry most of the dynamics of the hadron production process, points to the need of dedicated studies aimed to investigate the correlation dependence on the number of emission sources. Here we propose to study this dependence by grouping several e<sup>+</sup>e<sup>-</sup> $`hadrons`$ events to represent a multi-emission sources of particles. To obtain significant results, even when only few sources are considered, one needs a very high statistics, like that which can be supplied by a Monte Carlo (MC) generated sample, to minimise the calculation error and thus be sensitive to the correlations dependence on the number of sources. Another advantage of using MC generated events is in its possibility to generate a multihadron sample free from contamination of other processes like, for example, e<sup>+</sup>e<sup>-</sup> $`\tau ^+\tau ^{}hadrons`$.
In this Letter we present the results of a MC study on the effect of several emission sources on the genuine higher order multiparticle correlations. The study was based on a generated sample of about $`5\times 10^6`$ events of the reaction e<sup>+</sup>e<sup>-</sup> $``$ $`\mathrm{Z}^0`$ $`hadrons`$ which passed a full simulation of the OPAL detector at LEP and did reproduce rather well the measured genuine high order correlations present in the OPAL hadronic $`\mathrm{Z}^0`$ decay data . Moreover, the e<sup>+</sup>e<sup>-</sup> $``$ $`\mathrm{Z}^0`$ $`hadrons`$ annihilations should represent well the one emission source situation in contrast to events produced in hadron-hadron interactions. Our analysis on the dependence of many sources on the genuine correlations was thus carried out in a way that avoided effects of other reaction features, like the multiplicity which was discussed recently in connection with two-particle BEC analyses in and . In as much that final state interactions between hadrons coming from different sources can be neglected, our correlation study based on e<sup>+</sup>e<sup>-</sup> $``$ $`\mathrm{Z}^0`$ $`hadrons`$ annihilations may be extended to other types of reactions since the hadronisation process is believed to rest on a common basis .
## 2 The analysis method
The analysis is based on a generated sample of hadronic Z<sup>0</sup> decays using the Jetset 7.4 MC program including a full simulation of the OPAL detector at LEP . The MC sample also included initial-state radiation and effects of finite lifetimes. The parameters of the program were tuned to yield a good description of the measured event shape and single particle distributions .
The selection criteria for multihadron events used here are identical to the ones previously utilised by OPAL in their recent data analysis of multiparticle correlations . In particular, selected events were required to have at least five charged tracks each having at least 20 measured points in the jet chamber where the first point had to be closer than 40 cm from the beam axis. The cosine of the polar angle of the event sphericity axis with respect to the beam direction was required to be less than 0.7 to ensure the event to be within the volume of the detector. The sphericity axis was calculated by using all accepted tracks and electromagnetic and hadronic calorimeter clusters.
To simulate several emission sources we did overlay several e<sup>+</sup>e<sup>-</sup> $``$ $`\mathrm{Z}^0`$ $`hadrons`$ generated events and analysed the correlations between pions as if they were created in a single event. The kinematic variables are defined within each generated e<sup>+</sup>e<sup>-</sup> event with respect to its own sphericity axis. For correlation analyses of variables like rapidity this procedure is equivalent to the one where the events are rotate to a common sphericity axis. This simulates multi-sources’ events with hadrons, here all taken to be pions, emerging from a common emitter. To note is that in this procedure the average event multiplicity is directly proportional to the number of sources. In our analysis each generated event was used only once, and hence required a very large MC event statistics.
To extract the genuine dynamical $`q`$-particle correlations, we used bin-averaged normalised factorial cumulant moments, or cumulants, first proposed in Ref. as a tool for search for genuine multiparticle correlation,
$$K_q=\frac{1}{M}\underset{m=1}{\overset{M}{}}_{\delta y}\underset{i}{}\mathrm{d}y_i\frac{C_q(y_1,\mathrm{},y_q)}{[_{\delta y}dy\rho _1(y)]^q},$$
(1)
where $`C_q(y_1,\mathrm{},y_q)`$ are the $`q`$-particle correlation functions given by the inclusive $`q`$-particle density distributions $`\rho _q(y_1,\mathrm{},y_q)`$ in terms of cluster expansion, e.g.,
$$C_3(y_1,y_2,y_3)=\rho _3(y_1,y_2,y_3)\underset{(3)}{}\rho _1(y_1)\rho _2(y_2,y_3)+2\rho _1(y_1)\rho _1(y_2)\rho _1(y_3).$$
(2)
Here $`M`$ is the number of equal bins of a width $`\delta y`$ into which the event phase-space is divided and the subscript (3) denotes the number of permutations. For simplicity we show all formulae in one-dimensional (e.g., rapidity) phase space.
The feature of the $`C_q`$-functions is that they vanish whenever there are no genuine correlations, i.e., the correlations are due to those existing in lower orders. The correlations extracted are of a dynamical nature since the cumulants share with normalised factorial moments (the intermittency analysis tool) the property of statistical noise suppression.
In this paper we computed the cumulants as they were used in experimental studies, in particular we used the form applied in Ref. namely,
$$K_q=\frac{𝒩^q_{m=1}^Mk_q^{(m)}}{_{m=1}^MN_m(N_m1)\mathrm{}(N_mq+1)}.$$
(3)
Here, the factors $`k_q^{(m)}`$ are the unnormalised factorial cumulant moments, or the Mueller moments , calculated for the $`m`$th bin. These factors represent the correlation functions $`C_q`$ integrated over the bin and $`N_m`$ is the number of particles in the $`m`$th bin summed over all the $`𝒩`$ events. The definition (3) takes into account the non-uniform shape of the single-particle distribution and the bias when the cumulants are computed at small bins.
The cumulant calculations were performed in the three-dimensional phase space of the kinematic variables commonly utilised in this kind of studies , namely:
* The rapidity, $`y=\mathrm{ln}\sqrt{(E+p_{})/(Ep_{})}`$, with $`E`$ and $`p_{}`$ being the energy and longitudinal momentum of the hadron in the interval $`2.0y2.0`$;
* The transverse momentum in the interval $`0.09p_T2.0`$ GeV/$`c`$;
* The azimuthal angle, $`0\mathrm{\Phi }<2\pi `$, calculated with respect to the eigenvector of the momentum tensor having the smallest eigenvalue, in the plane perpendicular to the sphericity axis.
These variables are defined with respect to the sphericity axis, in a way and within the intervals similar to those used in a recent OPAL analysis and in other cumulant studies .
## 3 Genuine correlations and the number of sources
### 3.1 Monte Carlo studies
In Fig. 1 we compare the MC based cumulants of orders $`q=2,3`$ and 4 calculated from a single e<sup>+</sup>e<sup>-</sup> $``$ $`\mathrm{Z}^0`$ $`hadrons`$ source (solid symbols) with those obtained by overlaying seven such events representing seven hadronic sources (open symbols). The calculations were performed in one-dimensional rapidity, in two-dimensional rapidity vs. azimuthal angle subspaces and in three-dimensional phase space of rapidity, azimuthal angle and transverse momentum.
The following observations can clearly be made from Fig. 1.
* The existence of a dynamical component, i.e. rise of the cumulants with increasing number of bins $`M`$, is seen to be present both in the single source as well as in the case of many sources. Although the slopes of this scaling behaviour are smaller for several sources than for a single source, they are still strongly present. It is also evident that the character of the scaling is kept as the number of sources increases e.g., for one source as well as for seven sources the one-dimensional (rapidity) cumulants level off at the same $`M`$ value. No such saturation exists for the one and seven sources cumulants in the two and three dimensions.
* The genuine dynamical correlations, measured by the cumulants, significantly decrease with the increase of the number of sources. This decrease is stronger for higher order correlations namely, whereas the two-particle cumulants suffer a reduction of an order of magnitude as the number of sources increases from one to seven, the four-particle cumulants diminish by three or four orders of magnitude.
* The hierarchy of the $`K_q`$ cumulants is reversed as the number of sources increases. The cumulants derived from the single-source events increase with increasing $`q`$-order so that $`K_2^{(1)}<K_3^{(1)}<K_4^{(1)}`$, whereas the hierarchy in the cumulants calculated for seven sources is reversed namely, $`K_2^{(7)}>K_3^{(7)}K_4^{(7)}`$. In addition, the multi-sources cumulants of order $`q>2`$ have almost the same reduced value namely, $`K_3^{(7)}K_4^{(7)}\stackrel{<}{_{}}𝒪(0.1)`$. This last feature does not change as the dimension increases.
* The overall dominant feature of the analysis results is the diminishing value of the higher order cumulants as the sources number increases leaving the $`K_2`$ to be the dominant genuine multiparticle correlations.
The observed diminished correlations is further illustrated in Fig. 2 where the $`M`$-averaged one dimensional (rapidity) $`K_q^{\mathrm{av}}`$ cumulants, are plotted with their errors against the number of sources. The cumulants were averaged over the $`M10`$ region where they are seen in Fig. 1 to approach an almost constant value. The values of these $`K_q^{\mathrm{av}}`$ are also listed in Table 1 together with the OPAL measured data cumulants of single e<sup>+</sup>e<sup>-</sup> $``$ $`\mathrm{Z}^0`$ $`hadrons`$ events. These data cumulants for $`q3`$ are seen to agree with those derived from the MC sample. The one source MC based $`q=4`$ cumulant lies lower than the measured data value but is still consistent within errors.
It is obvious from Fig. 2 that the $`M`$-averaged rapidity cumulants of order $`q>2`$ decrease fast with the increase of the number of sources. Already for two sources the hierarchy changes and the two-particle correlations visibly dominate over the higher order ones. At higher number of sources the dominant role of the two-particle correlations is even more pronounced.
### 3.2 Correlation dilution due to source mixing
Within the procedure adopted here for the simulation of multi-source events, it is clear that if a genuine correlation exists it can only be detected in groups of $`q`$ pions emerging from the very same source. In those $`q`$-group combinations which emerge from at least two sources, genuine correlations should not be present. This means that for $`K_q`$ cumulants that are calculated over all possible $`q`$-pion groups, the higher the number of sources the more diluted will be the signal for genuine correlations.
For a given $`q`$-order the genuine correlation dilution factor is thus:
$$R_q=\frac{P_q^\mathrm{G}}{(P_q^\mathrm{G}+P_q^{\mathrm{NG}})},$$
(4)
where $`P_q^\mathrm{G}`$ denotes the number of $`q`$particle groups, e.g., pairs or triplets of pions, which emerge from the same source. The term $`P_q^{\mathrm{NG}}`$ stands for the number of all possible combinations of $`q`$particle groups which emerge from at least two sources.
Since all sources are produced in the same reaction and at the same energy, they do have an identical average charged multiplicity. For the estimation of $`R_q`$ we assume that all the $`S`$ sources have the same fixed charged multiplicity $`n`$. In this case one has $`P_q^\mathrm{G}=S\left(\genfrac{}{}{0pt}{}{n}{q}\right)`$, and the dilution factors at $`q=2`$, 3 and 4 are given by
$`R_2`$ $`=`$ $`{\displaystyle \frac{\left(\genfrac{}{}{0pt}{}{n}{2}\right)S}{\left(\genfrac{}{}{0pt}{}{n}{2}\right)S+n^2\left(\genfrac{}{}{0pt}{}{S}{2}\right)}}\stackrel{n1}{}{\displaystyle \frac{1}{S}},`$
$`R_3`$ $`=`$ $`{\displaystyle \frac{\left(\genfrac{}{}{0pt}{}{n}{3}\right)S}{\left(\genfrac{}{}{0pt}{}{n}{3}\right)S+n^3\left(\genfrac{}{}{0pt}{}{S}{3}\right)+2n\left(\genfrac{}{}{0pt}{}{n}{2}\right)\left(\genfrac{}{}{0pt}{}{S}{2}\right)}}\stackrel{n2}{}{\displaystyle \frac{1}{S^2}},`$ (5)
$`R_4`$ $`=`$ $`{\displaystyle \frac{\left(\genfrac{}{}{0pt}{}{n}{4}\right)S}{\left(\genfrac{}{}{0pt}{}{n}{4}\right)S+n^4\left(\genfrac{}{}{0pt}{}{S}{4}\right)+3n^2\left(\genfrac{}{}{0pt}{}{n}{2}\right)\left(\genfrac{}{}{0pt}{}{S}{3}\right)+\left(\genfrac{}{}{0pt}{}{n}{2}\right)^2\left(\genfrac{}{}{0pt}{}{S}{2}\right)+2n\left(\genfrac{}{}{0pt}{}{n}{3}\right)\left(\genfrac{}{}{0pt}{}{S}{2}\right)}}\stackrel{n3}{}{\displaystyle \frac{1}{S^3}},`$
where the denominators include the number of all possible $`q`$-particle combinations in $`S`$ sources of charged multiplicity $`n`$. This dilution factor dependence on the number of sources can also be derived in terms of cumulants .
From these $`R_q`$ relations one can show that as long as $`nq`$ one obtains a general expression for the dilution factor,
$$R_q\stackrel{nq}{}\frac{1}{S^{q1}}.$$
(6)
To compare the dilution factors $`R_q`$ with our correlation results shown in Fig. 2, they do have to be multiplied by $`K_q^{\mathrm{av}(1)}`$ which is a measure of the genuine $`q`$-order correlation present in a single sources. The solid lines shown in Fig. 2 thus represent the diluted cumulants $`K_q^{\mathrm{av}}=K_q^{\mathrm{av}(1)}\times R_q`$. The striped areas in which the lines are embedded are the allowed regions when $`q`$ is not neglected with respect to the multiplicity $`n`$. The agreement between the cumulant calculations and the dilution factors predictions is really remarkable for $`q=2`$ and $`3`$ and certainly is still well within the rather large errors of the $`q=4`$ cumulants.
For the order $`q=2`$ one can relax the fixed charged multiplicity assumption and allow them to be different and still retain the $`R_q1/S^{q1}`$ relation as long as the multiplicity distribution is of a Poisson nature. This however is not the case for orders higher than 2. Nevertheless for order $`q=3`$ the relation $`R_3=1/S^2`$, derived from the fixed multiplicity assumption, is still valid as it describes well the $`K_3^{\mathrm{av}}`$ values up to at least thirteen sources (see Fig. 2). The large cumulants’ errors associated with the $`q=4`$ order prohibits to judge how accurate is the $`R_4=1/S^3`$ relation.
An additional interesting and useful application of the relation $`R_q1/S^{q1}`$ is that it offers a method to estimate the average number of sources $`S`$ via the cumulant averaged values over the large $`M`$ region of two sequential $`q`$-orders through the ratio,
$$S\frac{K_{q+1}^{\mathrm{av}(1)}}{K_q^{\mathrm{av}(1)}}\times \frac{K_q^{\mathrm{av}}}{K_{q+1}^{\mathrm{av}}}.$$
(7)
### 3.3 Comparison with hadron and nucleus induced reactions
As is already mentioned in the introduction, the genuine correlations measured in e<sup>+</sup>e<sup>-</sup> annihilations are found to be weaker in hadronic interactions and even more so in nuclear collisions . In nucleus-nucleus collisions at ultra-relativistic energies only the second-order correlations were so far detected .
In Table 2 we list the results obtained by several experiments on the $`M`$-averaged rapidity cumulant values for $`q=2`$ and 3. These average values were taken over the $`M`$-region where the cumulants are seen in the published figures to reach a constant level. The hadronic reactions and their cumulants values are ordered according to their reported mean charged multiplicity, from the lowest value to the highest one.
Table 2 shows that in hadron including nucleus induced reactions the two-particle correlations decrease rather fast as the mean multiplicity increases. However the three-particle correlations are found to be essentially non-existing even at moderately small mean multiplicity. Notwithstanding the possibility that production of hadrons in e<sup>+</sup>e<sup>-</sup> annihilation may well be simpler than in hadron induced reactions, it may nevertheless be instructive to relate our findings to the measured correlation data listed in Table 2. In as much that the mean multiplicity increases with the number of sources, the decrease in the two-particle correlations and the absence of three-pion correlations in nucleus induced reactions is consistent with our findings which demonstrated the dilution of the correlations with increased number of sources. A quantitative comparison between our findings and the correlations in nucleus-nucleus and hadron induced reactions is hard mainly because of the large errors associated with the average cumulant values. In particular the application of relation (7) is prohibited because most of the $`K_3^{\mathrm{av}}`$ are consistent within errors with zero.
Recently the two-pion BEC have been studied in $`\overline{\mathrm{p}}`$p reaction at centre of mass energy of 630 GeV as a function of multiplicity by using the normalised cumulants method similar to the one used here. In that analysis it has been found that the correlations of the cumulants of the like-sign pions as well as the opposite-sign pions decrease with the multiplicity $`n_{ch}`$. From our analysis we expect the pair correlation to decrease as $`1/S`$, where $`S`$ is the number of sources. This indicates that indeed the multiplicity is at least partially proportional to $`S`$. The BEC dependence has also been investigated in the framework of the totally coherent emission picture and in the quantum optical approach where the conclusions were that these correlations are weaker as the multiplicity increases.
## 4 Summary and conclusions
To investigate the effect of many emission sources on the genuine correlations in mutihadron final state we adopted a procedure which should minimise the confusion introduced by other variables like charged multiplicity. For the genuine correlations measurement we utilised the normalised cumulant method. To simulate the situation of many sources event we did overlay Monte Carlo generated hadronic Z<sup>0</sup> events treating them as one event. This Jetset7.4 MC sample of some five million events, tuned to the OPAL data taken at LEP1 on the Z<sup>o</sup> mass, has previously described rather well the measured correlations in the e<sup>+</sup>e<sup>-</sup> $``$ $`\mathrm{Z}^0`$ $`hadrons`$ data.
The results obtained here show that the cumulants, obtained from a single-source events and from events of many sources, almost do not change their structure with the decrease of the width of phase-space bins. This means that the scaling is preserved although larger slopes are seen to be in the case of one source compared with those for several sources. However, when the number of sources increases, the cumulants of order higher than two are suppressed and diminish to zero due to source mixing. The two-particle cumulants are also somewhat reduced in their value but they are way above the higher order ones and they are seen to completely dominate when the source number $`S`$ exceeds the value ten.
The one-dimensional (rapidity) correlations are very well reproduced by assuming that genuine correlation of the order $`q`$ can only be present when all the $`q`$ hadrons are emerging from the same source. Therefore the dilution of the genuine correlation signal is proportional to the ratio of the probability that the $`q`$ hadrons will come from the same source. From simple combinatorial considerations this probability is approximately equal to $`1/S^{q1}`$. Thus a measurement of the one-dimensional correlation for two sequential orders renders the number of sources.
The genuine correlations measured in hadron and nucleus induced reactions do follow qualitatively the findings of our work. In particular in nucleus-nucleus reactions, where many sources are expected to contribute to the final hadronic state, the $`q>2`$ orders are very small and indeed consistent with zero. The $`q=2`$ order is still present but it is also getting smaller as the atomic number of the nuclei increases. The general decrease of the second order cumulants with the increase of multiplicity re-confirms the belief that the higher the multiplicity the larger the number of sources. Our results may also be useful for the understanding of other types of measured correlations like the Bose-Einstein interference of two and more identical bosons. It has been previously pointed out that in the absence of final state interactions the BEC of the e<sup>+</sup>e$`{}_{}{}^{}`$ W<sup>+</sup>W$`{}_{}{}^{}hadrons`$ will be half of that of the $`\mathrm{Z}^0`$ decay to hadrons. From our study it follows that the two-particle BEC, or any other correlations, in the two-source reaction e<sup>+</sup>e$`{}_{}{}^{}`$ W<sup>+</sup>W$`{}_{}{}^{}hadrons`$ should be reduced by a factor two as compared to that of the hadrons emerging from one W-boson.
## Acknowledgements
We would like to thank our colleagues from the OPAL Collaboration for allowing us to use the Monte Carlo sample of Z<sup>o</sup> decays into hadrons. Special thanks are also due G. Bella, S. Kananov and S. Nussinov for many helpful and inspiring discussions.
Calculations are based on the measured factorial moments.
<sup>∗∗</sup> Semicentral collisions.
<sup>∗∗∗</sup> Central collisions.
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# FERMILAB-PUB-00/035-T BNL-NT-00/14 Soft Double–Diffractive Higgs Production at Hadron Colliders
## 1 Introduction
In this letter we suggest a new mechanism for “soft” double–diffractive production of Higgs boson. We consider three reactions
$`p+p`$ $``$ $`p+[LRG]+H+[LRG]+p;`$ (1)
$`p+p`$ $``$ $`X_1+[LRG]+H+[LRG]+X_2;`$ (2)
$`p+p`$ $``$ $`p+[LRG]+H+[LRG]+X_2;`$ (3)
where LRG denotes the large rapidity gap between produced particles and $`X`$ corresponds to a system of hadrons with masses much smaller than the total energy. These reactions have such a clean signature for experimental search (see Fig.1, where the lego - plot is shown for reaction of Eq. (1)) that they have been the subject of continuing theoretical studies during this decade ( see Refs.).
The main idea behind all calculations, starting from the Bialas-Landshoff paper , is to describe the reactions of Eq. (1) and Eq. (2) as a double Pomeron (DP) Higgs production ( see Fig.2 ) . In Fig.2, the Pomerons are the so–called “soft” Pomerons for which one uses the phenomenological Donnachie-Landshoff form ( see Ref. ), while the vertex $`\gamma `$ can be calculated in perturbative QCD.
We can demonstrate the problems and uncertainties of such kind of approach by considering the simplest pQCD diagram for the double Pomeron Higgs production ( DPHP ) (see Fig.3-a). This diagram leads to the amplitude
$$M(qqqHq)=\frac{2}{9}\mathrm{\hspace{0.17em}2}g_H\frac{d^2Q_{}}{Q_{}^2Q_{1,}^2Q_{2,}^2}\mathrm{\hspace{0.17em}4}\alpha _S(Q_{}^2)(\stackrel{}{Q}_{1,}\stackrel{}{Q}_{2,}),$$
(4)
where $`g_H`$ is the Higgs coupling that has been evaluated in perturbative QCD . For the reaction of Eq. (1), $`|t_1|=|\stackrel{}{Q}_{}\stackrel{}{Q}_{1,}||t_2|=|\stackrel{}{Q}_{}\stackrel{}{Q}_{2,}|\mathrm{\hspace{0.17em}2}/B_{el}`$ and therefore,
$$M(q+qq+H+q)\frac{d^2Q_{}}{Q_{}^4}.$$
(5)
Eq. (5) has an infrared divergence which is regularized by the size of the colliding hadrons. In other words, one can see that already the simplest diagrams show that the DP Higgs production is, in a sense, a “soft” process. Taking into account the emission of extra gluons denoted in Fig.3-b as Pomeron builders, we recover the exchange of the “soft” Pomerons.
Nevertheless, the emission vertex for the Higgs boson can still be calculated in pQCD since the typical distances inside the quark triangle in fig.3-a are rather short $`1/M_T`$, where $`M_T`$ is the mass of t-quark. The coupling $`g_H`$ has been evaluated in Ref. and is given by
$$g_H^2=\sqrt{2}G_F\alpha _S^2(M_H^2)N^2/9\pi ^2,$$
(6)
where $`N`$ is a function of the ratio $`M_T/M_H`$ which was calculated in Refs. .
In this paper we consider an alternative approach to DPHP, in which we estimate the value of the cross section from non–perturbative QCD. In section 2 we review a non–perturbative method suggested by Shifman, Vainshtein and Zakharov for the evaluation of the coupling of the Higgs boson to hadrons; it is valid if the mass of the Higgs is smaller than the mass of the top quark. In section 3 we develop a method of obtaining the DPHP cross section using the approach of Ref. . The problem of survival of large rapidity gaps (LRG) will be discussed in section 4. We conclude in section 5 with discussion of our results and of the uncertainties inherent to our approach.
## 2 The coupling of Higgs boson to hadrons in <br>non-perturbative QCD
To evaluate the non–perturbative coupling of the Higgs boson to hadrons, we need to have a closer look at the properties of the energy–momentum tensor in QCD. The trace of this tensor is given by
$$\mathrm{\Theta }_\alpha ^\alpha =\frac{\beta (g)}{2g}G^{\alpha \beta a}G_{\alpha \beta }^a+\underset{l=u,d,s}{}m_l(1+\gamma _{m_l})\overline{q_l}q_l+\underset{h=c,b,t}{}m_h(1+\gamma _{m_h})\overline{Q_h}Q_h,$$
(7)
where $`\gamma _m`$ are the anomalous dimensions; in the following we will assume that the current quark masses are redefined as $`(1+\gamma _m)m`$. The appearance of the scalar gluon operator in (7) is the consequence of scale anomaly , . The QCD beta function can be written as
$$\beta (g)=b\frac{g^3}{16\pi ^2}+\mathrm{},b=9\frac{2}{3}n_h,$$
(8)
where $`n_h`$ is the number of heavy flavors ($`c,b,..`$). Since there is no valence heavy quarks inside light hadrons, at scales $`Q^2<4m_h^2`$ one expects decoupling of heavy flavors. This decoupling was consistently treated in the framework of the heavy-quark expansion ; to order $`1/m_h`$, only the triangle graph with external gluon lines contributes. Explicit calculation shows that the heavy-quark terms transform in the piece of the anomalous gluonic part of $`\mathrm{\Theta }_\alpha ^\alpha `$:
$$\underset{h}{}m_h\overline{Q_h}Q_h\frac{2}{3}n_h\frac{g^2}{32\pi ^2}G^{\alpha \beta a}G_{\alpha \beta }^a+\mathrm{}$$
(9)
It is immediate to see from (9), (7) and (8) that the heavy-quark terms indeed cancel the part of anomalous gluonic term associated with heavy flavors, so that the matrix element of the energy–momentum tensor can be rewritten in the form
$$\mathrm{\Theta }_\alpha ^\alpha =\frac{\stackrel{~}{\beta }(g)}{2g}G^{\alpha \beta a}G_{\alpha \beta }^a+\underset{l=u,d,s}{}m_l\overline{q_l}q_l,$$
(10)
where heavy quarks do not appear at all; the beta function in (3.10) includes the contributions of light flavors only:
$$\stackrel{~}{\beta }(g)=9\frac{g^3}{16\pi ^2}+\mathrm{}$$
(11)
Because the mass of the Higgs boson $`M_H`$ is presumably large, its coupling to hadrons involves the knowledge of hadronic matrix elements at the scale $`Q^2M_H^2`$, at which the heavy quarks in general are not expected to decouple. However, if the Higgs boson mass $`M_H`$ is smaller than the mass of the top quark $`M_T`$, one can still perform expansion in the ratio $`M_H/M_T`$; we expect this to be a reasonable procedure if $`M_H100`$ GeV. In this case, one finds
$$M_T\overline{t}t\frac{2}{3}\frac{g^2}{32\pi ^2}G^{\alpha \beta a}G_{\alpha \beta }^a+\mathrm{}$$
(12)
Since the mass of the hadron is defined as the forward matrix element of the energy–momentum tensor, the expression Eq. (12) leads to the following Yukawa vertex for the coupling of a Higgs boson to the hadron:
$$2^{\frac{1}{4}}G_F^{\frac{1}{2}}H\varphi _h^2h|M_T\overline{t}t|h=\mathrm{\hspace{0.17em}\hspace{0.17em}2}^{\frac{1}{4}}G_F^{\frac{1}{2}}H\varphi _h^2\frac{2M^2}{27};$$
(13)
this relation is valid in the chiral limit of massless light quarks (see Eq. (10)); $`M_T`$ is the mass of the heavy quark and $`M`$ is the hadron mass. We put the number of light quarks $`N_F=3`$ and the number of colors $`N_c=3`$; $`\varphi _h`$ and $`H`$ are hadron and Higgs operators. Note that, as a consequence of scale anomaly, Eq. (13) does not have an explicit dependence on the coupling $`\alpha _s`$.
## 3 Estimates for double Pomeron Higgs production cross sections
### 3.1 General formulae for double Pomeron Higgs production
The amplitude for Higgs production in the Pomeron approach is given by ( see for example Refs.)
$$M(h+hh+H+h)=g_1(t_1)g_2(t_2)\gamma (t_1,t_2)\eta _+(t_2)\eta _+(t_1)\left(\frac{s}{s_2}\right)^{\alpha _P(t_2)}\left(\frac{s}{s_1}\right)^{\alpha _P(t_1)},$$
(14)
where $`s_1=(P_1+q_1)^2`$ and $`s_2=(P_2+q_1)^2`$ ( $`P_{1,2}`$ are momenta of incoming hadrons ); $`\eta _+(t_i)`$ is a signature factor, which for the Pomeron is
$$\eta _+(t_i)=i+tan^1\left(\frac{\pi \alpha _P(t_i)}{2}\right),$$
(15)
where $`\alpha _P(t)`$ is the Pomeron trajectory, $`\alpha _P(t)=1+\mathrm{\Delta }_P+\alpha _P^{}t`$, with $`\mathrm{\Delta }_P0.08`$ ; all other notations are evident from Fig.2.
The cross section for DPHP in the central rapidity region ($`y_H=0`$, where $`y_h`$ is the rapidity of the produced Higgs boson) can be written down as
$$\frac{d\sigma }{dy_Hdt_1dt_2}|_{y_H=0}=$$
(16)
$$\frac{1}{2s}|M(h+hh+H+h)|^2\underset{i=1,2}{}\frac{d^3P_i^{}}{(2\pi )^32P_{i,0}^{}}\frac{d^2p_{H,}}{2(2\pi )^3}(2\pi )^4\delta ^{(4)}(P_1+P_2P_1^{}P_2^{}p_H)$$
where $`P_i^{}`$ are momenta of recoil hadrons, while $`p_H`$ is the momentum of the produced Higgs boson.
Performing all integrations and recalling that $`s_1s_2=M_H^2s`$ we obtain
$$\frac{d\sigma }{dy_Hdt_1dt_2}|_{y_H=0}=\frac{2g_1^2(t_1)g_2^2(t_2)\gamma ^2(t_1,t_2)}{\pi (16\pi )^2}\left(\frac{s}{M_H^2}\right)^{2\mathrm{\Delta }_P}e^{\alpha _P^{}\mathrm{ln}(s/M_H^2)[t_1+t_2]}.$$
(17)
We will assume that $`\gamma (t_1,t_2)`$ is a smooth function of $`t_1`$ and $`t_2`$ in comparison with $`g_1(t_1)`$ and $`g_2(t_2)`$. Indeed, the t-dependence of $`g_i`$ is related to the quark distribution inside the hadron while the t-dependence of $`\gamma `$ is determined by the mean transverse of gluon inside the Pomeron. The typical scale for this momentum is $`1/\alpha _P\mathrm{\hspace{0.17em}4}GeV^2`$ which is much larger than the typical momentum of a quark in a hadron ( $`\mathrm{\hspace{0.17em}\hspace{0.17em}0.1}GeV^2`$ ).
Using this assumption together with the simplest Gaussian parameterization for the vertex $`g_i(t_i)=g_i(0)exp(R_0^2|t_i|)`$ we obtain
$$\frac{d\sigma }{dy_H}|_{y_H=0}=\frac{8g_1^2(0)g_2^2(0)}{\pi [16\pi B_{el}(s/M_H^2)]^2}\gamma ^2(t_1=0,t_2=0)\left(\frac{s}{M_H^2}\right)^{2\mathrm{\Delta }_P}.$$
(18)
Recalling now the well–known relation between the total and elastic cross sections for the one Pomeron exchange, namely,
$$R_{el}(s)=\frac{\sigma _{el}(s)}{\sigma _{tot}(s)}=\frac{g_1(0)g_2(0)}{16\pi B_{el}(s)}\left(\frac{s}{s_0}\right)^{\mathrm{\Delta }_P};$$
(19)
where $`B_{el}=4R_0^2+2\alpha _P^{}\mathrm{ln}s`$, one can derive
$$\frac{d\sigma }{dy_H}|_{y_H=0}=\gamma ^2(t_1=0,t_2=0)\times \frac{8}{\pi }\times R_{el}^2\left(\frac{s}{M_H^2}s_0\right).$$
(20)
There is only one unknown factor in Eq. (20), namely, $`\gamma ^2(t_1=0,t_2=0)`$. In the next subsection we present the estimates for this factor using the non-perturbative approach that has been discussed in the section 2.
### 3.2 The production vertex $`\gamma (t_1=0,t_2=0)`$
Our estimate of $`\gamma (t_1=0,t_2=0)`$ consists of two steps:
1. For positive values of $`t_1=t_2=m_{glueball}^2`$ we can obtain $`\gamma (t_1=m_{glueball}^2,t_2=m_{glueball}^2)`$ from Eq. (13);
2. Using Eq. (14) we can make the analytic continuation to the region $`t_1<0`$ and $`t_2<0`$, which corresponds to the scattering process.
We will assume that there exists a tensor $`2^{++}`$ glueball which lies on the Pomeron trajectory, namely, that its mass satisfies the following relation:
$$\alpha _P(t=m_{glueball}^2)=\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}1}+\mathrm{\Delta }+\alpha _P^{}(0)m_{glueball}^2=\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}2}.$$
(21)
There is no undisputed experimental evidence for such a meson but lattice calculations give for its mass $`m_{glueball}=2.4`$ GeV . This mass is a little bit higher than can be expected from Eq. (21) with the experimental $`\alpha _P^{}(0)=0.25GeV^2`$ . On the other hand, it is possible to describe experimental data using a smaller value of $`\alpha _P^{}(0)0.17GeV^2`$ which is needed to satisfy Eq. (21) with $`m_{glueball}=2.4GeV`$, assuming the presence of substantial shadowing corrections .
For the diagram in Fig.4 the vertex $`\gamma _{glueball}`$ can be easily evaluated from Eq. (13); it is equal to
$$\gamma _{glueball}=\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}2}^{\frac{1}{4}}G_F^{\frac{1}{2}}\frac{2m_{glueball}^2}{27}.$$
(22)
One can see that Eq. (14) leads to the contribution described by Fig.4. Indeed, for $`t_im_{glueball}^2`$
$$\eta _+(t_i)\frac{2}{\pi \alpha _P^{}\left(m_{glueball}^2t_i\right)}.$$
(23)
(A more detailed discussion of the analytic properties of the Reggeon exchange can be found in Ref. ). Using Eq. (23) and comparing Eq. (14) with the diagram of Fig.4, we conclude that
$$\gamma (t_1=m_{glueball}^2,t_2=m_{glueball}^2)=\frac{\pi }{2}\alpha _P^{}(0)\gamma _{glueball}.$$
(24)
The reggeon approach cannot tell us anything on the relation between $`\gamma (t_1=m_{glueball}^2,t_2=m_{glueball}^2`$ and $`\gamma (t_1=0,t_2=0)`$. The only thing that we can claim is that the signature factor takes into account the steepest part of $`t`$-behavior. Therefore, in the next subsection we will assume that
$$\gamma (t_1=m_{glueball}^2,t_2=m_{glueball}^2)=\gamma (t_1=0,t_2=0);$$
(25)
this is an extreme assumption which can be used to obtain an upper bound on the cross section. Uncertainties related to this and other assumptions we make will be discussed in detail in section 3.4 and in the summary, section 5.
### 3.3 The magnitude of the cross section
Using Eq. (22),Eq. (24) and Eq. (25) we can rewrite Eq. (20) in the simple form
$$\frac{d\sigma }{dy_H}|_{y_H=0}=\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}2}\pi \left(\alpha _P^{}m_{glueball}^2\right)^2\times \frac{4\sqrt{2}G_F}{27^2}\times R^2\left(\frac{s}{M_H^2}s_0\right).$$
(26)
For $`M_H=100GeV`$, the factor $`S/M_H^2s_0`$ is equal to $`400GeV^2`$ for $`s_0=1GeV^2`$. Therefore, we can take $`R_{el}0.175`$ ( see Fig.5 ) for the Tevatron energies. Eq. (26) leads to
$$\frac{d\sigma }{dy_H}|_{y_H=0}(M_h=100GeV,\sqrt{s}=1800GeV)=\mathrm{\hspace{0.17em}\hspace{0.17em}6.4}\mathrm{pbarn},$$
(27)
This is a very large number, especially if we recall that the total inclusive cross section for Higgs production in perturbation theory is on the order of $`1`$ pbarn . However this estimate does not yet contain the suppression due to the (small) probability of the rapidity gap survival, which will be discussed in section 4, where we present our final results. Since $`R_{el}`$ grows with energy, $`R_{el}s^\mathrm{\Delta }`$ we expect that the cross section at the LHC energy is approximately 2 times larger than the one in Eq. (27).
### 3.4 Uncertainties of our estimates
1. Let us start with the value of $`R_{el}`$. We took it from the experimental data, but we nevertheless have two uncertainties associated with it. First, Eq. (19) is written for one Pomeron exchange while in experimental data at $`\sqrt{s}\mathrm{\hspace{0.17em}20}GeV`$ we have about 30% contamination from the secondary Reggeons . If we try to extract the one Pomeron exchange from the data, it reduces the value of cross section for DPHP by 1.7 times. Therefore, the value for the cross section can be about $`3.8pbarn`$ rather than Eq. (27). The second uncertainty in evaluation of $`R_{el}`$ is the value of $`s_0`$; even though $`s_0=1GeV`$ appears in all phenomenological approaches , we have no theoretical argument for the value of $`s_0`$. However, since the ratio $`R_{el}`$ in Fig.5 is a rather smooth function of energy we do not expect that the uncertainty in the value of $`s_0`$ can introduce a large error.
2. We can take into account also the reactions of Eq. (2) and Eq. (3). In Eq. (26) we would then have to substitute
$$R_{el}R_D=R_{el}+\frac{\sigma ^{DD}(s)}{\sigma _{tot}},$$
(28)
where $`\sigma ^{DD}`$ is the cross section of the double diffraction dissociation. Unfortunately, we do not have conclusive data on this cross section. However, recent CDF measurements show that this cross section could be rather large ( about 4.7 mb at the Tevatron energy).
3. The principle uncertainty, however, is associated with the continuation from $`t=m_{glueball}^2`$ to $`t=0`$. This is a question which at present can only be addressed in the framework of different models. For example, in Veneziano model instead of $`\eta _+(t)`$ (see Eq. (15) ) a new factor appears, namely
$$\eta _+^V(t_i)=\mathrm{\Gamma }(2\alpha _P(t_i))e^{i\frac{\pi \alpha _P(t_i)}{2}}.$$
(29)
Eq. (29) does not give the factor of $`\pi /2`$ in Eq. (24) and, therefore, decreases the value of the cross section given by Eq. (26) by a factor of 2.5. We will return to the discussion of the analytic continuation in the summary section.
4. As we have discussed in section 2, we can evaluate the value of $`\gamma (t_1=m_{glueball}^2,t_2=m_{glueball}^2)=\gamma _{glueball}`$ only if $`M_H/M_T<1`$. The accuracy of Eq. (22) is $`O(M_H^2/M_T^2)`$ and we thus believe that Eq. (22) gives a reasonable estimate of $`\gamma _{glueball}`$ for Higgs meson with $`M_H100GeV`$.
## 4 Survival of large rapidity gaps
As has been discussed intensively during the past decade (see Refs.), the cross section of Eq. (26) has to be multiplied by a factor $`S_{spect}^2`$, which is the survival probability of large rapidity gap (LRG) processes. The “experimental” cross section is therefore given by
$$\frac{d\sigma (ppppH)}{dy}|_{y=0}=S_{spect}^2\frac{d\sigma _P(ppppH)}{dy}|_{y=0}.$$
(30)
Here, $`d\sigma _P(ppppH)/dy`$ denotes the cross section calculated in Eq. (26). The factor $`S_{spect}^2`$ has a very simple meaning – it is a probability of the absence of inelastic interactions of the spectators which could produce hadrons inside the LRG. We have rather poor theoretical control of the value of the survival probability; this fact reflects the lack of knowledge of the “soft” physics stemming from non-perturbative QCD. Different models exist, leading to the values about $`S_{spect}^2\mathrm{\hspace{0.17em}\hspace{0.17em}10}^1`$ at the Tevatron energies. For double Pomeron processes, this quantity has been discussed in Ref . The result of this analysis is that the value of the survival probability for double Pomeron production is almost the same as for “hard” dijet production with LRG between them. Fortunately, the value of $`S_{spect}^2`$ has been measured , and is equal to 0.07 for the highest Tevatron energy.
Multiplying Eq. (27) by $`S_{spect}^2`$ = 0.07 and taking into account suppression due to the factor of Eq. (29), we obtain
$$\frac{d\sigma }{dy_H}|_{y_H=0}(M_h=100GeV,\sqrt{s}=1800GeV)=\mathrm{\hspace{0.17em}\hspace{0.17em}0.2}\mathrm{pbarn}.$$
(31)
This estimate is not our final result yet, since we still have to correct it by the additional suppression factor $`S_{par}^2`$ which describes the probability of the absence of the parasite gluon emission around the Higgs production vertex (see Fig.4-b) . As was argued in Ref.,
$$S_{par}^2=e^{<N_G(\mathrm{\Delta }y=\mathrm{ln}(M_H^2/s_0))>},$$
(32)
with
$$<N_G(\mathrm{\Delta }y=\mathrm{ln}(M_H^2/s_0))>=\frac{N_{hadrons}(\mathrm{\Delta }y=\mathrm{ln}(M_H^2/s_0))}{N_{hadrons}(oneminijet)}\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}2}÷\mathrm{\hspace{0.17em}4}.$$
(33)
It gives $`S_{par}^2=0.14÷\mathrm{\hspace{0.17em}0.014}`$.
The appearance of this factor can be illustrated by the following argument: one of the most important differences between the diagrams of Fig.2 and in Fig.4 is the fact that the Pomeron exchange is almost purely imaginary while the glueball exchange leads to the real amplitude. Imaginary amplitude describes the production of particles and the Pomeron is associated with the inelastic process with large multiplicity. Therefore, normally, in a large rapidity region $`\mathrm{\Delta }y=\mathrm{ln}(M_H^2/s_0)`$ we expect to see a large number of produced particles while in Fig.2 we require that only one Higgs boson is produced. Therefore, it seems reasonable to expect a suppression for the double–diffractive Higgs production, and this suppression can be described by Eq. (32) and Eq. (33).
Finally, for the Tevatron energy we expect
$`{\displaystyle \frac{d\sigma (ppppH)}{dy}}|_{y=0}\left(\sqrt{s}=\mathrm{\hspace{0.17em}\hspace{0.17em}1.8}TeV\right)`$ $`=`$ $`S_{spect}^2\times S_{par}^2\times {\displaystyle \frac{d\sigma _P(ppppH)}{dy}}|_{y=0}`$
$`=`$ $`\left(\mathrm{\hspace{0.17em}0.0038}÷\mathrm{\hspace{0.17em}0.028}\right)\mathrm{pbarn}`$ (34)
Extrapolating to the LHC energy, we have two effects that work in different directions: the rise of the Pomeron contribution and the decrease of the $`S_{spect}^2`$ with energy. From Ref. we expect that $`S_{spect}^2(\sqrt{s}=\mathrm{\hspace{0.17em}\hspace{0.17em}16}TeV)/S_{spect}^2(\sqrt{s}=\mathrm{\hspace{0.17em}\hspace{0.17em}1.8}TeV)\mathrm{\hspace{0.17em}\hspace{0.17em}0.7}`$ while the rise of the Pomeron exchange leads to an extra factor of 2 in Eq. (34). Therefore, our final estimate for the LHC is
$`{\displaystyle \frac{d\sigma (ppppH)}{dy}}|_{y=0}\left(\sqrt{s}=\mathrm{\hspace{0.17em}16}TeV\right)`$ $`=`$ $`S_{spect}^2\times S_{par}^2\times {\displaystyle \frac{d\sigma _P(ppppH)}{dy}}|_{y=0}`$ (35)
$`=`$ $`\left(\mathrm{\hspace{0.17em}0.0015}÷\mathrm{\hspace{0.17em}0.042}\right)\mathrm{pbarn}.`$
Eq. (34) and Eq. (35) give significantly larger (by about $`5`$ times) larger cross sections than expected for double–diffractive production in pQCD . However, Ref. contains an estimate of the upper bound on double Pomeron Higgs production in pQCD obtained by choosing the largest possible value for $`S_{par}^2`$ ( see Eq. (33) ). This upper bound appears to be about 7 times larger than the highest value in Eq. (34).
## 5 Summary and discussion
The approach suggested in this paper is based entirely on non-perturbative QCD. We believe that such an approach is logically justified for diffractive Higgs production since even pQCD calculations show that this is, to large extent, a “soft” process (see Eq. (5) and the following discussion). However, just because of this, we have to stress again that the accuracy of our calculation is not very good. We feel, however, that our results support the idea that in pQCD approach to diffractive Higgs production the running QCD coupling has to be taken at the “soft” scale $`Q^21\mathrm{GeV}^2`$. As was argued in Ref., in BLM prescription of taking into account the running QCD coupling one can insert the quark bubbles only in the $`t`$-channel gluon lines in Fig. 3. Therefore, the running QCD coupling depends on the transverse momenta of these gluons, and they are determined by the “soft” scale<sup>1</sup><sup>1</sup>1 In this soft regime, the dependence on the coupling constant in the Pomeron can disappear as a consequence of scale anomaly .. The Eq. (13) indeed does not depend on the QCD coupling, demonstrating the non-perturbative, “soft” character of the discussed process.
We obtain quite large values for the cross section of the diffractive Higgs production – after integration over the Higgs rapidity $`y`$ in Eq. (34) and Eq. (35) we get
$$\sigma (ppppH)\left(\sqrt{s}=\mathrm{\hspace{0.17em}\hspace{0.17em}1.8}TeV\right)=\mathrm{\hspace{0.17em}\hspace{0.17em}0.019}÷\mathrm{\hspace{0.17em}0.14}\mathrm{pbarn}.$$
(36)
and
$$\sigma (ppppH)\left(\sqrt{s}=\mathrm{\hspace{0.17em}\hspace{0.17em}16}TeV\right)=\mathrm{\hspace{0.17em}\hspace{0.17em}0.01}÷\mathrm{\hspace{0.17em}0.27}\mathrm{pbarn}.$$
(37)
Comparing our estimates with the ones based on pQCD we conclude that the lowest of our values of the cross section of double Pomeron Higgs production is about the same as the highest one in the pQCD approach. However, both our approach and the pQCD one are suffering from large uncertainties, stemming from the analytical continuation in our approach and from the survival probability of rapidity gap and the absence of “parasite emission” $`S_{par}^2`$ in pQCD.
Let us point out that Eq. (36) shows that the double Pomeron Higgs production constitutes a substantial part of the total inclusive Higgs production. Moreover, our calculations lead to an additional contribution to the inclusive cross section which is shown in Fig.6.
(Note that the triple Pomeron interaction gives a very small contribution to the process in Fig. 6 due to the small real part in the Pomeron exchange .) Using the same approach as in derivation of Eq. (17) we obtain
$$\frac{d\sigma _{incl}(ppH+X)}{dy}|_{y_H=0}=\gamma _R^2(t_1=0,t_2=0)\frac{2g_1(0)g_2(0)(G_{RR}^P)^2}{\pi (16\pi B_R)^2}\frac{1}{\mathrm{\Delta }_R^2}\left(\frac{s}{M_H^2}\right)^{\mathrm{\Delta }_P},$$
(38)
where $`\mathrm{\Delta }_R0.5`$. As a first approximation we can take ( see Eq. (24) )
$$\gamma _R(t_1=0,t_2=0)=\frac{\pi }{2}\mathrm{\hspace{0.17em}2}^{\frac{1}{4}}G^{\frac{1}{2}}\frac{2\alpha _R^{}m_f^2}{27}$$
(39)
where $`m_f`$ is the mass of the $`f`$ \- meson which is the first resonance on the secondary Reggeon trajectory, and $`\alpha _R^{}m_f^2=1.5`$. Substituting Eq. (39) in Eq. (38) we obtain
$$\frac{d\sigma _{incl}(ppH+X)}{dy}|_{y_H=0}=\frac{(G_{RR}^P)^2}{8B_R}\times \frac{B_R}{B_{el}}\times \mathrm{\hspace{0.17em}2}^{\frac{1}{2}}G_F\frac{1}{3}\times R\left(\frac{s}{M_H^2}s_0\right).$$
(40)
Eq. (40) gives
$$\frac{d\sigma _{incl}(ppH+X)}{dy}|_{y_H=0}=\left(\frac{G_{RR}^P}{g}\right)^2\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}3.4\hspace{0.17em}10}^7mb$$
(41)
which does not yet contain the suppression arising from the analytical continuation. We take this suppression into account by multiplying Eq. (41) by factor $`S_{par}^2=0.140.014`$. Unfortunately, we do not know the value for the ratio $`G_{RR}^P/g`$. In the triple Pomeron parameterization of the cross section of diffractive dissociation in hadron reactions this ratio changes from 1 to 0. For $`G_{RR}^P/g=1`$ we get for the “soft” inclusive cross section the value of $`43÷430\mathrm{pbarn}`$. On the other hand, taking Field and Fox value for this ratio we obtain a much smaller, but still very sizeable value of $`0.43÷\mathrm{\hspace{0.17em}4.28}\mathrm{pbarn}`$. It is thus clear that the evaluation of the “soft” contribution to the inclusive Higgs production is plagued by large uncertainties; however, it might be bigger than the pQCD one .
We hope that this paper will help to look at diffractive Higgs production from a different viewpoint, and will stimulate a much needed further work. To our surprise, despite the very different non–perturbative method used here, our estimates for the double–diffractive production turn out to be not that far from the pQCD calculation ( the average is about $`5`$ times larger ). It adds some confidence in both approaches and gives us a hope that one will be able to perform a reliable calculation in the nearest future.
Acknowledgments: We are very grateful to Mike Albrow, Andrew Brandt, Al Mueller and all participants of the working group “Diffraction physics and color coherence” at QCD and Weak Boson Physics Workshop in preparation for Run II at the Fermilab Tevatron for stimulating discussions of Higgs production and encouraging criticism. We thank Stan Brodsky, Asher Gotsman, Valery Khoze, Larry McLerran, Uri Maor and Misha Ryskin for fruitful discussions on the subject. E.L. thanks the Fermilab theory department for creative atmosphere and hospitality during his stay when this paper was being written.
The work of D.K. was supported by the US Department of Energy (Contract # DE-AC02-98CH10886) and RIKEN. The research of E.L. was supported in part by the Israel Science Foundation, founded by the Israeli Academy of Science and Humanities, and BSF # 9800276.
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# 1 Introduction
## 1 Introduction
Non-commutative field theory is a model of a world with non-commuting spatial coordinates (space/space non–commutativity) and in fact such field theories do arise as the description of string theory in certain backgrounds . Non–commutative theories are very nonlocal in the non-commuting spatial directions but are quadratic in time derivatives. Nonlocality in spatial directions ruins Lorentz invariance but it is consistent with the basic rules of Hamiltonian quantum mechanics. Action at a distance may occur but events never precede their causes.
The situation is much less clear for field theories with non–commutativity between time and a space direction (space/time non–commutativity). The action is arbitrarily non-local in time with the evolution of fields at one time depending on the value of fields at both past and future times. The question then is whether the kind of unusual behavior found in space/time non–commutative field theories can ever occur in any consistent theory with a Hamiltonian and a unitary S-Matrix. We don’t know the answer to this question but we examine an obvious candidate, string theory in a background electric $`B_{\mu \nu }`$ field. This theory is manifestly unitary and may be expected to exhibit effects similar to those seen in the field theory example. However as we have seen in the theory in an electric field never becomes a field theory and retains its stringy excitations. We will see that the stringy effects cancel the acausal effects of space/time non–commutativity.
The plan of the paper is as follows. In section 2 we describe the scattering of wave packets in non–commutative field theory. In the case of space/space non–commutativity the scattering induces a sudden spatial displacement of the wave packet in the direction orthogonal to the momentum. The magnitude of displacement is proportional to the momentum. As expected the degree of nonlocality increases with momentum. The scattered particles behave like rigid rods oriented perpendicular to their momentum with a size proportional to their momentum .
In the space/time case the wave packets scatter in a manner that appears to violate causality. Two scattered packets appear. One of them is physically sensible and corresponds to a time delay proportional to the incoming momentum. The other wave packet has a negative time delay! As in the space/space case the effect increases with the momentum. Thus at very high energy one of the outgoing waves appears to originate long before the particles could have collided. We refer to this behavior as advanced. Alternatively the particles behave like rods oriented along the direction of motion. Again the length of the rod is proportional to its momentum. This increase of length with momentum is very counterintuitive and is quite opposite to the expected Lorentz contraction.
Having defined the distinctive signatures of space/time non–commutativity, in Section 3 we proceed to look for string theory realizations of these signatures. We investigate open string scattering with and without background electric fields. We find delayed wave packets with time delay proportional to momentum as expected in non–commutative theories. We emphasize that the delayed effect occurs with and without a background electric field. However in no case do we find the acausal signatures of space-time non–commutativity. The case of open strings in electric fields is particularly interesting. Although the amplitudes acquire Moyal phases the stringy effects mask the phases that would otherwise give rise to advanced effects. The scattering with the electric field does not seem appreciably different than without it even in the critical limit.
## 2 Scattering in Non–Commutative Field Theory
In this section we study the effect of space/space and space/time non–commutativity on the scattering of massless scalar particles. We begin with the space/space case. To illustrate the main points it is sufficient to consider 2+1 dimensional non–commutative scalar $`\varphi ^4`$ theory in lowest order perturbation theory. The coordinates are labeled $`(x,y,t)`$. Since this case is familiar we will just describe the scattering schematically. Let us consider two high energy particles moving along the $`x`$ axis with spatial momentum $`P_x`$. We will take the initial wave function to be
$$\mathrm{\Psi }(x,y)=\mathrm{exp}(iP_xx)\psi _{in}(y)$$
(2.1)
where
$$\psi _{in}(y)=𝑑P_y\widehat{\psi }_{in}(P_y)e^{iP_yy}.$$
(2.2)
The important feature of the scattering amplitude for our purposes is the Moyal phase factors which take the form
$$M=\mathrm{exp}\frac{i}{2}\theta (P_xQ_yP_yQ_x)$$
(2.3)
where $`Q`$ is the momentum transfer. We assume $`P_xP_y,Q_x,Q_y`$.
After the scattering, the scattered momentum space wave function is given by an expression of the form
$$\widehat{\psi }_{out}(P_y)=𝑑Q_y\widehat{\psi }_{in}(P_y+Q_y)\mathrm{exp}\frac{i}{2}\theta (P_xQ_y).$$
(2.4)
In coordinate space
$$\psi _{out}(y)\psi _{in}(y)\delta (y\frac{1}{2}\theta P_x).$$
(2.5)
In other words the outgoing scattered wave appears to originate from the displaced position $`y=\theta P_x/2`$.
An intuitive way to understand this effect is to think of the incident particles as extended rods oriented perpendicular to their momentum . The size of the rods is $`\theta P`$ and the rule is that they only interact if their ends touch.
Now we turn to the more interesting case of space/time non–commutativity which we will study in much more detail. For simplicity we will work in $`1+1`$ dimensions. We denote time by $`t`$ and the spatial variable by $`x`$.
Let us begin by reviewing the scattering of wave packets in 1+1 dimensions. A free scalar field in $`1+1`$ dimensions has the following Fourier decomposition
$$\varphi (x,0)=\frac{dp}{(2\pi )\sqrt{2E_p}}\left(a_pe^{ipx}+a_p^{}e^{ipx}\right),$$
(2.6)
with
$$[a_p,a_k^{}]=(2\pi )\delta (pk).$$
(2.7)
Because of the special infrared divergences of massless 1+1 dimensional scalar fields we will work with the derivative of $`\varphi `$ rather than $`\varphi `$ itself:
$$\varphi ^{}(x,0)=i\frac{dp}{(2\pi )\sqrt{2E_p}}\left(pa_pe^{ipx}pa_p^{}e^{ipx}\right).$$
(2.8)
Single particle states with momentum $`p`$ are normalized as follows
$$|p>=\sqrt{2E_p}a_p^{}|0>.$$
(2.9)
Then the norm
$$<p|k>=2E_p(2\pi )\delta (pk)$$
(2.10)
is Lorentz invariant. The wavefunction of such a state will be defined by
$$<0|\varphi ^{}(x)|p>=ipe^{ipx}.$$
(2.11)
Using the equation of motion for the free scalar field, we can find the wavefunction at all times.
Next, we turn on some interactions. For example, consider a commutative $`\varphi ^4`$ interaction. We are interested in the scattering of massless scalars, in particular 2–body to 2–body scattering. For sufficiently high energies, we can use perturbation theory to calculate an S–matrix. The S–matrix takes the following form
$$S=1+iT,$$
(2.12)
where
$$<p_1,p_2|iT|k_1,k_2>=(2\pi )^2\delta ^2(k_1+k_2p_1p_2)i(k_1,k_2p_1,p_2).$$
(2.13)
Here, $`k_1,k_2`$ denote the 2–momenta of the incoming particles and $`p_1,p_2`$ the 2–momenta of the outgoing particles. The invariant amplitude $`i`$ is computed in the usual way using Feynman diagrams. For the simple case of a $`\varphi ^4`$ interaction,
$$i=ig$$
(2.14)
to leading order in perturbation theory. In $`1+1`$ dimensions the only effect one expects to see in 2–body to 2–body scattering is time delays.
Now, consider an incoming state consisting of correlated pairs of particles with opposite momenta:
$$|\varphi >_{in}=\frac{dk}{(2\pi )2E_k}\varphi _{in}(k)|k,k>,$$
(2.15)
with
$$\varphi _{in}(k)=\varphi _{in}(k).$$
(2.16)
The wavefunction of such a state is given by
$$\mathrm{\Phi }_{in}(x)<0|\varphi ^{}(x_1)\varphi ^{}(x_2)|\varphi >_{in}=2\frac{dkk^2}{(2\pi )2E_k}\varphi _{in}(k)e^{ikx},$$
(2.17)
where $`x=x_1x_2`$ is the relative separation of the two particles. There is no dependence on the center of mass position, since the overall center of mass momentum is zero. Let us also choose $`\varphi _{in}(k)`$ so that at the time of the collision $`t=0`$, the wave-packet is well concentrated at $`x=0`$. Then the incoming particles are close together at $`t=0`$. For example, we may choose
$$\varphi _{in}(k)=e^{\frac{(kk_0)^2}{\lambda }}+e^{\frac{(k+k_0)^2}{\lambda }}.$$
(2.18)
The wavepacket is concentrated at energies closed to $`k_0`$. The width of the packet in space is given by $`1/\lambda ^{1/2}`$. We let $`\lambda k_0^2`$ and take $`k_0`$ large. At earlier times, $`t<0`$, we can use the free equations of motion to find that the packet is concentrated at $`x=2t`$. This means that the incoming particles are far apart in the past and they collide at $`t=0`$.
Similarly, the outgoing state is taken to be
$$|\varphi >_{out}=\frac{dp}{(2\pi )2E_p}\varphi _{out}(p)|p,p>.$$
(2.19)
Then,
$$|\varphi >_{out}=S|\varphi >_{in}=|\varphi >_{in}+iT|\varphi >_{in}.$$
(2.20)
Therefore, we have that
$$\frac{\varphi _{out}(p)}{(2\pi )2E_p}=\frac{\varphi _{in}(p)}{(2\pi )2E_p}+\frac{<p,p|iT|\varphi >_{in}}{8(2\pi )^2E_p^2\delta (0)}.$$
(2.21)
Now, using the form of the matrix element $`<p,p|iT|k,k>`$, eq. (2.13), we find that the non–trivial part of $`\varphi _{out}(p)`$ is given by
$$\frac{dk}{(2\pi )2E_k}\varphi _{in}(k)(\frac{i}{8E_p^2})\delta (2E_k2E_p)=\frac{\varphi _{in}(p)}{(2\pi )2E_p}\frac{i}{8E_p^2}.$$
(2.22)
Therefore, using eq. (2.17), the non–trivial part of the outgoing wavefunction can be obtained by
$$\mathrm{\Phi }_{out}(x)<0|\varphi ^{}(x_1)\varphi ^{}(x_2)|\varphi >_{out}=2\frac{dpp^2}{(2\pi )2E_p}\varphi _{in}(p)\frac{i}{8E_p^2}e^{ipx}.$$
(2.23)
In the case of the $`\varphi ^4`$ theory, we see that nothing much happens. Choose $`\varphi _{in}(p)`$ to be a polynomial in $`p`$ times a Gaussian so that the integral converges. Then $`\mathrm{\Phi }_{in}(x)`$ is concentrated at $`x=0`$. Since $`ig`$, at time $`t=0`$, the outgoing wavefunction will also be concentrated at $`x=0`$. Therefore, there are no large time delays. Using the free equations of motion, we find that at later times the wave-packet is concentrated at $`x=2t`$ and so the outgoing particles separate in the far future.
Consider now the effect of space/time non–commutativity.
$$[t,x]=i\theta .$$
(2.24)
The theory is defined by replacing the ordinary product by a $``$–product given by
$$\varphi _1\varphi _2(x,t)=e^{i\frac{\theta }{2}[_0^y_1^z_1^y_0^z]}\varphi (y)\varphi (z)|_{y=z=(x,t)}.$$
(2.25)
The $`\varphi ^4`$ Lagrangian contains now an infinite number of time derivatives from the interaction term
$$g\varphi \varphi \varphi \varphi .$$
(2.26)
Therefore the theory is not local in time. It is not clear that such a theory has a well defined Hamiltonian. One plus one dimensional Lorentz invariance however, is undisturbed by the non–commutativity. This is easily seen from the fact that the defining commutation relation has the form
$$[x^\mu ,x^\nu ]=i\theta ϵ^{\mu \nu }.$$
(2.27)
The effect of the $``$–product is to produce phases in the interaction vertex that depend on the energies of the particles. The tree–level scattering amplitude is now given by
$$ig[\mathrm{cos}(p_1p_2)\mathrm{cos}(p_3p_4)+23+24],$$
(2.28)
where $`p_1,p_2,p_3,p_4`$ are the 2–momenta of the particles satisfying
$$p_1+p_2+p_3+p_4=0.$$
(2.29)
Here $`pk=\theta (p^0k^1k^0p^1)`$. Note that we have used conventions with all particles taken to be incoming in the vertices; i.e. energies of outgoing particles are negative. In the center of mass frame with the incoming particles (and outgoing) having equal and opposite spatial momenta, the amplitude becomes
$$ig[\mathrm{cos}(4p^2\theta )+2].$$
(2.30)
The pattern is similar in more general non–commutative theories but depending on the spins and polarizations of the particles the periodic functions may be sines in place of cosines.
We remark that such a theory fails to be unitary at the 1–loop level . However, let us just consider tree–level scattering amplitudes and in particular the effect of non–commutativity on the outgoing wave-packets. We choose for the incoming wave-packet $`\varphi _{in}(p)`$ a gaussian function:
$$\varphi _{in}(p)E_p\left(e^{\frac{(pp_0)^2}{\lambda }}+e^{\frac{(p+p_0)^2}{\lambda }}\right).$$
(2.31)
(The extra factor of $`E_p`$ is added to simplify the integrals but does not change the qualitative behavior of our results.) Using eq. (2.19), we can find the outgoing wavefunction
$$\mathrm{\Phi }_{out}(x)g𝑑p[\mathrm{cos}(4p^2\theta )+2]\left(e^{\frac{(pp_0)^2}{\lambda }}+e^{\frac{(p+p_0)^2}{\lambda }}\right)e^{ipx}.$$
(2.32)
To compute the integral, we need to calculate the following Fourier transform
$`{\displaystyle 𝑑pe^{4ip^2\theta }e^{\frac{(pp_0)^2}{\lambda }}e^{ipx}}e^{\frac{p_0^2}{\lambda }}{\displaystyle \frac{1}{\sqrt{\frac{1}{\lambda }4i\theta }}}e^{(\frac{2p_0}{\lambda }+ix)^2\frac{1}{4(\frac{1}{\lambda }4i\theta )}}`$ (2.33)
$`=`$ $`{\displaystyle \frac{1}{\sqrt{\frac{1}{\lambda }4i\theta }}}\mathrm{exp}\left[{\displaystyle \frac{\lambda \left(x+8p_0\theta \right)^2}{4(1+16\theta ^2\lambda ^2)}}\right]\mathrm{exp}\left[{\displaystyle \frac{i\theta \lambda ^2\left(x\frac{p_0}{2\lambda ^2\theta }\right)^2}{1+16\theta ^2\lambda ^2}}\right]\mathrm{exp}i{\displaystyle \frac{p_0^2}{4\lambda ^2\theta }}.`$ (2.34)
We take $`p_0\lambda ^{1/2}1/p_0\theta `$, and also assume that $`\lambda \theta 1`$. Then, eq. (2.33) simplifies as follows
$$\frac{1}{\sqrt{4i\theta }}e^{\frac{(x+8p_0\theta )^2}{64\theta ^2\lambda }}e^{i\frac{(x\frac{p_0}{2\lambda ^2\theta })^2}{16\theta }}e^{i\frac{p_0^2}{4\lambda ^2\theta }}F(x;\theta ,\lambda ,p_0).$$
(2.35)
Then the outgoing wavefunction is given by
$$\mathrm{\Phi }_{out}(x)g\left[F(x;\theta ,\lambda ,p_0)+4\sqrt{\lambda }e^{\lambda \frac{x^2}{4}}e^{ip_0x}+F(x;\theta ,\lambda ,p_0)\right]+(p_0p_0).$$
(2.36)
We see that the wave-packet splits into three parts, one concentrated at $`x=8p_0\theta `$, one at $`x=0`$ and the other at $`x=8p_0\theta `$. The width of the first and third packet is given by $`8\lambda ^{1/2}\theta `$ while the one concentrated at $`x=0`$ has width $`2/\lambda ^{1/2}`$. Therefore, the packets are well separated for $`p_0\lambda ^{1/2}1/p_0\theta `$. The separation of the two displaced packets is proportional to $`p_0`$ which is the energy of the particles. The bigger the energy is the bigger the separation.
The packet at $`x=0`$ oscillates with frequency $`p_0`$. The other two packets oscillate with phases $`\mathrm{exp}[i(x+p_0/2\lambda ^2\theta )^2/16\theta ]`$ and $`\mathrm{exp}[i(xp_0/2\lambda ^2\theta )^2/16\theta ]`$. Locally, near the maxima at $`x=\pm 8p_0\theta `$ the phases in the other two packets become $`\mathrm{exp}[ip_0(1+1/16\lambda ^2\theta ^2)\mathrm{\Delta }x]`$ and so they oscillate with frequency $`p_0`$ since $`\lambda \theta 1`$. This was expected from energy conservation.
All three wavepackets propagate towards $`x\mathrm{}`$. They correspond to particles $`3`$ and $`4`$ moving apart. In our conventions, particle $`4`$ has momentum opposite to that of particle $`1`$. We can think of it as particle $`1`$ back–scattered. The first packet is an advanced wave. It appears at $`x=0`$ at some time before the incoming wave arrives at the origin. The phase responsible for the acausal behavior is $`e^{4i\theta p^2}`$. The third packet is delayed. The opposite phase causes the delay. Similarly the terms we get from $`p_0p_0`$ are waves moving towards $`x\mathrm{}.`$ It is easy to check that the advanced wave is again produced by the phase $`e^{4i\theta p^2}`$.
Thus the collision is described as follows: The center of mass back scattering is isomorphic to bouncing off a wall. An incoming wave packet of spatial width $`\lambda ^{1/2}`$ is arranged to arrive at the wall at time $`t=0`$. The outgoing wave consists of three terms. One term appears to originate from the wall at time $`t=8p_0\theta `$, well after the incoming packet reached the wall. It is odd that the wave is delayed for so long a time as the energy increases but it is not acausal. A second term is neither significantly delayed or advanced. We will ignore it.
The other term is an “advanced” wave which appears to leave the wall before the incoming packet arrived. What is worse, the effect increases with energy so that the advance is proportional to the energy. This certainly seems acausal.
In itself, an advance does not violate causality. A simple non-relativistic model illustrates the point. Picture the incoming particles as rigid rods of length $`L`$. Assume the rod reflects when its leading end strikes the wall. In this case the center of mass of the rod will appear to reflect before it reaches the wall. In a Newtonian world a physicist measuring an advance would conclude that the scattering objects resembled rigid rods.
The problem with such rigid rods is that they conflict with the combined constraints of causality and Lorentz invariance. In fact the required properties of the rod are completely at variance with the usual expectations of special relativity. For example one usually assumes that perfectly rigid bodies can not exit. The reason is that by suddenly displacing one end of a rod, the signal would instantly appear at the other end. Since such a rod is spacelike, this is usually thought to lead to action at a distance, nonlocality and violation of causality.
Equally peculiar is the behavior of the rods under boost. Suppose the momentum is increased. The conventional expectation is that the rod will Lorentz contract thus decreasing the advance. This is precisely the opposite of what space/time non–commutativity implies. The rod seems to expand as its momentum increases.
Another phenomenon predicted by eq. (2.35) is that the outgoing packet is much broader than the incoming. Let the incoming packet be of spatial width $`\lambda ^{\frac{1}{2}}`$. By contrast, the outgoing packet has spatial width $`\lambda ^{\frac{1}{2}}\theta `$. In the limit we study of large $`\lambda \theta `$ this is broader than the incoming packet. How is this explained?
To understand this effect we return to the rod model. The advance is of order the rod size $`L`$. If we take $`L=p\theta `$ then the uncertainty in the rod size is
$$\mathrm{\Delta }L=\theta \mathrm{\Delta }p=\theta \lambda ^{\frac{1}{2}}.$$
(2.37)
This means that the advance is also uncertain by the same amount. This obviously broadens the outgoing packet by the required amount.
All three terms in eq. (2.35) can be interpreted in terms of the rod model. Each of the incoming rods has two ends, a leading and a trailing end. The advanced term is due to the scattering of the two leading ends while the retarded contribution originates from the interaction of the trailing ends. The interaction of a leading and a trailing end contributes the second term in eq. (2.35).
What are we to make out of this behavior? The most obvious response is to dismiss it as pathological and declare space/time non–commutativity to be unphysical. Our opinion is that this is prematurely pessimistic. The main reason is that some of the properties of the amplitude largely follow from the uncertainty principle implied by eq. (2.24)
$$\mathrm{\Delta }t\mathrm{\Delta }x\theta .$$
(2.38)
This uncertainty principle has the same form as the stringy uncertainty principle
$$\mathrm{\Delta }t\mathrm{\Delta }x\alpha ^{}.$$
(2.39)
It therefore behooves us to inquire into the structure of string theory amplitudes to see if they produce any behavior similar to what we find in theories with space/time non–commutativity.
## 3 Scattering in Open String Theory
In this section, we analyze tree level scattering amplitudes of open strings on branes in the presence of a background electric $`B_{\mu \nu }`$ field. In the presence of a background electric field the underlying spacetime is non–commutative. It is interesting to ask how scattering experiments similar to those studied in the previous example can probe the space/time non–commutativity. In particular we would like to investigate whether the acausal behavior we found in the simple field theory model is present. As we shall see, the amplitudes produce causal behavior and exhibit large time delays proportional to the momentum. The later phenomenon persists even without the electric field. The reader may think of the problem in a $`1+1`$ dimensional context by considering open string scattering on a stack of D1-branes . Throughout this section, we denote by $`g_s,l_s`$ the string coupling constant and length scale.
At the level of disc amplitudes the inclusion of the electric field is simple. All we need is to start with the amplitudes at $`E=0`$, replace the metric $`\eta _{\mu \nu }`$ by the effective open string metric $`G_{\mu \nu }`$, $`g_s`$ by $`G_s`$ and multiply the answer by the phase factors with non–commutativity parameter $`\theta `$. In terms of the electric field these parameters are given by
(3.1)
$`G_{\mu \nu }=(1\stackrel{~}{E}^2)\eta _{\mu \nu },\mu \nu =0,1`$ (3.2)
$`G_{\mu \nu }=\delta _{\mu \nu },\mu \nu 0,1`$ (3.3)
$`\theta ^{01}=2\pi l_s^2{\displaystyle \frac{\stackrel{~}{E}}{1\stackrel{~}{E}^2}}`$ (3.4)
$`G_s=g_s(1\stackrel{~}{E}^2)^{\frac{1}{2}}.`$ (3.5)
Here, $`\stackrel{~}{E}=E/E_{cr}1`$. The critical electric field is given by $`E_{cr}=1/2\pi l_s^2`$.
Let us consider the Veneziano amplitude describing massless open string scattering. In terms of open string parameters the amplitude has the following form:
$`A_4`$ $``$ $`G_s\left(K_{st}e^{i(p_1p_2+p_3p_4)}+K_{st}^{}e^{i(p_1p_4+p_3p_2)}\right){\displaystyle \frac{\mathrm{\Gamma }(2sl_s^2)\mathrm{\Gamma }(2tl_s^2)}{\mathrm{\Gamma }(1+2ul_s^2)}}`$ (3.6)
$`+`$ $`G_s\left(K_{su}e^{i(p_1p_2+p_4p_3)}+K_{su}^{}e^{i(p_1p_4+p_2p_3)}\right){\displaystyle \frac{\mathrm{\Gamma }(2sl_s^2)\mathrm{\Gamma }(2ul_s^2)}{\mathrm{\Gamma }(1+2tl_s^2)}}`$ (3.7)
$`+`$ $`G_s\left(K_{tu}e^{i(p_1p_3+p_2p_4)}+K_{tu}^{}e^{i(p_1p_3+p_4p_2)}\right){\displaystyle \frac{\mathrm{\Gamma }(2tl_s^2)\mathrm{\Gamma }(2ul_s^2)}{\mathrm{\Gamma }(1+2sl_s^2)}}.`$ (3.8)
The amplitude $`A_4`$ is obtained by integrating four vertex operators around the disc. We denote the two incoming particles by 1 and 2 and the two outgoing particles by 3 and 4 and let all momenta be incoming. Using Mobius invariance the vertex operators of particles 1, 2 and 3 can be put at three fixed points on the boundary of the disc – mapping it to the upper half plane these are usually taken to be $`z_1=0`$, $`z_2=1`$ and $`z_3=\mathrm{}`$ respectively. The location of the vertex operator of particle number 4, $`z_4`$, is then integrated over the real axis. Since a Mobius transformation does not change the cyclic ordering of the vertex operators, we need to add another piece obtained from fixing $`z_4=1`$ and integrating the location of particle number 2. The three terms in the answer correspond to $`\mathrm{}<z_4<0`$, $`1<z_4<\mathrm{}`$ and $`0<z_4<1`$ respectively and similarly for $`24`$. Here $`pk=\theta ^{01}(p_0k_1k_0p_1)`$.
The kinematic factors $`K`$ in (3.2) involve momenta, $`p_i`$, polarization vectors, $`\xi _i`$, and also traces over Chan Paton factors $`\lambda _i`$. The quantities $`s,t,u`$ are the Mandelstam variables
$$s=2p_1p_2,t=2p_1p_4,u=2p_1p_3$$
(3.9)
satisfying the mass shell constraint $`s+t+u=0`$. Scattering in the backward direction is defined by $`u=0`$. In eq. (3.3) we used the open string metric to contract the indices.
For the case of backward scattering, $`u=0`$, the kinematics are such that only the first term corresponding to the $`s`$–channel exchange gets multiplied by phases. One phase occurs when particles $`1`$, $`2`$ and $`3`$ are placed at $`z_1=0`$, $`z_2=1`$ and $`z_3=\mathrm{}`$ respectively and the location of particle $`4`$ is integrated from $`\mathrm{}<z_4<0`$. The opposite phase occurs when $`24`$. No phases multiply the other two terms. One of the two phases, $`e^{2\pi i\stackrel{~}{E}sl_s^2}`$, caused the appearance of the advanced waves in the non–commutative field model.
Setting $`u`$ to zero, the first term in the amplitude takes the form
$$A_{st}G_s\left(K_{st}e^{2\pi i\stackrel{~}{E}sl_s^2}+K_{st}^{}e^{2\pi i\stackrel{~}{E}sl_s^2}\right)\mathrm{\Gamma }(2sl_s^2)\mathrm{\Gamma }(2sl_s^2).$$
(3.10)
Using the identity
$$y\mathrm{\Gamma }(y)\mathrm{\Gamma }(y)=\frac{\pi }{\mathrm{sin}(\pi y)},$$
(3.11)
we can write this as follows
$$A_{st}G_s\left(K_{st}e^{2\pi i\stackrel{~}{E}sl_s^2}+K_{st}^{}e^{2\pi i\stackrel{~}{E}sl_s^2}\right)\frac{1}{s\mathrm{sin}(2\pi sl_s^2)}.$$
(3.12)
The kinematic factors $`K`$ are also simple in this case. They are proportional to $`s^2`$ times products of polarization vectors and traces over Chan Paton factors. Therefore,
$$A_{st}G_ss\left(a_1e^{2\pi i\stackrel{~}{E}sl_s^2}+a_2e^{2\pi i\stackrel{~}{E}sl_s^2}\right)\frac{1}{\mathrm{sin}(2\pi sl_s^2)},$$
(3.13)
where the constants $`a_1`$ and $`a_2`$ are independent of $`s`$.
This term has poles at $`s=n/2l_s^2`$ with $`n`$ being an integer. The divergence of the amplitude at the poles is an essential physical feature of the amplitude, a resonance corresponding to the propagation of an intermediate string state over long spacetime distances. To define the poles we use the correct $`ϵ`$ prescription replacing $`ss+iϵ`$. This has the effect of shifting the poles off the real axis. Then the function $`1/\mathrm{sin}(2\pi sl_s^2)`$ can be expanded as a power series in $`y=e^{2i\pi sl_s^2ϵ}`$. In all, this term in the amplitude takes the form
$$A_{st}G_ss\underset{n>0odd}{}a_1e^{2\pi i(n+\stackrel{~}{E})sl_s^2}+a_2e^{2\pi i(n\stackrel{~}{E})sl_s^2}+O(ϵ).$$
(3.14)
Comparing with eq. (2.30), we see that the amplitude looks similar to the case of a non–commutative field theory. We get a sum of phases with the identification
$$\theta _n^{}=2\pi (n\pm \stackrel{~}{E})l_s^2,n>0,odd.$$
(3.15)
What is interesting is that the non–commutativity parameter $`\theta `$ gets modified by stringy oscillator effects. In fact the phases persist even in the absence of the electric field. We see that $`\theta _n^{}`$ are positive for all positive odd integers.
In contrast with the field theory case, here we get phases that cause time delays only. The “acausal” phase, $`e^{2\pi i\stackrel{~}{E}s}`$, gets multiplied by powers of $`y`$ from the Gamma functions. The net effect is to produce phases which, for $`\stackrel{~}{E}<1`$, or for $`E<E_{cr}`$, cause time delays. Evidently, even in the presence of a background electric field, string scattering amplitudes produce causal behavior only. The acausal behavior due to the non–commutativity parameter $`\theta `$ is cancelled by phases from the Gamma functions. It seems that the oscillators are crucial for the causal behavior of the theory. The effects of the non–commutativity are always mixed with the effects of the string oscillators. We see another reason why space/time non–commutative field theories cannot be obtained as limits of string theory in background electric fields, as was found in . Such theories show pathological acausal behavior and are not unitary. What is interesting, however, is that the onset of the acausal behavior we found occurs as the electric field approaches its critical value.
The other two terms in formula (3.2) can be analyzed in a similar way. It is easier to write
$$A_{su}+A_{tu}=G_s(A_1+A_2),$$
(3.16)
where
$`A_1\left(K_{su}+K_{tu}\right)\left({\displaystyle \frac{\mathrm{\Gamma }(2sl_s^2)\mathrm{\Gamma }(2ul_s^2)}{\mathrm{\Gamma }(1+2tl_s^2)}}+{\displaystyle \frac{\mathrm{\Gamma }(2tl_s^2)\mathrm{\Gamma }(2ul_s^2)}{\mathrm{\Gamma }(1+2sl_s^2)}}\right)`$ (3.17)
and
$`A_2\left(K_{su}K_{tu}\right)\left({\displaystyle \frac{\mathrm{\Gamma }(2sl_s^2)\mathrm{\Gamma }(2ul_s^2)}{\mathrm{\Gamma }(1+2tl_s^2)}}{\displaystyle \frac{\mathrm{\Gamma }(2tl_s^2)\mathrm{\Gamma }(2ul_s^2)}{\mathrm{\Gamma }(1+2sl_s^2)}}\right).`$ (3.18)
Then we can analyze the sum and differences of the two combinations of Gamma functions that appear in (3.2) as $`u0`$. Again, no Moyal phases multiply these two terms.
Setting $`u=0`$, we find that
$$A_1a_3s\frac{\mathrm{cos}(2\pi sl_s^2)}{\mathrm{sin}(2\pi sl_s^2)}.$$
(3.19)
Shifting $`ss+iϵ`$, we can expand $`1/\mathrm{sin}(2\pi sl_s^2)`$ in powers of $`y`$. Thus we find
$$A_1a_3s(1+e^{4\pi isl_s^2})\underset{n0even}{}e^{2\pi insl_s^2}.$$
(3.20)
The first term in the series is just proportional to $`s`$ and produces no large time delays. The other terms are phases responsible for time delays. The phases in this term are independent of $`\theta `$. They are present even in the absence of a background field. Ordinary scattering of open strings shares features with scattering in non–commutative field theory with effective non–commutativity parameter the string length squared. However the amplitude produces only causal behavior.
The other term produces a pole at $`u=0`$. We have that
$$A_2a_4s\frac{1}{ul_s^2+iϵ}.$$
(3.21)
The pole in $`u`$ corresponds to the exchange of a massless particle and we will ignore it. This term has no oscillations in $`s`$.
The effect on the outgoing wave-packet is similar to the previous example except that the advanced waves are absent. Let us consider the case of no electric field for simplicity. Then $`\theta _n^{}=2\pi nl_s^2`$ and the open string metric $`G_{\mu \nu }=\eta _{\mu \nu }.`$ If we use eq. (2.23) for the same $`\varphi _{in}(p)`$ given in eq. (2.31), we find for a typical phase in the amplitude
$$\mathrm{\Phi }_{out}(x)G_s𝑑pp^2e^{4i\theta _n^{}p^2}\left(e^{\frac{(pp_0)^2}{\lambda }}+e^{\frac{(p+p_0)^2}{\lambda }}\right)e^{ipx}.$$
(3.22)
For $`p_0\lambda ^{1/2}1/2\pi p_0l_s^2`$, this is proportional to
$$\mathrm{\Phi }_{out}(x)G_s\frac{d^2}{dx^2}F(x;\theta _n^{},\lambda ,p_0)+(p_0p_0)$$
(3.23)
with $`F(x;\theta _n,\lambda ,p_0)`$ given by eq. (2.34).
We see that the outgoing wave-packet splits into a series of packets, one localized at $`x=0`$, and a series at $`x=8p_0\theta _n^{}n>0`$. The advanced waves are absent. Only delayed waves are present. Each delayed packet has width given by $`8\lambda ^{1/2}\theta _n^{}`$. The packets are not overlapping for $`p_0\lambda ^{1/2}1/2\pi p_0l_s^2`$. The $`n`$–th packets are more spread. After the derivatives are performed we find that the contributions to the sum are dominated by the small $`n`$ packets. The amplitude of the packets falls like $`n^{5/2}`$. Again the time delays are proportional to the energy $`p_0`$. It is interesting that the large time delays persist even in the absence of the electric field.
The interpretation is different than before. The scattering is causal. We would like to suggest the following to explain the series of time delays. As the two strings come together, an intermediate stretched string state is formed. The string state has total energy $`p_0`$. The state is oscillating from small size to a large size proportional to $`p_0`$. To see this we write
$$p_0=\frac{L}{l_s^2}+\frac{N}{L},$$
(3.24)
where $`N`$ is some oscillation number. We see that this is minimized for $`Lp_0l_s^2`$. The state begins from small size and grows to a string of maximal size of order $`p_0l_s^2`$, storing the energy as potential energy. This repeats itself periodically. With each oscillation there is an amplitude for the string to split. Thus there is an infinite sequence of delayed wave packets. The delay is proportional to $`L`$ since the string ends move with the speed of light. The intermediate state has size proportional to the energy $`p_0`$. This is a manifestation of the stringy uncertainty relation.
In the case of a background electric field we find time delays (in closed string units) proportional to
$$\mathrm{\Delta }t=\frac{p_0l_s^2}{1\stackrel{~}{E}^2}.$$
(3.25)
The time delays are proportional to $`1/T_{eff}`$, where $`T_{eff}=(1\stackrel{~}{E}^2)/l_s^2`$ is the effective tension of the open strings in the presence of the electric field. The effect of the electric field is to reduce the tension of the strings . As the electric field approaches its critical value, the time delays become longer. The extent of the intermediate state in space is also bigger. However, we note that as the field approaches its critical value, the effective coupling constant $`G_s`$ tends to zero and the amplitudes are suppressed.
We have illustrated the violations of causality in space/time non–commutative field theory and its restoration in string theory by considering the evolution of wave packets. Evidently the scattering amplitudes of the field theory violate some principle of S-matrix theory that string theory preserves. In fact it is not difficult to see what principle is involved. Macroscopic causality is usually assumed to follow from two properties of amplitudes. The first involves the location of singularities in the Mandelstam $`s`$ variable; namely, the amplitude should be analytic in the upper half plane. In the case of non–commutative field theory the amplitude in eq.(2.30) is an entire function and satisfies this rule. In the case of the string theory tree diagrams there is an infinite sequence of poles on the real axis. However with the conventional $`iϵ`$ prescription the poles are displaced to the lower half plane and lead to no violation of causality. The second requirement is that the amplitudes should not exponentially diverge along any direction in the upper half plane in order to insure that certain contours of integration can be closed. This is what is violated in the non–commutative field theory. The cosine term in eq(2.30) exponentially diverges in the upper half plane. By contrast the non–commutative Moyal phases in string theory are compensated for by the factor $`1/\mathrm{sin}(2\pi sl_s^2)`$ in formulae like eq(3.7) as long as the electric field is smaller than critical.
## 4 Conclusion
Space/time non–commutativity is a more subtle phenomenon than its space/space counterpart. If we define the space/time non–commutative deformation of a theory by multiplying its tree diagrams by space/time Moyal phases then an ordinary quantum field theory becomes acausal as well as non–unitary. The acausality is easily seen in the scattering of wave packets by the appearance of an outgoing signal that originates before the incoming particles reach each other.
By contrast, the space/time non–commutative deformation of open string theory is not acausal. The theory does not have a limit in which stringy effects disappear. These stringy effects conspire to shift the Moyal phases so that they become causal. Thus the peculiar advanced effects found in the field theory should not be though of as the signature of non–commutativity.
The delayed effects of non–commutativity in a collision process are also interesting. The space/time non–commutativity manifests itself by time delays which grow linearly with increasing momentum. As we have seen in Section 3 the time delay of the leading delayed wave is governed by a parameter $`\theta _\pm ^{}=2\pi (1\pm \stackrel{~}{E})l_s^2`$. Nothing special seems to happen to $`\theta ^{}`$ as the electric field is turned off. However $`\theta _{}^{}`$ vanishes at the critical electric field. One possible interpretation of this is that open string theory exhibits the signature of space/time non–commutativity without any electric field. Indeed this interpretation is suggested by the well known space/time uncertainty principle i.e., it appears like the string grows in the longitudinal direction to a length of order $`pl_s^2`$.
Acknowledgements
This paper is a revised version of an incorrect earlier version. We are very grateful to Igor Klebanov and Juan Maldacena for pointing out the error in the previous version. L.S. and N.T. would also like to acknowledge Igor Klebanov’s help in analyzing the open string amplitudes in Section 3.
N.S. would like to thank the University of Texas Theory Group for hospitality during part of this work. We thank R. Gopakumar, S. Minwalla, D. Morrison, M. Peskin, S. Shenker and E. Witten for useful discussions. The work of N.S. was supported in part by DOE grant #DE-FG02-90ER40542. The work of L.S. and N.T. was supported in part by NSF grant 980115.
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# Special Lagrangian cones
## 1. Introduction
Let $`Y`$ be Calabi-Yau manifold of complex dimension $`n`$ with Kähler form $`\omega `$ and non-zero parallel holomorphic $`n`$-form $`\mathrm{\Omega }`$ satisfying the normalization condition $`\omega ^n/n!=(1)^{n(n1)/2}(i/2)\mathrm{\Omega }\overline{\mathrm{\Omega }}`$. Then $`\text{Re }(\mathrm{\Omega })`$ is a calibrated form, whose calibrated submanifolds are called special Lagrangian submanifolds .
Moduli spaces of special Lagrangian submanifolds (and of other calibrated submanifolds) have appeared recently in string theory , . On physical grounds, Strominger, Yau and Zaslow argued that a Calabi-Yau manifold $`Y`$ with a mirror partner $`\widehat{Y}`$ admits a (singular) fibration by special Lagrangian tori, and that $`\widehat{Y}`$ should be obtained by compactifying the dual fibration. To make this idea rigorous one needs to have control over the singularities and compactness properties of families of special Lagrangian submanifolds. In dimensions three and higher these properties are not well understood.
Motivated by these problems we study the simplest isolated singularities of special Lagrangian varieties – homogeneous cones in $`^n`$ with an isolated singularity. These are also local models for more general singularities in that they are possible tangent cones to special Lagrangian currents at singular points.
We introduce the notion of a $`\theta `$-special Legendrian submanifold – a special class of minimal $`(n1)`$-dimensional submanifolds – of $`S^{2n1}(1)`$, and characterize $`\theta `$-special Lagrangian cones in $`^n`$ as those cones $`C`$ whose links $`L=CS^{2n1}`$ are $`\theta `$-special Legendrian submanifolds of $`S^{2n1}`$ (Proposition 2.5). From any special Legendrian link, in addition to a special Lagrangian cone, we obtain a one-parameter family of asymptotically conical special Lagrangian (possibly immersed) submanifolds.
###### Theorem A.
Let $`\mathrm{\Sigma }^{n1}`$ be a $`\theta `$-special Legendrian submanifold of $`S^{2n1}(1)^n`$. Let $`\mathrm{\Sigma }_d`$ ($`d`$) denote the set
$$\{(zp^n):p\mathrm{\Sigma },z,\text{where}\mathrm{Im}(z^n)=d,\mathrm{arg}z[0,\frac{\pi }{n}]\}.$$
Then
(i) $`\mathrm{\Sigma }_d`$ is a $`\theta `$-special Lagrangian variety.
(ii) $`\mathrm{\Sigma }_0=C(\mathrm{\Sigma })C(e^{i\pi /n}\mathrm{\Sigma })`$ where $`C(\mathrm{\Sigma })`$ denotes the cone on $`\mathrm{\Sigma }`$.
(iii) $`\mathrm{\Sigma }_d`$ is asymptotically conical, with two ends $`\mathrm{\Sigma }`$ and $`e^{i\pi /n}\mathrm{\Sigma }`$.
In the case of special Lagrangian cones in $`^3`$, results of Yau and others put restrictions on three-dimensional special Lagrangian cones. For example, we obtain:
###### Theorem B.
Let $`C`$ be a homogeneous special Lagrangian cone in $`^3`$, with $`L=CS^5(1)`$ a (possibly immersed) sphere. Then $`C`$ must be a special Lagrangian plane.
A simple corollary of this theorem is a regularity result for homogeneous solutions of the special Lagrangian graph equation in dimension three. The theorem is also sharp in the following two senses. The analogous statement in $`^4`$ is false, as recent examples of Chen et al. demonstrate. Moreover, if the link type is a torus not a sphere then even in dimension three there are nontrivial special Lagrangian cones. Our main result gives an abundance of such cones.
###### Theorem C.
There exists a countably infinite family of non-isometric special Lagrangian cones in $`^3`$. Each cone has link an embedded torus which is invariant under some $`S^1\text{SU(3)}`$.
Each special Lagrangian cone in $`^3`$ also gives rise to other calibrated cones (Lemma 2.9). For example, from each special Lagrangian cone in $`^3`$ with an isolated singularity we associate: an associative cone in $`^7`$ with an isolated singularity, a coassociative cone in $`^7`$ and a Cayley cone in $`^4`$ with a line of singularities and special Lagrangian cones in $`^{n+3}`$ with singularities along a real $`n`$-plane.
The strategy for the construction of the special Lagrangian cones in $`^3`$ is as follows. By exploiting the connection between harmonic maps and minimal surfaces we construct a two-parameter family $`u_{\alpha ,J}`$ of special Legendrian immersions $`^2S^5(1)`$. Harmonic maps from two dimensional domains to Lie groups and symmetric spaces have a rich structure, with relations to infinite dimensional completely integrable systems and loop groups , . In the case of $`S^1`$-equivariant harmonic maps to spheres, a finite dimensional completely integrable system – the C. Neumann system – appears. Several geometric features of the harmonic map have nice interpretations in terms of conserved quantities of this system. The mechanical interpretation of the Legendrian condition is not so clear, but nonetheless we are able to obtain minimal Legendrian immersions from certain solutions of the Neumann system.
###### Theorem D.
For each $`\theta [0,2\pi )`$ there exists a $`2`$-parameter family $`u_{\alpha ,J}`$, $`(\alpha ,J)[0,1]\times [0,1/3\sqrt{3}]`$, of $`\theta `$-special Legendrian immersions $`^2S^5(1)`$ with the following properties:
(i) The immersion $`u_{\alpha ,J}`$ is invariant under the $`1`$-parameter subgroup of SU(3) generated by $`A=\mathrm{diag}i(1,\alpha ,1\alpha )su(3)`$.
(ii) For $`\alpha =1`$ and $`J=1/3\sqrt{3}`$ these immersions all describe Clifford tori, but otherwise all the immersions are geometrically distinct.
(iii) The family $`u_{\alpha ,J}`$ contains all $`\theta `$-special Legendrian Legendrian immersions of the form (3.4) which cover a torus.
To obtain tori from these immersions, in Proposition 5.3 we examine the conditions under which $`u_{\alpha ,J}`$ is doubly periodic with respect to some lattice. By examining the special cases $`u_{J,0}`$ and $`u_{0,\alpha }`$ we deduce
###### Theorem E.
(i) For $`\alpha (0,1]`$, the immersion $`u_{0,\alpha }`$ is doubly periodic and hence gives rise to a minimal Legendrian torus.
(ii) For a dense set of $`J(0,1/3\sqrt{3})`$ the immersion $`u_{J,0}`$ is doubly periodic and hence gives rise to a minimal Legendrian torus.
For the family $`u_{0,\alpha }`$, referred to in part (i) of the previous theorem, we give detailed information about the geometry (e.g. conformal structure, maximum and minimum values of the Gauss curvature and embeddedness) of the corresponding surfaces. As a corollary we find (Theorem 5.5) that there are embedded ‘almost flat’ minimal Legendrian tori. These tori demonstrate sharpness of two pinching results of Yau on minimal Lagrangian (Legendrian) immersions into $`P^2`$ ($`S^5`$).
The paper is organized as follows. In Section 2 we recall basic facts about special Lagrangian geometry in $`^n`$, introduce the notion of special Legendrian in $`S^{2n1}`$ and characterize special Lagrangian cones in terms of special Legendrian links. In Section 3 we recall basic facts from harmonic map theory: principally the relation with minimal surfaces and the appearance of the C. Neumann system in $`S^1`$-equivariant harmonic maps into spheres. In Section 4 we study which solutions of the Neumann system give rise to special Legendrian immersions, give explicit parametrisations of these solutions and study the geometry of these immersions. In Section 5 we study the periodicity conditions for these immersions and hence are able to deduce our main results.
## 2. special Lagrangian cones in $`^n`$
### 2.1. Special Lagrangian geometry in $`^n`$
Special Lagrangian geometry is an example of a calibrated geometry . We review some elementary facts about calibrations and special Lagrangian geometries in $`^n`$ in particular (see for further details).
Each calibrated geometry is a distinguished class of minimal submanifolds of a Riemannian manifold $`(M,g)`$ associated with a closed differential $`p`$-form $`\varphi `$ of comass one.
For each $`mM`$, the comass of $`\varphi `$ is defined to be
$$\varphi _m^{}=sup\{<\varphi _m,\xi _m>:\xi _m\text{is a unit simple }p\text{-vector at m}\}.$$
In other words, $`\varphi _m^{}`$ is the supremum of $`\varphi `$ restricted to the Grassman of oriented $`p`$-dimensional planes $`G(p,T_mM)`$, regarded as a subset of $`\mathrm{\Lambda }^pT_mM`$.
To any form of comass one there is a natural subset of $`G(p,TM)`$
$$G_m(\varphi )=\{\xi _mG(p,T_mM):<\varphi _m,\xi _m>=1\},$$
that is, the collection of oriented $`p`$-planes on which $`\varphi `$ assumes its maximum. These planes are the planes calibrated by $`\varphi `$. An oriented $`p`$-dimensional submanifold of $`(M,g)`$ is calibrated by $`\varphi `$ if its tangent plane at each point is calibrated.
The key property of calibrated submanifolds is that they are homologically volume minimizing
###### Lemma 2.1 (Harvey and Lawson ).
Let $`(M,g,\varphi )`$ be a calibrated geometry, and suppose $`S`$ is a calibrated submanifold (possibly with boundary). Then for any oriented $`p`$-dimensional submanifold $`\widehat{S}`$ homologous to $`S`$
$$\text{vol}(S)\text{vol}(\widehat{S})$$
with equality if and only if $`\widehat{S}`$ is also calibrated (by $`\varphi `$).
Let $`z_1,\mathrm{},z_n`$ denote standard complex coordinates on $`^n`$. For any $`\theta [0,2\pi )`$ the real $`n`$-form
$$\alpha _\theta =\text{Re}(e^{i\theta }dz^1\mathrm{}dz^n)$$
is a calibrated form, called the $`\theta `$-special Lagrangian calibration on $`^n`$.
For the proof that $`\alpha _\theta `$ has comass one see . A $`\theta `$-special Lagrangian plane (we will sometimes abbreviate this as $`\theta `$-SLG) is an oriented $`n`$-plane calibrated by the form $`\alpha _\theta `$. A useful characterization of the $`\theta `$-special Lagrangian planes is
###### Lemma 2.2.
An oriented $`n`$-plane $`\xi `$ in $`^n`$ is $`\theta `$-special Lagrangian (for the correct choice of orientation) if and only if
1. $`\xi `$ is Lagrangian with respect to the standard symplectic form $`\omega =dx^idy^i`$, (i.e. $`\omega `$ restricts to zero on $`\xi `$) and
2. $`\beta _\theta :=\text{Im}(e^{i\theta }dz^1\mathrm{}dz^n)`$ restricts to zero on $`\xi `$.
One reason for considering the whole $`S^1`$-family of special Lagrangian calibrations is the following result of Harvey and Lawson (Proposition 2.17 of ):
###### Proposition 2.3.
A connected oriented Lagrangian submanifold $`S^n`$ is minimal (i.e. it is a critical point of volume, or its mean curvature $`H`$ vanishes) if and only if $`S`$ is $`\theta `$-special Lagrangian for some $`\theta `$.
### 2.2. Regular cones and special Legendrian links
For any compact connected oriented embedded submanifold $`\mathrm{\Sigma }S^{n1}(1)^n`$ define the cone on $`\mathrm{\Sigma }`$,
$$C(\mathrm{\Sigma })=\{tx:t^0,x\mathrm{\Sigma }\}.$$
A cone $`C`$ in $`^n`$ is regular if there exists $`\mathrm{\Sigma }`$ as above so that $`C=C(\mathrm{\Sigma })`$, in which case we call $`\mathrm{\Sigma }`$ the link of the cone $`C`$. $`C(\mathrm{\Sigma })0`$ is an embedded smooth submanifold, but $`C(\mathrm{\Sigma })`$ has an isolated singularity at $`0`$ unless $`\mathrm{\Sigma }`$ is a totally geodesic sphere.
To characterize the links of regular special Lagrangian cones we need to introduce some geometric structures on the unit sphere $`S^{2n1}`$ in $`^n`$. As a convex hypersurface in a Kähler manifold , $`S^{2n1}(1)`$ inherits a contact form, that is, a $`1`$-form $`\gamma `$ so that
(2.1)
$$\gamma d\gamma ^{n1}0.$$
Let $`X`$ denote the Euler vector field $`x/x`$ on $`^n`$ and $`\omega `$ denote the standard symplectic form on $`^n`$. Then the contact form on $`S^{2n1}(1)`$ is
$$\gamma =\iota _X\omega |_{S^{2n1}}.$$
Associated with $`\gamma `$ is the contact distribution, the hyperplane field $`\mathrm{ker}\gamma TS^{2n1}`$. The condition (2.1) on $`\gamma `$ ensures that the distribution $`\mathrm{ker}\gamma `$ is not integrable. The maximal dimensional integral submanifolds (i.e. submanifolds on which $`\gamma `$ restricts to zero) of the distribution are $`(n1)`$-dimensional and are called Legendrian submanifolds.
The relevance of Legendrian submanifolds of the sphere can be see from the next result whose proof is standard.
###### Lemma 2.4.
Let $`\mathrm{\Sigma }`$ be an $`(n1)`$-dimensional submanifold of $`S^{2n1}(1)`$. Then $`C(\mathrm{\Sigma })`$ is Lagrangian if and only if $`\mathrm{\Sigma }`$ is Legendrian.
For any $`p`$-form $`\varphi `$ on $`^n`$ define the normal part of $`\varphi `$ by
$$\varphi _N=\iota _X\varphi ,$$
where $`X`$ again denotes the Euler vector field on $`^n`$. In particular, $`\alpha _{\theta ,N}`$ denotes the normal part of the $`\theta `$-special Lagrangian calibration $`\alpha _\theta `$.
An oriented $`(n1)`$-dimensional submanifold $`\mathrm{\Sigma }`$ of $`S^{2n1}(1)`$ is a $`\theta `$-special Legendrian submanifold if at each point of $`\mathrm{\Sigma }`$, $`\alpha _{\theta ,N}`$ restricts to the volume form on $`\mathrm{\Sigma }`$.
###### Proposition 2.5.
A regular cone $`C=C(\mathrm{\Sigma })`$ in $`^n`$ is $`\theta `$-special Lagrangian if and only if $`\mathrm{\Sigma }`$ is $`\theta `$-special Legendrian.
###### Proof.
This is essentially a special case of Theorem 5.6 of . We sketch the proof. For any constant $`p`$-form $`\varphi `$ on $`^n`$ define the tangential part of $`\varphi `$ to be
$$\varphi _T=\iota _X\left(\frac{x}{|x|}dx\varphi \right).$$
Then $`\varphi `$ decomposes as
(2.2)
$$\varphi =\varphi _T+\frac{x}{|x|}dx\varphi _N$$
where $`\varphi _N`$ is the normal part of $`\varphi `$ defined previously. When $`d\varphi =0`$, restricting to the unit sphere it follows that
$$d\varphi _T=0\text{and}d\varphi _N=p\varphi _T.$$
Since $`\varphi _N(\xi )=\varphi (x\xi )`$ and $`\xi =x\xi `$ for any simple $`(p1)`$-vector in $`\mathrm{\Lambda }^{p1}x^{}`$, $`\varphi _N`$ still has comass one (but since $`\varphi _N`$ is not closed it is not a calibration itself). Moreover, submanifolds $`\mathrm{\Sigma }`$ of $`S^{n1}(1)`$ on which $`\varphi _N`$ restricts to the volume form are exactly those for which $`C(\mathrm{\Sigma })`$ is calibrated by $`\varphi `$. Hence the result follows by taking $`\varphi `$ to be any of the $`\theta `$-special Lagrangian calibrations $`\alpha _\theta `$. ∎
We also have the following Legendrian analogue of Proposition 2.3
###### Proposition 2.6.
A connected oriented Legendrian submanifold of $`S^{2n1}(1)`$ is minimal if and only if it is $`\theta `$-special Legendrian for some $`\theta `$.
###### Proof.
Let $`\mathrm{\Sigma }`$ be a minimal Legendrian submanifold of $`S^{2n1}(1)`$. It is a standard fact that $`C(\mathrm{\Sigma })`$ is minimal if and only if $`\mathrm{\Sigma }`$ is minimal in the unit sphere. Thus $`C(\mathrm{\Sigma })`$ is a minimal Lagrangian cone which from Proposition 2.3 must be $`\theta `$-special Lagrangian for some $`\theta `$. By the previous proposition this implies $`\mathrm{\Sigma }`$ is $`\theta `$-special Legendrian. The converse is similar. ∎
We finish the section by proving Theorem A, which gives a one-parameter family of asymptotically conical special Lagrangian varieties modeled on any special Lagrangian cone. This result generalizes Theorem 3.5 of which is our result in the special case $`\mathrm{\Sigma }=\{(x_1,\mathrm{},x_n)^n:x_i\text{with}x_i^2=1\}`$.
###### Proof of Theorem A.
(i) By Proposition 2.5, $`C(\mathrm{\Sigma })`$ is a $`\theta `$-SLG cone. By rotating $`\mathrm{\Sigma }`$ by $`A=\text{diag}(e^{i\theta /n},\mathrm{},e^{i\theta /n})`$ we can assume $`C(\mathrm{\Sigma })`$ is $`0`$-SLG. Thus $`\beta |_{C(\mathrm{\Sigma })}=0`$. Let $`\varphi :\mathrm{\Sigma }S^{2n1}`$ denote the inclusion of $`\mathrm{\Sigma }`$ in the sphere, and let $`x_1,\mathrm{},x_{n1}`$ be local coordinates on $`\mathrm{\Sigma }`$. Then $`\beta |_{C(\mathrm{\Sigma })}=0`$ is equivalent to
(2.3)
$$\text{Im}(det{}_{}{}^{}(\varphi ,\frac{\varphi }{x_1},\mathrm{},\frac{\varphi }{x_{n1}}))=0.$$
Let $`\mathrm{\Phi }:\times \mathrm{\Sigma }^n`$ be given by $`\mathrm{\Phi }(t,x)=f(t)\varphi (x)`$ where $`f:`$ is some nonconstant smooth complex valued function. It is straightforward to check that any such $`\mathrm{\Phi }`$ gives rise to a Lagrangian immersion to $`^n`$.
Let $`x_0=t`$ and for $`i=0,\mathrm{},n1`$ denote $`\mathrm{\Phi }/x_i`$ by $`\mathrm{\Phi }_i`$. Then for $`j=1,\mathrm{},n1`$
$$\omega (\mathrm{\Phi }_0,\mathrm{\Phi }_j)=\omega (\dot{f}\varphi ,f\varphi _j)=\mathrm{Re}(\overline{f}\dot{f})\omega (\varphi ,\varphi _i)+\mathrm{Im}(\overline{f}\dot{f})<\varphi ,\varphi _i>=0$$
where the first term vanishes because $`\varphi `$ is Legendrian and the second because $`|\varphi |^2=1`$. For $`j,k=1,\mathrm{},n1`$ we have
$$\omega (\mathrm{\Phi }_j,\mathrm{\Phi }_k)=\omega (f\varphi _j,f\varphi _k)=|f|^2\omega (\varphi _j,\varphi _k)=0$$
and so $`\mathrm{\Phi }`$ is Lagrangian as claimed.
Now we claim that $`\mathrm{\Phi }`$ is $`0`$-SLG if and only if $`f`$ satisfies
(2.4)
$$\mathrm{Im}(f^n)=d$$
for some real constant $`d`$. To prove this it is enough to show that $`\beta |_\mathrm{\Phi }=0`$ holds if and only if (2.4) is satisfied. But $`\beta |_\mathrm{\Phi }=0`$ is equivalent to
(2.5)
$$\text{Im}(det{}_{}{}^{}(\mathrm{\Phi }_0,\mathrm{},\mathrm{\Phi }_{n1}))=0.$$
Since $`C(\mathrm{\Sigma })`$ is $`0`$-SLG, we have $`\text{Im}(det{}_{}{}^{}(\varphi ,\varphi _1,\mathrm{},\varphi _{n1}))=0`$ and hence
$`\text{Im}(det{}_{}{}^{}(\mathrm{\Phi }_0,\mathrm{},\mathrm{\Phi }_{n1}))`$ $`=`$ $`\text{Im}(det{}_{}{}^{}(\dot{f}\varphi ,f\varphi _1,\mathrm{},f\varphi _{n1}))`$
$`=`$ $`\text{Im}(\dot{f}f^{n1}det{}_{}{}^{}(\varphi ,\varphi _1,\mathrm{},\varphi _{n1}))`$
$`=`$ $`\text{Re}(det{}_{}{}^{}(\varphi ,\varphi _1,\mathrm{},\varphi _{n1}))\times \text{Im}(\dot{f}f^{n1}).`$
Thus (2.5) holds if and only if
$$\text{Im}(\dot{f}f^{n1})=\frac{1}{n}\left(\text{Im}\frac{d}{dt}f^n\right)=\frac{1}{n}\frac{d}{dt}\text{Im}(f^n)=0.$$
Hence $`\mathrm{\Phi }`$ is $`0`$-SLG if and only if $`\mathrm{Im}(f^n)=d`$ as claimed.
Parts (ii) and (iii) are straightforward to verify. ∎
### 2.3. Minimal Legendrian surfaces
In dimension two, any special Lagrangian cone must be a union of special Lagrangian planes (since its link must be a union of Legendrian geodesics in $`S^3`$). In the first interesting case, namely special Lagrangian cones in $`^3`$, restrictions on the geometry and topology of the allowable links follow from the next result, essentially due to Yau.
###### Theorem 2.7.
Let $`\mathrm{\Sigma }`$ be a minimal Legendrian surface of $`S^5(1)`$. Then:
(i) If $`\mathrm{\Sigma }`$ has genus zero, $`\mathrm{\Sigma }`$ is totally geodesic.
(ii) If $`\mathrm{\Sigma }`$ is a complete nonnegatively curved surface, $`\mathrm{\Sigma }`$ is a totally geodesic sphere or a flat torus.
(iii) If $`\mathrm{\Sigma }`$ is complete nonpositively curved surface then $`\mathrm{\Sigma }`$ is a flat torus.
This theorem is the Legendrian analogue of a result of Yau on minimal Lagrangian immersions into Kähler surfaces of constant holomorphic sectional curvature (e.g. $`P^2`$ with the Fubini-Study metric). In fact, using Reckziegel’s observation about the local correspondence between minimal Legendrian immersions into $`S^5`$ and minimal Lagrangian immersions into $`P^2`$, one can deduce the Legendrian result from the Lagrangian one.
As a corollary of part (i) of Theorem 2.7 we deduce Theorem B. Applying this theorem to the special case of special Lagrangian graphs we deduce
###### Corollary 2.8.
Any homogeneous degree 1 solution $`u`$ of the 3-dimensional special Lagrangian graph equation
$$\mathrm{\Delta }u=det\mathrm{Hess}(u)$$
is a quadratic function.
Theorem B is sharp in the following two senses. Firstly, in $`^4`$ the analogous result is false as recent examples of Chen et al. show . Secondly, in $`^3`$ there are nontrivial special Lagrangian cones with link type a torus, the simplest example of which is the cone on a generalized Clifford torus.
Let $`T`$ be the Lagrangian product $`3`$-torus contained in $`S^5(1)`$
$$T=\{z^3:|z_i|^2=1/3,i=1,2,3\},$$
and $`T_\theta `$ be the $`2`$-torus
$$T_\theta =\{zT:\mathrm{arg}z_i=\theta \}.$$
The $`T_\theta `$, the generalized Clifford tori, are all flat minimal Legendrian tori and $`T_{\theta /3}`$, $`T_{\pi +\theta /3}`$ are $`\theta `$-special Legendrian. Harvey and Lawson discovered the cones on these tori in a family of special Lagrangian level sets invariant under the maximal torus $`T^2\text{SU(3)}`$. This high degree of symmetry allowed them to explicitly write down solutions. In the next three sections we shall find a family of nonisometric minimal Legendrian tori in $`S^5(1)`$, which are invariant under an $`S^1\text{SU(3)}`$. This symmetry is still enough to allow us to give quite explicit descriptions of these tori.
Special Lagrangian cones in dimension three with isolated singularities, also give rise naturally to several related singular calibrated varieties. For example, remarks of Harvey-Lawson (IV.2.C. Remark 2.12) and Donaldson-Thomas show the following:
###### Lemma 2.9.
If $`X^3`$ is a $`0`$-special Lagrangian variety, then:
(i) $`X\times \{pt\}^3\times `$ is an associative variety
(ii) $`X\times ^3\times `$ is a coassociative variety
(iii) $`X\times ^3\times `$ is a Cayley variety.
In cases (ii) and (iii) starting with a special Lagrangian cone with an isolated singularity we obtain cylindrical cones, which have a whole line of singularities. One also gets special Lagrangian cylindrical cones in $`^{n+3}`$ by taking the Cartesian product of a $`3`$-dimensional cone in $`^3`$ with a real $`n`$-plane in $`^n`$.
## 3. Harmonic Maps, Minimal Surfaces and the Neumann System
We shall construct $`S^1`$-invariant minimal Legendrian tori in $`S^5(1)`$ by exploiting two relationships. The first is the connection between harmonic maps and minimal surfaces. The second is the link between $`S^1`$-equivariant harmonic maps into spheres and the C. Neumann system describing motion on the sphere under a quadratic potential.
### 3.1. Harmonic Maps
We recall some definitions and basic facts from harmonic map theory. Suppose $`M`$ and $`N`$ are Riemannian manifolds. For any $`C^1`$ map $`u:MN`$ define a smooth function $`e(u)`$, the energy density of $`u`$, by $`e(u)(x)=Tr(du_x^2)`$. Define a functional on $`C^1(M,N)`$, the total energy, by $`E(u)=_Me(u)\mu _M`$, where $`\mu _M`$ is the Riemannian volume element of $`M`$. Critical points of $`E`$ are harmonic maps from $`M`$ into $`N`$.
If $`N`$ is isometrically embedded in $`^K`$ then we can view a function $`u:MN`$ as a function $`u=(u^1,\mathrm{},u^K)`$ into $`^K`$ with the constraint that $`u(x)N`$ for all $`xM`$. Then
(3.1)
$$E(u)=\underset{i=1}{\overset{K}{}}_M|u^i|^2\mu _M.$$
Extremals of $`E`$ subject to the constraint that $`u(M)N`$ give us the harmonic maps to $`N`$. From this we see that the harmonic map equations can be written simply as
$$\mathrm{\Delta }u(x)T_{u(x)}N,u(x)N,xM.$$
In the case that $`N=S^n(1)^{n+1}`$ (with the metric induced by this inclusion) this implies $`\mathrm{\Delta }u=\lambda u`$ for some function $`\lambda `$ on $`M`$. Taking the inner product of both sides with $`u`$ and using the constraint equation $`|u|^2=1`$ we determine that $`\lambda =(u,\mathrm{\Delta }u)=|du|^2`$. Summarizing we have
###### Lemma 3.1.
A smooth map $`u:MS^n(1)^{n+1}`$ is harmonic iff and only if $`u`$ satisfies the equation
(3.2)
$$\mathrm{\Delta }u=|du|^2u.$$
Finally we recall what happens to the harmonic map equations when we make a conformal change of metric on the domain $`M`$. If $`\stackrel{~}{g}=\lambda ^2g`$ then $`\stackrel{~}{g}^1=\lambda ^2g^1`$, and $`\stackrel{~}{\mu }_M=\lambda ^m\mu _M`$. Hence $`E_{\stackrel{~}{g}}(u)=\lambda ^{m2}E_g(u)`$ and we see that $`E`$ is conformally invariant if and only if dim $`M=2`$. Therefore in dimension 2 harmonicity depends only on the structure of $`M`$ as a Riemann surface. In particular there is a natural quadratic differential $`\mathrm{\Phi }`$, the Hopf differential. If $`z`$ is a local complex coordinate on $`M`$ then $`\mathrm{\Phi }=\varphi (z)dz^2`$ where
(3.3)
$$\varphi (z)=(u_z,u_z)=\frac{1}{4}\left(|u_x|^2|u_y|^22i(u_x,u_y)\right).$$
Harmonicity of $`u`$ implies that $`\mathrm{\Phi }`$ is holomorphic. If the Hopf differential vanishes the map $`u`$ is conformal. Moreover, we have the following connection with minimal surfaces:
###### Proposition 3.2.
() $`u`$ is harmonic and conformal if and only if $`u`$ is a (branched) minimal immersion.
### 3.2. Equivariant Harmonic Maps and the Neumann System
For harmonic maps from $`^2`$ to $`S^5(1)^3`$ (where both $`^2`$ and $`S^5(1)`$ are given their standard metrics) of the special form
(3.4)
$$u(s,t)=e^{As}z(t)$$
where $`Aso(6)`$ and $`z:S^5(1)`$, it follows from (3.2) that $`u`$ is harmonic if and only if $`z`$ satisfies
(3.5)
$$\ddot{z}+A^2z=(|\dot{z}|^2+|Az|^2)z$$
where $`\dot{}`$ denotes differentiation with respect to $`t`$. As Uhlenbeck noted these are the equations of motion for the C. Neumann problem of motion of a particle on a sphere under the quadratic potential $`|Az|^2`$.
Define $``$ actions on $`^2`$ and $`S^5(1)`$ by
$$\gamma (s,t)=(s+\gamma ,t)$$
$$\gamma p=e^{A\gamma }p$$
where $`s,t,\gamma `$ and $`pS^5(1)`$. These induce an action in the usual manner on the Banach manifold $`C^1(^2,S^5)`$ by
(3.6)
$$\left(\gamma u\right)(x)=\gamma u(\gamma ^1x)$$
the fixed points of which are precisely maps of the form (3.4). Since $``$ acts by isometries on both $`^2`$ and $`S^5`$, it follows from the definition of $`E`$ that it is an $``$-invariant function on $`C^1(^2,S^5)`$. Hence we could also appeal to Palais’s Principle of Symmetric Criticality to find the equations satisfied by $`z`$, as in .
From now on we consider only the case that $`Au(3)`$, so that the one-parameter group $`e^{As}`$ preserves both the metric and the symplectic structure. Then by conjugation we may assume that $`A=\text{diag}i(\lambda _1,\lambda _2,\lambda _3).`$ In this case equation (3.5) becomes
(3.7)
$$\ddot{z_j}\lambda _j^2z_j=\lambda z_j,z_j,j=1,2,3$$
where
(3.8)
$$\lambda =|Az|^2+|\dot{z}|^2.$$
It will be convenient to rewrite the equations slightly. Writing $`z_j=R_je^{i\theta _j}`$ we see that (3.7) is equivalent to
(3.9)
$$\ddot{R_j}\frac{J_j^2}{R_j^3}=(\lambda _j^2\lambda )R_j,j=1,2,3$$
where $`\theta _j`$ is determined up to a constant by the relation $`J_j=R_j^2\dot{\theta _j}`$.
There are some obvious conserved quantities. From conservation of energy we have
(3.10) $`H=|\dot{z}|^2|Az|^2`$
and conservation of the quantities
(3.11) $`J_j=x_j\dot{y_j}y_j\dot{x_j},j=1,2,3`$
expresses the fact that angular momentum in each of the three complex planes $`z_1,z_2,z_3`$ is conserved. For details of other less obvious conserved quantities of the Neumann system we refer the reader to .
The condition that $`u`$ be conformal is conveniently expressed in terms of the integrals of motion. Namely, $`u`$ is conformal if and only if
(3.12) $`|u_s|^2|u_t|^2=|\dot{z}|^2|Az|^2=H=0`$
and
(3.13) $`(u_s,u_t)=(\dot{z},Az)={\displaystyle \underset{i=1}{\overset{3}{}}}\lambda _iJ_i=0.`$
Summarizing we have
###### Proposition 3.3.
$`u(s,t)=e^{As}z(t):^2S^5(1)`$ is a minimal immersion if and only if $`z`$ satifies the equations of motion of the Neumann system (3.7) and the conserved quantities $`H,J_j`$ satisfy the constraints (3.12) and (3.13).
## 4. $`S^1`$ Equivariant minimal Legendrian immersions
### 4.1. The Legendrian constraints
For $`u(s,t)=e^{As}z(t)`$ to be a minimal Legendrian immersion, besides the conditions of Proposition 3.3, two further constraints must hold
(4.1) $`\alpha (u_s)=\omega (u,u_s)=\omega (z,Az)={\displaystyle \underset{i=1}{\overset{3}{}}}\lambda _iR_i^2=0,`$
and
(4.2) $`\alpha (u_t)=\omega (u,u_t)=\omega (z,\dot{z})={\displaystyle \underset{i=1}{\overset{3}{}}}J_i=0.`$
Note that the second equation corresponds merely to further constraints on the values of the integrals of the Neumann system. The first equation is more mysterious and is in general not preserved under the flow of the Neumann system. In fact, we have
###### Lemma 4.1.
There are minimal Legendrian immersions of the form given in (3.4) if and only if $`A\text{su(3)}`$.
###### Proof.
From Proposition 2.6 any minimal Legendrian immersion is $`\theta `$-special Legendrian for some $`\theta `$ and hence $`\beta _\theta `$ restricts to zero on the cone. At a point $`xu`$ on the cone (where $`x^+`$) we have
$$\beta _\theta |_{C(u)}=x^2\text{Im}(e^{i\theta }det{}_{}{}^{}(u,u_s,u_t))=x^2\text{Im}(e^{i{\scriptscriptstyle \lambda _is}}e^{i\theta }det{}_{}{}^{}(z(t),Az(t),\dot{z}(t))).$$
Since this must hold for all real $`s`$ and $`t`$, for $`\beta _\theta `$ to restrict to zero we must have $`\lambda _i=0`$ as claimed.
One can also show necessity directly from the equations for a minimal Legendrian equation by showing that the constraints (3.12,3.13,4.1,4.2) are not consistent with the equations of motion of the Neumann system unless $`Asu(3)`$.
To see this let us compute the second derivative of the mysterious constraint $`c:=\omega (z,Az)`$ for a solution of the Neumann system at an instant when all the constraints and their first derivatives are satisfied. One finds
$$\ddot{c}=\omega (Az,\ddot{z})+\omega (A\dot{z},\dot{z})=\omega (Az,A^2z)+\omega (A\dot{z},\dot{z}).$$
Let $`c_1=\omega (A^2z,Az)`$ and $`c_2=\omega (A\dot{z},\dot{z})`$. Then $`c_1`$ may be expressed in terms of the symmetric polynomials in the $`\lambda _i`$ as
$$c_1=\lambda _i^3R_i^2=(\lambda _j)(\lambda _i^2R_i^2)(\lambda _j\lambda _k)(\lambda _iR_i^2)+\lambda _1\lambda _2\lambda _3(R_i^2).$$
Hence using the constraints we have
$$c_1=(\lambda _j)|Az|^2+\lambda _1\lambda _2\lambda _3.$$
A calculation shows
$$c_2=\lambda _1\lambda _2\lambda _3\frac{|\dot{z}|^2}{|Az|^2}$$
and so
(4.3)
$$\ddot{c}=(\lambda _j)|Az|^2H\frac{\lambda _1\lambda _2\lambda _3}{|Az|^2}.$$
Clearly once we have imposed the constraint $`H=0`$, $`\ddot{c}=0`$ if and only if $`Asu(3)`$. Moreover, by differentiating (4.3) it is easy to verify that all higher derivatives of the constraint $`c`$ also vanish when $`Asu(3)`$. Thus to show existence of minimal Legendrian immersions we need only show there exist initial conditions for the Neumann system which satisfy all the constraints together with the first derivative of the mysterious constraint. We will see that this is indeed the case in the proof of Theorem D which we now give. ∎
###### Proof of Theorem D.
Let $`u`$ be a minimal Legendrian immersion of the form (3.4), i.e. $`u(s,t)=e^{As}z(t)`$ where $`Asu(3)`$. By conjugation we may assume $`A=i\text{diag}(\lambda _1,\lambda _2,\lambda _3)`$ where $`\lambda _1\lambda _20>\lambda _3`$. Let $`\alpha =\lambda _2/\lambda _1`$, then $`\alpha [0,1]`$ and $`A=i\text{diag}\lambda _1(1,\alpha ,1\alpha )`$. Moreover, by rescaling $`s`$ and $`t`$ we may assume that $`\lambda _1=1`$. Let
$$\stackrel{}{1}=(1,1,1),\stackrel{}{J}=(J_1,J_2,J_3),\stackrel{}{\lambda }=(\lambda _1,\lambda _2,\lambda _3),\stackrel{}{R^2}=(R_1^2,R_2^2,R_3^2).$$
Then the constraints (3.13,4.1,4.2) together with the constraint that $`z`$ lie on the unit sphere can be written as
(4.4) $`\stackrel{}{1}\stackrel{}{J}=0,`$ $`\stackrel{}{\lambda }\stackrel{}{J}=0,`$
(4.5) $`\stackrel{}{1}\stackrel{}{R^2}=0,`$ $`\stackrel{}{\lambda }\stackrel{}{R^2}=1`$
and $`Asu(3)`$ is equivalent to $`\stackrel{}{1}\stackrel{}{\lambda }=0`$. Let $`\stackrel{}{\mu }`$ be the cross product of $`\stackrel{}{1}`$ and $`\stackrel{}{\lambda }`$
$$\stackrel{}{\mu }=\stackrel{}{1}\times \stackrel{}{\lambda }=(12\alpha ,2+\alpha ,\alpha 1).$$
The constraints in (4.4) are equivalent to
(4.6)
$$\stackrel{}{J}=J\stackrel{}{\mu }$$
for some constant $`J`$, while the constraints in (4.5) are equivalent to
(4.7)
$$\stackrel{}{R^2}(t)=\gamma (t)\stackrel{}{\mu }+\frac{1}{3}\stackrel{}{1}$$
for some function $`\gamma (t)`$. The remaining constraint $`|\dot{z}|^2=|Az|^2`$ then becomes
(4.8)
$$\frac{\dot{\gamma }^2}{4}+J^2=R_1^2R_2^2R_3^2$$
or in terms of $`\gamma `$
(4.9)
$$\frac{\dot{\gamma }^2}{4}+J^2=\gamma ^3\mu _1\mu _2\mu _3+\frac{\gamma ^2}{3}\underset{ij}{}\mu _i\mu _j+\frac{1}{27}.$$
Since we seek periodic solutions we may assume that $`\gamma (0)=\gamma _0>0`$, $`\dot{\gamma }(0)=0`$. Then at $`t=0`$, (4.9) becomes
(4.10)
$$\gamma ^3\mu _1\mu _2\mu _3+\frac{\gamma ^2}{3}\underset{ij}{}\mu _i\mu _j+\frac{1}{27}=J^2$$
and thus specifying $`\gamma _0`$ determines $`J^2`$ (and vice versa).
Let us fix $`\alpha [0,1]`$ and consider the case where $`J0`$. Given $`J(0,\frac{1}{3\sqrt{3}}]`$, (4.10) has a unique smallest nonnegative root $`\gamma _+(J)`$. Let $`\gamma (0)=\gamma _+(J),\dot{\gamma }(0)=0`$. Then up to a translation in time any periodic solution of (4.10) (except possibly a solution corresponding to $`J=0`$ which we shall treat later) arises from such an initial condition. Once the initial conditions for $`\gamma `$ and $`\dot{\gamma }`$ have been specified, (4.7) fixes $`R_j(0)`$ and $`\dot{R}_j(0)`$ for $`j=1,2,3`$. Given $`J`$ and $`\alpha `$, (4.6) fixes $`\stackrel{}{J}`$. If we define $`\dot{\theta }_j=J_j/R_j^2`$, then $`\dot{\theta _j}(0)`$ is determined by (4.6) and (4.7) for $`j=1,2,3`$. By a global rotation in SU(3) (e.g. replacing $`z(t)`$ by $`Bz(t)`$ where $`B=\mathrm{exp}(i\mathrm{diag}(\sigma _1,\sigma _2,\sigma _3))\text{SU(3)}`$) we may rotate $`z(t)`$ so that $`\theta _2(0)=\theta _3(0)=0`$. We may not assume also that $`\theta _1(0)=0`$ without allowing $`B\text{U(3)}`$ in which case we will change the value of $`\theta `$ for which $`u`$ is $`\theta `$-special Legendrian. In the case $`J0`$ we shall verify later that choosing $`\theta _1(0)=\theta `$ gives rise to a $`\theta `$-special Legendrian immersion (in the case $`J=0`$ choosing $`\theta _1(0)=\theta +\pi /2`$ gives rise to a $`\theta `$-special Legendrian immersion).
Thus for each $`\theta [0,2\pi )`$, and $`(\alpha ,J)[0,1]\times (0,1/3\sqrt{3}]`$ there is a unique solution of the Neumann equation (given by specifying initial data in the manner above) which satisfies the constraints (3.12,3.13,4.1,4.2). Hence by the proof of the previous lemma it gives rise to a minimal Legendrian immersion which we denote $`u_{\alpha ,J}`$.
In the case $`J=0`$ we will explicitly exhibit solutions later in this section and see that as $`\alpha 0`$ the period of $`\gamma `$ becomes infinite, and that the limiting solution $`u_{0,0}`$ describes a minimal Legendrian sphere (which as previously noted is necessarily totally geodesic).
To see which immersions $`u_{\alpha ,J}`$ are geometrically distinct consider in greater detail the geometry of these immersions. Since the immersions are all conformal, the metric $`g`$ induced on $`^2`$ can be described by a single positive function $`y=|Az|^2=|\dot{z}|^2`$, where $`g=y|dz|^2`$. A calculation shows that $`\gamma `$ and $`y`$ are related by
(4.11)
$$y=\gamma \mu _1\mu _2\mu _3+\frac{1}{3}\lambda _i^2.$$
It follows from (4.11) and (4.9) that $`y`$ satisfies
(4.12)
$$\dot{y}^2+4y^32y^2\lambda _i^2=4C$$
where
$$C=\lambda _1^2\lambda _2^2\lambda _3^2+J^2\mu _1^2\mu _2^2\mu _3^2.$$
The Gauss curvature of the immersion satisfies
(4.13)
$$K=\frac{(\mathrm{ln}y)^{^{\prime \prime }}}{2y}=1+2Cy^3.$$
In the case $`J=1/3\sqrt{3}`$, the corresponding solution of (4.9) is $`\gamma 0`$ independent of the choice of $`\alpha [0,1]`$ and hence $`u_{\alpha ,1/3\sqrt{3}}`$ has $`K0`$. It follows that $`u`$ must be (a piece of) a generalized Clifford torus. Similarly, if $`\alpha =1`$ then $`\mu _3=0`$ and it follows from (4.11) that $`y2`$. Once again $`K0`$ and hence $`u_{1,J}`$ is a (piece of a) generalized Clifford torus.
All other immersions $`u_{\alpha ,J}`$ are geometrically distinct. To begin with, note that the remaining immersions are all invariant under a unique $`1`$-parameter family of SU(3) – the subgroup generated by $`A=i\mathrm{diag}(1,\alpha ,1\alpha )`$. For $`\alpha [0,1)`$ these are all inequivalent, hence $`u_{\alpha ,J}`$ and $`u_{\stackrel{~}{\alpha },\stackrel{~}{J}}`$ are distinct when $`\alpha \stackrel{~}{\alpha }`$. Now fix $`\alpha `$ and consider $`u_{\alpha ,J}`$ for $`J(0,1/3\sqrt{3})`$. We claim that the minimum and maximum values of the Gauss curvature $`K`$ are respectively strictly decreasing and increasing functions of $`J`$ on $`(0,1/3\sqrt{3})`$. It follows that $`u_{\alpha ,J}`$ and $`u_{\alpha ,\stackrel{~}{J}}`$ are geometrically distinct when $`\alpha \stackrel{~}{\alpha }`$.
To proof the previous claim, note that (4.13) shows that for a given immersion $`u_{\alpha ,J}`$ the minimum (maximum) value of $`K`$ occurs at the minimum (maximum) value of $`y`$. From (4.12) it is clear that for fixed $`\alpha `$, $`y_{min}`$ and $`y_{max}`$, the minimum and maximum values attained by y, are strictly decreasing and increasing functions of $`J`$ respectively. Since $`C`$ is a decreasing function of $`J`$, from (4.13) we see that the minimum and maximum values of $`K`$ are, like $`y`$, strictly decreasing and increasing functions of $`J`$ respectively as claimed. ∎
It is also possible to write down explicit solutions in terms of elliptic functions. Let us express $`\gamma `$ in terms of the Jacobi elliptic functions. Recall that $`\gamma `$ satisfies the equation
$$\frac{\dot{\gamma }^2}{4}+J^2=\gamma ^3\mu _1\mu _2\mu _3+\frac{\gamma ^2}{3}\underset{ij}{}\mu _i\mu _j+\frac{1}{27}$$
and that for $`J^2[0,1/27)`$ and $`\alpha 1`$ there are three solutions $`\mathrm{\Gamma }_1,\mathrm{\Gamma }_2,\mathrm{\Gamma }_3`$ to this equation when $`\dot{\gamma }=0`$. Let us label these solutions so that $`\mathrm{\Gamma }_20\mathrm{\Gamma }_1\mathrm{\Gamma }_3`$. Then we can rewrite the previous equation as
(4.14)
$$\dot{\gamma }^2=4\mu _1\mu _2\mu _3(\gamma \mathrm{\Gamma }_1)(\gamma \mathrm{\Gamma }_2)(\gamma \mathrm{\Gamma }_3).$$
###### Proposition 4.2.
$`\gamma (t)=\mathrm{\Gamma }_2(\mathrm{\Gamma }_2\mathrm{\Gamma }_1)\mathrm{sn}^2(rt,k)`$ is a solution of (4.14) provided
$$r^2=\mu _1\mu _2\mu _3(\mathrm{\Gamma }_3\mathrm{\Gamma }_2),k^2=\frac{\mathrm{\Gamma }_2\mathrm{\Gamma }_1}{\mathrm{\Gamma }_2\mathrm{\Gamma }_3}$$
where $`\mathrm{sn}`$ is the Jacobi elliptic sn-noidal function.
###### Proof.
The proof is a straightforward computation using the basic properties of the Jacobi elliptic functions (for details see ). ∎
From this proposition and (4.7) we derive expressions for $`R_j^2`$
(4.15)
$$R_j^2=\mu _j(\gamma \gamma _j)=\mu _j\left((\mathrm{\Gamma }_2\gamma _j)(\mathrm{\Gamma }_2\mathrm{\Gamma }_1)\mathrm{sn}^2(rt,k)\right)$$
where $`\gamma _j=1/\mu _j`$.
As promised in the proof of Theorem D we now provide explicit solutions for the $`J=0`$ case.
###### Proposition 4.3.
For each $`\theta [0,2\pi )`$, there exists a family of $`\theta `$-special Legendrian immersions $`u_{\alpha ,0}:^2S^5(1)`$, for $`\alpha [0,1]`$, whose Gauss curvature $`K`$ satisfies (4.20) and (4.21) (and hence are all distinct). Moreover, $`u_{0,0}`$ gives rise to a $`\theta `$-special Legendrian sphere and is the only member of the family $`u_{\alpha ,J}`$ to do so.
###### Proof.
In the case $`J=0`$ we know explicitly the values of the $`\mathrm{\Gamma }_i`$
$$\mathrm{\Gamma }_i=\gamma _i=\frac{1}{3\mu _i},i=1,2,3$$
and hence
(4.16)
$$r^2=(1+2\alpha ),k^2=\frac{1\alpha ^2}{1+2\alpha }.$$
Equation (4.15) specializes to
(4.17) $`R_1`$ $`=`$ $`\mu _1(\gamma _2\gamma _1)\mathrm{cn}(rt,k)`$
(4.18) $`R_2`$ $`=`$ $`\mu _2(\gamma _1\gamma _2)\mathrm{sn}(rt,k)`$
(4.19) $`R_3`$ $`=`$ $`\mu _3(\gamma _2\gamma _3)\mathrm{dn}(rt,k).`$
Define $`u_{\alpha ,0}`$ by the formula
$$u_{\alpha ,0}(s,t)=e^{As}(e^{i(\theta +\pi /2)}R_1(t),R_2(t),R_3(t))$$
where as previously we set $`A=i\mathrm{diag}(1,\alpha ,1\alpha )`$ for $`\alpha [0,1]`$. Then $`u`$ is a $`\theta `$-special Legendrian immersion invariant under $`e^{As}`$.
To find the extreme values taken on by the Gauss curvature, note that in the case $`J=0`$ we have $`y_{min}=\lambda _2\lambda _3=\alpha (1+\alpha )`$ and $`y_{max}=\lambda _1\lambda _3=1+\alpha `$. Thus
(4.20)
$$K_{min}=1+\frac{2\lambda _1^2}{\lambda _2\lambda _3}=1\frac{2}{\alpha (1+\alpha )}$$
and
(4.21)
$$K_{max}=1+\frac{2\lambda _2^2}{\lambda _1\lambda _3}=1\frac{2\alpha ^2}{1+\alpha }.$$
From (4.16) we see that $`k^21`$ as $`\alpha 0`$, and $`k^20`$ as $`\alpha 1`$. In these two limits $`\mathrm{sn}`$ reduces to $`\mathrm{tanh}`$ and $`\mathrm{sin}`$ respectively. Thus in the limiting case $`\alpha =0,J=0`$ we have
$$\gamma =\frac{1}{6}+\frac{1}{2}\mathrm{tanh}^2t,R_1=R_3=\frac{1}{\sqrt{2}}\mathrm{sech}t,R_2=\mathrm{tanh}t.$$
Finally, one can show that in order for any immersion of the form $`u(s,t)=e^{As}z(t)`$ to describe a harmonic sphere, the limit of $`z(t)`$ as $`t\pm \mathrm{}`$ must be a fixed point of the action $`e^{As}`$ . Moreover, all the conserved quantities of the Neumann system must also be zero (since they are zero at a fixed point). For $`Asu(3)`$ as above, $`e^{As}`$ has nonzero fixed points if and only if $`\alpha =0`$, in which case any point of the form $`(0,z_2,0)^3`$ is fixed. From equation (4.6), all three angular momenta $`J_j`$ are zero if and only if $`J=0`$. Thus $`u_{0,0}`$ is the only $`u_{\alpha ,J}`$ which could describe a minimal sphere. In this case $`u`$ (in the $`0`$-special Legendrian case) has the explicit form
(4.22)
$$u(s,t)=(\frac{1}{\sqrt{2}}ie^{is}\mathrm{sech}t,\mathrm{tanh}t,\frac{1}{\sqrt{2}}e^{is}\mathrm{sech}t)$$
and we can see directly that the $`2`$-sphere described is the intersection of the plane
$$i\overline{z_1}=z_3,\text{Im}z_2=0$$
with the $`5`$-sphere (and hence is totally geodesic). ∎
## 5. Periodicity conditions
In order to analyze the periodicity of the immersions $`u_{\alpha ,J}`$ we need the following lemma whose proof is a short computation (see for full details).
###### Lemma 5.1.
For $`J0`$ the sum of the angles $`\theta _i`$ and $`\dot{\gamma }`$ satisfy
(5.1)
$$\dot{\gamma }(t)=2J\mathrm{tan}\left(\theta _i(t)\theta _i(0)\right).$$
Since we chose $`\gamma `$ so that $`\dot{\gamma }(0)=0`$, this lemma has the following obvious corollary:
###### Corollary 5.2.
If $`T`$ is the period of $`\gamma `$ then $`\theta _i(T)=\theta _i(0)+n\pi `$, for some integer $`n`$.
The previous lemma is also useful in verifying what conditions on $`\theta _1(0)`$ ensure that the immersions $`u_{\alpha ,J}`$ are $`\theta `$-special Legendrian. For this we need to compute $`\beta _\theta `$ restricted to the cone on $`u`$. At a point $`(x,s,t)`$ on the cone
$$\beta _\theta =x^2\text{Im}(e^{i\theta }det{}_{}{}^{}(u,u_s,u_t))=x^2\text{Im}(e^{i\theta }det{}_{}{}^{}(z,Az,\dot{z})).$$
If $`J=0`$, so that $`\theta _j`$ are all constant we have
$$det{}_{}{}^{}(z,Az,\dot{z})=ie^{i{\scriptscriptstyle \theta _j(0)}}|Az|^2$$
and hence the immersion is $`\theta `$-special Legendrian where $`\theta `$ depends only on the initial sum of the angles $`\theta _j`$. For example, if we choose $`\theta _1(0)=\pi /2`$ or $`\theta _1(0)=3\pi /2`$ (and $`\theta _2(0)=\theta _3(0)=0`$) then the cones are $`0`$-SLG.
If $`J0`$ a short computation using Lemma 5.1 shows that
$$det{}_{}{}^{}(z,Az,\dot{z})=\frac{i|Az|^2e^{i{\scriptscriptstyle \theta _i(0)}}}{R_1^2R_2^2R_3^2}(J+i\dot{\gamma }/2)(\dot{\gamma }/2+iJ)=|Az|^2e^{i{\scriptscriptstyle \theta _i(0)}}$$
so that now choosing $`\theta _1(0)=0`$ or $`\theta _1(0)=\pi `$ gives $`0`$-special Legendrian immersions.
Suppose that $`(\sigma ,\tau )`$ is a period of $`u(s,t)=e^{As}z(t)`$, i.e.
(5.2)
$$u(s+\sigma ,t+\tau )=u(s,t)s,t.$$
Then the periodicity properties of $`u`$ are characterized by
###### Proposition 5.3.
(a) $`(\sigma ,\tau )`$ is a period of $`u_{\alpha ,J}`$ implies $`\tau `$ is an integer multiple of $`T_{\alpha ,J}`$, the basic period of $`y_{\alpha ,J}=|Az|^2`$
(b) If $`u`$ admits two independent periods then it admits a period of the form $`(\sigma ,0)`$
(c) $`u`$ admits a period of the form $`(\sigma ,0)`$ if and only if $`\alpha `$
(d) $`u`$ admits two independent periods if and only if
(5.3)
$$\alpha ,\frac{1}{2\pi }(\alpha \theta _1(T)\theta _2(T)).$$
###### Proof.
(a) Differentiating (5.2) with respect to $`s`$ and taking the norm of both sides implies $`|Az(t+\tau )|=|Az(t)|`$.
(b) If the periods are $`(\sigma _1,n_1T)`$ and $`(\sigma _2,n_2T)`$ then $`(n_1\sigma _2n_2\sigma _1,0)`$ is also a period.
(c) $`(\sigma ,0)`$ is a period implies $`e^{i\sigma \lambda _j}=1`$, for $`j=1,2,3`$. So $`\sigma \lambda _j2\pi `$. In particular, $`\alpha =\frac{\lambda _2}{\lambda _1}`$. Conversely if $`\alpha =\frac{m}{n}`$ then $`(\frac{2n\pi }{\lambda _1},0)`$ is a period.
(d) If $`u`$ admits two independent periods then $`\alpha `$ is rational by (b) and (c). By (a) any period is of the form $`(\sigma ,mT)`$. Now since $`\dot{\theta _j}`$ is $`T`$-periodic we have $`\theta _j(nT)=n\theta _j(T)`$. Then periodicity with respect to $`(\sigma ,mT)`$ is equivalent to $`e^{i\lambda _j\sigma +im\theta _j(T)}=1`$. Hence $`\sigma \lambda _j+m\theta _j(T)2\pi `$, for $`j=1,2,3`$. Together with rationality of $`\alpha `$ this implies $`\alpha \theta _1(T)\theta _2(T)2\pi `$.
Conversely, by (c) the rationality of $`\alpha `$ gives us one period $`(\sigma _1,0)`$. From above $`(\sigma ,mT)`$ is a period if and only if $`\sigma \lambda _j+m\theta _j(T)2\pi ,j=1,2,3`$. By assumption $`\frac{1}{2\pi }(\alpha \theta _1(T)\theta _2(T)=\frac{M}{N}`$ for some integers $`M`$ and $`N`$. With $`\sigma =\frac{2N\theta _1(T)}{\lambda _1}`$ and $`m=2N`$ the period condition becomes
$$2N\left(\frac{\lambda _j}{\lambda _1}\theta _1(T)+\theta _j(T)\right)2\pi ,j=1,2,3.$$
For $`j=1`$ this condition is trivial, while it holds for $`j=2`$ because
$$2N\left(\frac{\lambda _2}{\lambda _1}\theta _1(T)+\theta _2(T)\right)=4N\pi \left(\alpha \theta _1(T)\theta _2(T)\right)=4M\pi $$
Since $`\lambda _i=0`$ and by Corollary 5.2, $`\theta _3(T)=n\pi \theta _1(T)\theta _2(T)`$, we have
$$2N\left(\frac{\lambda _3}{\lambda _1}\theta _1(T)+\theta _3(T)\right)=2N\left(\alpha \theta _1(T)\theta _2(T)+n\pi \right)=4M\pi +2\pi Nn.$$
So the $`j=3`$ period condition also holds, and hence $`(\frac{2N\theta _1(T)}{\lambda _1},2NT)`$ is a second period of $`u`$. ∎
Two cases of the previous proposition are particularly interesting: when $`J=0`$ or $`\alpha =0`$. For the case $`J=0`$ we prove the following result which implies Theorem C and part (i) of Theorem E.
###### Proposition 5.4.
For $`\alpha (0,1]`$, the immersion $`u_{0,\alpha }`$ is doubly periodic and hence gives rise to a minimal Legendrian torus. Further, let $`\alpha =\frac{m}{n}`$, where $`m<n`$ and $`(m,n)=1`$. If $`mn`$ is even, then the period lattice of $`u_{\alpha ,0}`$ is rectangular with basis $`\omega _1=(2n\pi ,0)`$, $`\omega _2=(0,4\mathrm{Ke}(k)/r)`$. Otherwise the period lattice is not rectangular and is generated by $`\omega _1=(2n\pi ,0)`$ and $`\omega _3=(n\pi ,2\mathrm{Ke}(k)/r)`$. In either case each such torus $`T_{m,n}`$ is embedded and its Gauss curvature satisfies (4.20) and (4.21).
Notation in the proposition: $`k`$ and $`r`$ are defined as a functions of $`\alpha `$ by (4.16), and $`\mathrm{Ke}`$ is the complete elliptic integral defined by
$$\mathrm{Ke}(k)=_0^{\pi /2}\frac{dx}{\sqrt{1k^2\mathrm{sin}^2x}}$$
(the period of $`\mathrm{sn}(t,k)`$ is $`4\mathrm{Ke}(k)`$).
###### Proof.
Since $`J=0`$, the $`\theta _i`$ are constant and the second condition of part (d) of the previous proposition is superfluous. Thus the immersion is doubly periodic if and only if $`\alpha `$. Let $`\alpha =\frac{m}{n}`$. It is easy to see that $`\omega _1=(2n\pi ,0)`$ and $`\omega _2=(0,4\mathrm{Ke}/r)`$ belong to the period lattice of $`u_{\alpha ,0}`$. To find the full period lattice it is sufficient to find all periods in the rectangle $`R`$ formed by $`0`$, $`\omega _1`$, $`\omega _2`$ and $`\omega _1+\omega _2`$. Let $`(\sigma ,\tau )`$ be such a period. By part (a) of the previous proposition $`\tau `$ must be a integer multiple of $`2\mathrm{Ke}/r`$, the basic period of $`y`$. It is easy to see that the smallest period of the form $`(\sigma ,0)`$ occurs when $`\sigma =2n\pi `$ and so we need only deal with the case $`\tau =2\mathrm{Ke}/r`$. Using the fact that $`\mathrm{cn}(t+2\mathrm{Ke})=\mathrm{cn}(t)`$, $`\mathrm{sn}(t+2\mathrm{Ke})=\mathrm{sn}(t)`$, $`\mathrm{dn}(t+2\mathrm{Ke})=\mathrm{dn}(t)`$ we find that $`(\sigma ,2\mathrm{Ke}/r)`$ is a period if and only iff $`\sigma `$ satisfies
(5.4)
$$e^{i\sigma }=1,e^{i\sigma \alpha }=e^{i\sigma m/n}=1,e^{i(1+\alpha )}=1.$$
Clearly the third equation is implied by the first two. Moreover, the first equation implies $`e^{im\sigma }=(1)^m`$, whereas the second implies $`e^{im\sigma }=(1)^n`$. Hence if either $`m`$ or $`n`$ is even (both cannot be even since we assumed $`(m,n)=1`$) then these two equations are inconsistent. Thus there are no further periods and the period lattice is generated by $`\omega _1`$ and $`\omega _2`$. If both $`m`$ and $`n`$ are odd, then one can check that $`\sigma =n\pi `$ is the unique solution in $`[0,2n\pi )`$. Hence $`w_3=(n\pi ,2\mathrm{Ke}/r)`$ is the only new period in the rectangle $`R`$ and in this case the period lattice is generated by $`\omega _1`$ (or $`\omega _2`$) and $`\omega _3`$.
Let us show embeddedness in the case where one of $`m`$, $`n`$ is even. The other case is similar, but a little more involved since the period lattice is not rectangular. We need to show that if $`s,\stackrel{~}{s}[0,2n\pi )`$ and $`t,\stackrel{~}{t}[0,4\mathrm{Ke}/r)`$ and $`u(s,t)=u(\stackrel{~}{s},\stackrel{~}{t})`$ then $`s=\stackrel{~}{s}`$ and $`t=\stackrel{~}{t}`$. From our explicit formulae for $`R_i`$ we see that $`u(s,t)=u(\stackrel{~}{s},\stackrel{~}{t})`$ is equivalent to
(5.5) $`e^{is}\mathrm{cn}(t/r)`$ $`=`$ $`e^{i\stackrel{~}{s}}\mathrm{cn}(\stackrel{~}{t}/r)`$
(5.6) $`e^{is}\mathrm{sn}(t/r)`$ $`=`$ $`e^{i\stackrel{~}{s}}\mathrm{sn}(\stackrel{~}{t}/r)`$
(5.7) $`e^{i(1+\alpha )s}\mathrm{dn}(t/r)`$ $`=`$ $`e^{i(1+\alpha )\stackrel{~}{s}}\mathrm{dn}(\stackrel{~}{t}/r).`$
Certainly this implies $`|\mathrm{cn}t/r|=|\mathrm{cn}\stackrel{~}{t}/r|`$, which implies there exists some $`T[0,\mathrm{Ke}]`$ such that $`rt,r\stackrel{~}{t}\{T,2\mathrm{Ke}T,2\mathrm{Ke}+T,4\mathrm{Ke}T\}`$. If $`t`$ and $`\stackrel{~}{t}`$ are distinct, there are essentially two different cases, depending on whether $`\mathrm{cn}t=\mathrm{cn}\stackrel{~}{t},\mathrm{sn}t=\mathrm{sn}\stackrel{~}{t}`$ or $`\mathrm{cn}t=\pm \mathrm{cn}\stackrel{~}{t},\mathrm{sn}t=\mathrm{sn}\stackrel{~}{t}`$. In the first case the three equations above reduce to
$$e^{i\sigma }=1,e^{i\alpha \sigma }=1,e^{i\sigma (1+\alpha )}=1$$
where $`\sigma =s\stackrel{~}{s}`$. That is, we have the same equations as occurred in the periodicity part of the proof. Since we assumed one of $`m`$ and $`n`$ was even, the first two equations are inconsistent unless $`t=\stackrel{~}{t}`$ in which case $`s=\stackrel{~}{s}`$ is also forced. In the second case the equations reduce to
$$e^{i\sigma }=\pm 1,e^{i\alpha \sigma }=1,e^{i\sigma (1+\alpha )}=1.$$
Clearly, the first two equations are inconsistent with the third one. ∎
In the case $`J0`$, $`\alpha =0`$, the conditions in part (d) reduce to $`\theta _2(T)`$. Using the explicit expressions given in the previous chapter and properties of elliptic functions one can show that viewed as a function of $`J`$, $`\theta _2(T)`$ is strictly monotone. Part (ii) of Theorem E follows.
We conclude with the following result which demonstrates the sharpness of the pinching results on minimal Legendrian immersions given in parts (ii) and (iii) of Theorem 2.7.
###### Theorem 5.5.
For any $`ϵ>0`$ there exists an embedded minimal Legendrian torus $`T`$ in $`S^5`$ which is not flat, but for which $`sup_{xT}|K(x)|<ϵ`$, where $`K`$ is the Gauss curvature of $`T`$.
###### Proof.
Consider an immersion $`u_{\alpha ,J}`$ with $`J=0`$ and $`\alpha =1\delta `$. From (4.20) and (4.21) the minimum and maximum values of the Gauss curvature are given by
(5.8)
$$K_{min}(u_{1\delta ,0})=\frac{\delta (3+\delta )}{(1\delta )(2\delta )}$$
and
(5.9)
$$K_{max}(u_{1\delta ,0})=\frac{\delta (32\delta )}{2\delta }.$$
Certainly for $`\delta <\frac{1}{2}`$ we have $`|K_{min}|<7\delta `$ and similarly for $`K_{max}`$. Since $`u_{1\delta ,0}`$ gives rise to an embedded minimal Legendrian torus whenever $`\alpha `$, just choose $`\delta (0,ϵ/7)`$ and the result is proved. ∎
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# 1 Introduction
## 1 Introduction
The Minimal Supersymmetric Standard Model (MSSM) provides one of the first realistic attempts to describe physics beyond the Standard Model. Although there is a huge amount of new fields and parameters in any SUSY-extended model, the MSSM is able to reproduce all the successful predictions of the original GWS theory with a very good accuracy. The MSSM Higgs sector is enlarged due to the fact that supersymmetry forbids the usual mechanism of generation of the up-type quark masses. This problem requires introduction of another Higgs doublet with opposite charges, which produces all the necessary quark mass-terms. Consequently there are five Higgs particles and three unphysical Goldstone bosons in the theory.
The global $`N=1`$ supersymmetry, being a new symmetry in addition to the original $`SU(3)_cSU(2)_\mathrm{l}U(1)`$ local gauge invariance, puts constraints on the Higgs sector, which is rather indefinite in nonsupersymmetric theories. In the case of MSSM this leads to the famous tree-level relation $`m_hm_Z|\mathrm{cos2}\beta |`$, which bounds the mass of the lightest Higgs scalar. Note that up to now there is no experimental evidence of such state; the present day lower limit is $`m_h88.3GeV,`$ see ,. However, radiative corrections modify the upper bound substantially ,. This bound is closely related to the tree-level sum-rule
$$m_h^2+m_H^2m_A^2m_Z^2=0$$
(1)
Renormalising this rule one can obtain the shift of the unsatisfactory bound descending from the radiative corrections. Up to now there are many papers performing the computation on the one-lop level ,,, and also some two-loop results were already obtained . Most of these papers are concentrated on the evaluation of the finite part of the correction only without explicit discussion of the divergences. It is expected that the ultraviolet divergences in such relations originating from the additional symmetry cancel, but there is still no general proof of this based on the relevant Ward identities. From this point of view it may be useful to demonstrate explicitly the mechanism of the compensation in the particular case of relation (1); this, in fact, is one of the goals of this work. In this respect, this paper supplements the calculations presented in the cited literature.
The whole analysis is performed for the top-stop sector only. It is sufficient due to the observation that the leading term coming from any fermion-sfermion cluster is proportional to the fourth power of the fermionic mass ,. Note that the contributions coming from the chargino and neutralino sectors are negligible .
The paper is organized as follows: the definition of the renormalised parameters is briefly reviewed in section 1; the U-gauge allows us to simplify the matter essentially compared to the choice of . Section 2 is devoted to discussion of the UV-divergences originating from the one-loop Feynman graphs renormalising the relevant 2-point Green functions. The leading logarithmic term is recovered in the third section. Most of the technical details are deferred to Appendices.
## 2 Definition of the renormalised parameters
As was stated before the whole computation will be performed in the unitary gauge; this particular choice reduces the number of diagrams to be considered and simplifies the renormalisation scheme. On the other hand, the presence of Goldstone bosons in $`R_\xi `$ gauges causes for example the total cancellation of contributions descending from the tadpole diagrams (see ), which does not occur in U-gauge.
The diagrams to be considered are listed in Appendix A. The ultraviolet divergences coming from the loops are handled using the standard technique of dimensional regularisation, see .
### 2.1 Renormalized (pseudo)scalar masses:
The renormalized masses $`m_X`$of (pseudo)scalars $`h`$,$`H`$ and $`A`$ are defined generically by
$$\left(1+\delta Z_X\right)m_{XB}^2=m_X^2+\delta m_X^2;$$
(2)
Next, let us denote the sum of all (relevant) graphs by $`i\mathrm{\Pi }_X(q)=i\mathrm{\Pi }_X^Y(q)`$. Then the 2-point Green functions can be written in the form (see ).
$$i\mathrm{\Gamma }_X^{(2)}(q,q)=q^2m_X^2i\mathrm{\Pi }_X(q)+i\left(\delta Z_Xq^2\delta m_X^2\right)+\mathrm{higher}\mathrm{order}$$
We adopt the so-called on-shell renormalisation scheme in which all the external momenta ($`q`$’s) are taken to be on the mass-shell, i.e. $`q^2=m_X^2`$ where $`m_X^2`$ denotes the squared mass of the considered particle. In this scheme we use the following renormalisation conditions
$$\mathrm{\Gamma }_X^{(2)}(q,q)=0,\frac{\mathrm{\Gamma }_X^{(2)}}{q^2}(q,q)=1\mathrm{at}q^2=m_X^2$$
This particular choice implies
$$\delta Z_X=0+\mathrm{higher}\mathrm{order};\delta m_X^2=\mathrm{\Pi }_X\left(q^2=m_X^2\right)+\mathrm{higher}\mathrm{order}$$
(3)
The physical mass $`m_X`$ can then be expressed (using (2) and (3)) as
$$m_X^2=m_{BX}^2+\mathrm{\Pi }_X(q^2=m_X^2)+\mathrm{higher}\mathrm{order}$$
### 2.2 Renormalized $`Z`$-boson mass:
Let us denote the sum of all the one-loop ($`Z`$) ’vacuum polarisation graphs’ by $`i\mathrm{\Pi }_Z^{\mu \nu }(q)`$. This quantity renormalizes the $`Z`$-boson mass to the new value
$$m_Z^2=m_{ZB}^2A_Z(q^2=m_Z^2)+\mathrm{higher}\mathrm{order}$$
(we have again used the on-shell conditions) where $`A_Z(q^2)`$ is defined by
$$\mathrm{\Pi }_Z^{\mu \nu }(q^2)A_Z(q^2)g^{\mu \nu }+B_Z(q^2)q^\mu q^\nu $$
(i.e. corresponds to the coefficient of the transverse part of $`\mathrm{\Pi }_Z^{\mu \nu }`$).
### 2.3 Renormalised sum-rule:
Having defined renormalised quantities we can recast the relation (1) in the renormalised form
$$m_h^2+m_H^2m_A^2m_Z^2=\mathrm{\Delta }+\mathrm{higher}\mathrm{order}$$
where
$$\mathrm{\Delta }\mathrm{\Pi }_h(q^2=m_h^2)+\mathrm{\Pi }_H(q^2=m_H^2)\mathrm{\Pi }_A(q^2=m_A^2)+A_Z(q^2=m_Z^2)$$
(4)
This is the most important relation of this section. In the following part we attempt to evaluate the one-loop leading-log term of $`\mathrm{\Delta }`$. As was already stated before the leading term descends from the graphs involving top and supertop loops so the rest of this computation will be performed for this sector only. Note that the full quantity includes contributions from almost all the particles in the theory, which would complicate the calculation essentially without any impact on the leading term so the other contributions are simply omitted.
## 3 Cancellation of UV-divergences
In this section we show that the UV-divergent parts of the diagrams listed in Appendix A cancel. To proceed we put the external momenta on-shell and substitute in (4). For the sake of brevity there will be no difference between the symbols used for the divergent parts of the considered expressions and the full contributions in this section; moreover the overall factors $`\mathrm{C}_{\mathrm{uv}}`$ and $`N_c`$ are suppressed too. For example $`B_h^t`$ (see Appendix A) corresponds here to $`g_{htt}^2(4\pi ^2)^1\left(3m_t^2\frac{1}{2}m_h^2\right)`$. To simplify the reader’s insight the definitions of the partial sums of divergences (denoted by UV with relevant sub- and superscripts) take care of the sign of the corresponding expressions in (4).
### 3.1 UV divergences in graphs involving top loops
Let us start with the divergences descending from the graphs involving one top-quark loop. The first three graphs in (A.1) give
$$B_h^t+B_H^tB_A^t=$$
(5)
$$=\frac{g^2m_t^2}{16\pi ^2m_W^2\mathrm{sin}^2\beta }\left[3m_t^2m_t^2\mathrm{cos}^2\beta +\frac{1}{2}\left(m_h^2\mathrm{cos}^2\alpha m_H^2\mathrm{sin}^2\alpha +m_A^2\mathrm{cos}^2\beta \right)\right]$$
Next, the divergent part of the fourth graph in (A.1) contributing to (4) is after some algebra
$$B_Z^t=\frac{g^2m_t^2m_Z^2}{32\pi ^2m_W^2}+\frac{g^2m_Z^4}{24\pi ^2m_W^2}\left(\epsilon _{L}^{t}{}_{}{}^{2}+\epsilon _{R}^{t}{}_{}{}^{2}\right)$$
Utilising relations(36) of Appendix B the tadpole graphs in (A.1) give
$$T_A^{ht}T_A^{Ht}=\frac{g^2m_t^4}{16\pi ^2m_W^2}$$
(6)
Summing up the partial results (5)-(6) one obtains the total divergence coming from the graphs involving one top-quark loop:
$$\mathrm{UV}_{\mathrm{top}}=\frac{g^2}{16\pi ^2m_W^2\mathrm{sin}^2\beta }\left[2m_t^4m_t^2m_Z^2\mathrm{sin}^2\beta +\frac{2}{3}m_Z^4\mathrm{sin}^2\beta \left(\epsilon _{L}^{t}{}_{}{}^{2}+\epsilon _{R}^{t}{}_{}{}^{2}\right)\right]$$
(7)
### 3.2 UV divergences in graphs involving supertop loops
The same brief list of divergences will be now built up for the diagrams with one supertop loop. The first type graphs in (A.2) give
$$\mathrm{UV}_{\mathrm{stop}}^{(1)}B_h^{\stackrel{~}{t}_1\stackrel{~}{t}_1}+B_h^{\stackrel{~}{t}_2\stackrel{~}{t}_2}+B_h^{\stackrel{~}{t}_1\stackrel{~}{t}_2}+B_H^{\stackrel{~}{t}_1\stackrel{~}{t}_1}+B_H^{\stackrel{~}{t}_2\stackrel{~}{t}_2}+B_H^{\stackrel{~}{t}_1\stackrel{~}{t}_2}B_A^{tildet_1\stackrel{~}{t}_2}=$$
$$=\frac{1}{16\pi ^2}\left(g_{h\stackrel{~}{t}_1\stackrel{~}{t}_1}^2+g_{h\stackrel{~}{t}_2\stackrel{~}{t}_2}^2+2g_{h\stackrel{~}{t}_1\stackrel{~}{t}_2}^2+g_{H\stackrel{~}{t}_1\stackrel{~}{t}_1}^2+g_{H\stackrel{~}{t}_2\stackrel{~}{t}_2}^2+2g_{H\stackrel{~}{t}_2\stackrel{~}{t}_2}^2+2g_{A\stackrel{~}{t}_1\stackrel{~}{t}_2}^2\right)$$
The result of the computation is
$$\mathrm{UV}_{\mathrm{stop}}^{(1)}=\frac{g^2}{16\pi ^2m_W^2\mathrm{sin}^2\beta }[2m_t^4\frac{1}{2}m_t^2(A_tm_6\mathrm{sin}\beta +\mu \mathrm{cos}\beta )^2+$$
$$+m_t^2m_Z^2\mathrm{sin}^2\beta m_Z^4\mathrm{sin}^2\beta (\epsilon _{L}^{t}{}_{}{}^{2}+\epsilon _{R}^{t}{}_{}{}^{2})]$$
(8)
Next, the total contribution descending from the second type graphs in (A.2) reads
$$\mathrm{UV}_{\mathrm{stop}}^{(2)}B_Z^{\stackrel{~}{t}_1\stackrel{~}{t}_1}+B_Z^{\stackrel{~}{t}_2\stackrel{~}{t}_2}+B_Z^{\stackrel{~}{t}_1\stackrel{~}{t}_2}=\frac{g^2m_Z^2}{16\pi ^2m_W^2}\{\frac{1}{3}m_Z^2(\epsilon _{L}^{t}{}_{}{}^{2}+\epsilon _{R}^{t}{}_{}{}^{2})$$
(9)
$$2[m_{\stackrel{~}{t}_1}^2(\epsilon _{L}^{t}{}_{}{}^{2}\mathrm{cos}^2\theta _t+\epsilon _{R}^{t}{}_{}{}^{2}\mathrm{sin}^2\theta _t)+m_{\stackrel{~}{t}_2}^2(\epsilon _{L}^{t}{}_{}{}^{2}\mathrm{sin}^2\theta _t+\epsilon _{R}^{t}{}_{}{}^{2}\mathrm{cos}^2\theta _t)]\}$$
The divergence coming from the tadpole sector of (A.2) can be written in the form (note that the minus sign corresponds to the sign of $`\mathrm{\Pi }_A`$ in (4))
$$\mathrm{UV}_{\mathrm{stop}}^{(3)}T_A^{h\stackrel{~}{t}_1}T_A^{h\stackrel{~}{t}_2}T_A^{H\stackrel{~}{t}_1}T_A^{H\stackrel{~}{t}_2}$$
with the result
$$\mathrm{UV}_{\mathrm{stop}}^{(3)}=\frac{g^2}{32\pi ^2m_W^2}\{m_t^2(m_{\stackrel{~}{t}_1}^2+m_{\stackrel{~}{t}_2}^2)+m_t^2(A_tm_6+\mu \mathrm{cot}\beta )^2+$$
(10)
$$+m_Z^2\mathrm{cos2}\beta [m_{\stackrel{~}{t}_1}^2(\epsilon _L^t\mathrm{cos}^2\theta _t\epsilon _R^t\mathrm{sin}^2\theta _t)+m_{\stackrel{~}{t}_2}^2(\epsilon _L^t\mathrm{sin}^2\theta _t\epsilon _R^t\mathrm{cos}^2\theta _t)]\}$$
The last part of the total UV divergence originates from the seagull-type diagrams in (A.2):
$$\mathrm{UV}_{\mathrm{stop}}^{(4)}S_h^{\stackrel{~}{t}_1}+S_h^{\stackrel{~}{t}_2}+S_H^{\stackrel{~}{t}_1}+S_H^{\stackrel{~}{t}_2}S_A^{\stackrel{~}{t}_1}S_A^{\stackrel{~}{t}_2}+S_Z^{\stackrel{~}{t}_1}+S_Z^{\stackrel{~}{t}_2}$$
After some algebra one gets
$$\mathrm{UV}_{\mathrm{stop}}^{(4)}=\frac{g^2m_t^2}{32\pi ^2m_W^2}\{(m_{\stackrel{~}{t}_1}^2+m_{\stackrel{~}{t}_2}^2)+$$
$$+m_Z^2\mathrm{cos2}\beta \left[m_{\stackrel{~}{t}_1}^2\left(\epsilon _L^t\mathrm{cos}^2\theta _t\epsilon _R^t\mathrm{sin}^2\theta _t\right)+m_{\stackrel{~}{t}_2}^2\left(\epsilon _L^t\mathrm{sin}^2\theta _t\epsilon _R^t\mathrm{cos}^2\theta _t\right)\right]$$
$$4m_Z^2[m_{\stackrel{~}{t}_1}^2(\epsilon _{L}^{t}{}_{}{}^{2}\mathrm{cos}^2\theta _t+\epsilon _{R}^{t}{}_{}{}^{2}\mathrm{sin}^2\theta _t)+m_{\stackrel{~}{t}_2}^2(\epsilon _{L}^{t}{}_{}{}^{2}\mathrm{sin}^2\theta _t+\epsilon _{R}^{t}{}_{}{}^{2}\mathrm{cos}^2\theta _t)]\}$$
(11)
With all the partial results (7),(8),(9),(10) and (11) at hand it is already easy to state (utilising (36) from Appendix B) that the divergent parts of the considered diagrams contributing to the relation (4) exactly cancel :
$$\mathrm{UV}_{\mathrm{top}}+\mathrm{UV}_{\mathrm{stop}}^{(1)}+\mathrm{UV}_{\mathrm{stop}}^{(2)}+\mathrm{UV}_{\mathrm{stop}}^{(3)}+\mathrm{UV}_{\mathrm{stop}}^{(4)}=0$$
It is the main result of this section. Such a cancellation of divergences in relations originating from the supersymmetry is a typical feature of SUSY-theories .
## 4 Finite part of $`\mathrm{\Delta }`$
The topic of this section is to compute the finite part of expression (4). Having proved the total cancellation of divergences it can be easily shown that the finite part does not depend on the mass-scale $`\mu `$. This can be seen from the fact, that the only $`\mu `$-dependent factor $`\mathrm{ln}\mu ^2`$ can be joined to the divergent factor $`\mathrm{C}_{\mathrm{uv}}`$ which drops out. In the explicit expressions we may, for convenience, put $`\mu =1`$ (mass unit).
The idea of the following computation is to split the set of the considered diagrams with respect to the magnitude of the typical mass and discuss these clusters of contributions separately. At this place it is indeed necessary to mention several restrictions put on the parameters of the theory:
* The masses of squarks and the top mass are much bigger than the other masses involved in the computation - $`m_{\stackrel{~}{t}_1},m_{\stackrel{~}{t}_2},m_tm_W,m_A,m_H,m_h`$)
* The off-diagonal entries in the $`\stackrel{~}{t}_L\stackrel{~}{t}_R`$ mass-matrix are very small with respect to $`m_t^2`$ which implies that there is no significant mixing in the supertop sector (see Appendix B for clarification). This condition can be translated in the mathematical form as $`A_tm_6+\mu \mathrm{cot}\beta m_t`$.
The first condition is relevant for (almost) the whole (experimentally admissible) area of the MSSM parametric space with the only exceptional case that $`m_A`$ is very massive. However, the mass $`m_A`$ often appears together with the factor $`\mathrm{cos}\beta `$, which is expected to be small enough to suppress such term. The second condition is more speculative but seems to be true for all the squark species (see ); for our purposes this will be taken as an assumption .
From the previous lines and the explicit form of the contributions listed in Appendix A it can be easily seen that the most important terms should be of order $`m_{\stackrel{~}{t}}^2`$. However, due to the previous assumptions, this term turns out to be small compared to the contribution coming from the terms of order $`\frac{m_t^4}{m_Z^2}`$.
The discussion below proves this statement. The notation is again abbreviated as in the previous section i.e. there will be no difference between the symbols for the finite part and the full contribution of the examined diagram.
### 4.1 Contributions of magnitude $`m_{\stackrel{~}{t}_i}^2`$
Looking at the coupling constants in Appendix A and taking into account the relation (39) from Appendix B one can check that the only contributions proportional to $`m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2`$ come from the second, third and fourth type graphs in (A.2). The relevant expression is defined by
$`\mathrm{F}_{m_{\stackrel{~}{t}_i}^2}`$ $``$ $`[B_Z^{\stackrel{~}{t}_1\stackrel{~}{t}_2}+B_Z^{\stackrel{~}{t}_2\stackrel{~}{t}_2}+B_Z^{\stackrel{~}{t}_1\stackrel{~}{t}_2}T_A^{h\stackrel{~}{t}_1}T_A^{h\stackrel{~}{t}_2}T_A^{H\stackrel{~}{t}_1}T_A^{H\stackrel{~}{t}_2}+`$ (12)
$`+S_h^{\stackrel{~}{t}_1}+S_h^{\stackrel{~}{t}_2}+S_H^{\stackrel{~}{t}_1}+S_H^{\stackrel{~}{t}_2}S_A^{\stackrel{~}{t}_1}S_A^{\stackrel{~}{t}_2}+S_Z^{\stackrel{~}{t}_1}+S_Z^{\stackrel{~}{t}_2}]_{m_{\stackrel{~}{t}_i}^2only}`$
After some manipulations one obtains
$`\mathrm{F}_{m_{\stackrel{~}{t}_i}^2}={\displaystyle \frac{N_cg^2m_Z^2}{16\pi ^2m_W^2}}\{2(\epsilon _{L}^{t}{}_{}{}^{2}\mathrm{cos}^2\theta _t+\epsilon _{R}^{t}{}_{}{}^{2}\mathrm{sin}^2\theta _t)m_{\stackrel{~}{t}_1}^2(1\mathrm{ln}m_{\stackrel{~}{t}_1}^2)+`$
$`+2\left(\epsilon _{L}^{t}{}_{}{}^{2}\mathrm{sin}^2\theta _t+\epsilon _{R}^{t}{}_{}{}^{2}\mathrm{cos}^2\theta _t\right)m_{\stackrel{~}{t}_2}^2\left(1\mathrm{ln}m_{\stackrel{~}{t}_2}^2\right)`$
$`\left(\epsilon _L^t\mathrm{cos}^2\theta _t+\epsilon _R^t\mathrm{sin}^2\theta _t\right)^2\left[2m_{\stackrel{~}{t}_1}^22{\displaystyle _0^1}dxD_{m_Z}^{\stackrel{~}{t}_1\stackrel{~}{t}_1}(x)\mathrm{ln}D_{m_Z}^{\stackrel{~}{t}_1\stackrel{~}{t}_1}(x)\right]`$
$`\left(\epsilon _L^t\mathrm{sin}^2\theta _t+\epsilon _R^t\mathrm{cos}^2\theta _t\right)^2\left[2m_{\stackrel{~}{t}_2}^22{\displaystyle _0^1}dxD_{m_Z}^{\stackrel{~}{t}_2\stackrel{~}{t}_2}(x)\mathrm{ln}D_{m_Z}^{\stackrel{~}{t}_2\stackrel{~}{t}_2}(x)\right]`$
$$\frac{1}{2}\mathrm{sin}^22\theta _t(\epsilon _L^t\epsilon _R^t)^2[m_{\stackrel{~}{t}_1}^2+m_{\stackrel{~}{t}_2}^22_0^1\mathrm{d}xD_{m_Z}^{\stackrel{~}{t}_1\stackrel{~}{t}_2}(x)\mathrm{ln}D_{m_Z}^{\stackrel{~}{t}_1\stackrel{~}{t}_2}(x)]\}$$
$$\frac{N_cg^2m_t^2}{32\pi ^2m_W^2\mathrm{sin}^2\beta }\left(A_tm_6\mathrm{sin}\beta +\mu \mathrm{cos}\beta \right)^2_0^1dx\mathrm{ln}D_0^{\stackrel{~}{t}_1\stackrel{~}{t}_2}(x)$$
(13)
Note that the last term must indeed be taken into account here because of (39). Fortunately it is strongly suppressed by the assumption of a small mixing in the supertop sector, see Appendix B.
First, it can be checked immediately that the terms of the form $`const.\times m_{\stackrel{~}{t}_i}^2`$ exactly cancel. The remaining structure is already not so easy to handle. The assumption of a small mixing allows us to approximate $`\mathrm{sin}\theta _t0`$, which drops out of the penultimate term in (13). The key role of this observation consists in the fact that the rest of (13) already does not involve any mixed $`D_{m_Z}^{\stackrel{~}{t}_1\stackrel{~}{t}_2}(x)`$ term. Thus the expressions proportional to $`m_{\stackrel{~}{t}_1}^2`$ and $`m_{\stackrel{~}{t}_2}^2`$ split into two independent clusters. To proceed one can use the expansion
$$\mathrm{ln}D_{m_Z}^{\stackrel{~}{t}_i\stackrel{~}{t}_i}(x)=\mathrm{ln}m_{\stackrel{~}{t}_i}^2+\mathrm{ln}\left[1\frac{m_Z^2}{m_{\stackrel{~}{t}_i}^2}x(1x)\right]=\mathrm{ln}m_{\stackrel{~}{t}_i}^2\frac{m_Z^2}{m_{\stackrel{~}{t}_i}^2}x(1x)+O\left(\frac{m_Z^4}{m_{\stackrel{~}{t}_i}^4}\right)$$
(14)
which originates from the definition of $`D_{m_Z}^{\stackrel{~}{t}_i\stackrel{~}{t}_i}`$, see Appendix A. Neglecting the contributions proportional to $`m_Z^2`$ one can check that the terms of the type $`m_{\stackrel{~}{t}_i}^2\mathrm{ln}m_{\stackrel{~}{t}_i}^2`$ cancel too. The previous discussion leads to the following result:
* In the case of no significant mixing within the supertop sector the contribution of the magnitude $`m_{\stackrel{~}{t}_i}^2`$ turns out to be negligible compared to the correction proportional to $`\frac{m_t^4}{m_Z^2}`$, which is investigated in the next subsection.
If the mixing in the supertop sector is not negligible one can obtain large negative contribution proportional to $`m_{\stackrel{~}{t}_i}^2`$ from this cluster of diagrams; for more comprehensive discussion see .
### 4.2 Contribution proportional to $`m_t^4m_Z^2`$, leading $`log`$-term
Looking into (A.1) and (A.2), one can immediately write down the sum of relevant terms :
$`\mathrm{F}_{m_t^4\times M^2}`$ $``$ $`[B_h^t+B_H^tB_A^tT_A^{ht}T_A^{Ht}+`$ (15)
$`+B_h^{\stackrel{~}{t}_1\stackrel{~}{t}_1}+B_h^{\stackrel{~}{t}_2\stackrel{~}{t}_2}+B_H^{\stackrel{~}{t}_1\stackrel{~}{t}_1}+B_H^{\stackrel{~}{t}_2\stackrel{~}{t}_2}]_{m_t^4\times M^2only}`$
Note that the graph $`B_Z^t`$ produces only a factor of order $`m_t^2`$ and therefore does not belong into this sum.
Let us first deal with the graphs involving top quark loop (the first five terms in (15)). One can easily check that they contribute by
$$\mathrm{F}_{m_t^4\times M^2}^{top}=\frac{N_cg^2m_t^2}{16\pi ^2m_W^2\mathrm{sin}^2\beta }\times $$
$$\times \{\mathrm{cos}^2\alpha [m_t^2+_0^1\mathrm{d}x\mathrm{ln}D_{m_h}^{tt}(x)[3m_t^2+3m_h^2x(1x)]]+$$
$$+\mathrm{sin}^2\alpha \left[m_t^2+_0^1dx\mathrm{ln}D_{m_H}^{tt}(x)\left[3m_t^2+3m_H^2x(1x)\right]\right]+$$
(16)
$$+\mathrm{cos}^2\beta [m_t^2+_0^1\mathrm{d}x\mathrm{ln}D_{m_A}^{tt}(x)[m_t^23m_A^2x(1x)]]m_t^2\mathrm{sin}^2\beta (1\mathrm{ln}m_t^2)\}$$
Note that all the irrelevant parts of orders $`m_t^2`$, $`m_Z^2`$,$`m_h^2`$ and $`m_H^2`$ were neglected; the possibly large factor $`m_A^2m_t^2\times M^2`$ is suppressed by $`\mathrm{cos}\beta `$. The other factors like $`m_h^2m_t^2\times M^2`$ or $`m_H^2m_t^2\times M^2`$ are assumed not to be high above $`m_Z^2m_t^2\times M^2`$, moreover they are put down by an overall factor coming from the integration over $`x`$. The magnitude of the total error is approximately 10 %. In addition it is easy to see that the non-logarithmic factors cancel.
The situation in the supertop-loop cluster (corresponding to the last four terms in (15)) is again quite complicated due to the structure of the appropriate squares of the couplings (28). Extracting only the relevant parts involving $`m_t^4`$ one gets
$$\mathrm{F}_{m_t^4\times M^2}^{stop}=\frac{N_cg^2m_t^4}{16\pi ^2m_W^2\mathrm{sin}^2\beta }\{\mathrm{cos}^2\alpha [_0^1\mathrm{d}x\mathrm{ln}D_{m_h}^{\stackrel{~}{t}_1\stackrel{~}{t}_1}(x)+_0^1\mathrm{d}x\mathrm{ln}D_{m_h}^{\stackrel{~}{t}_2\stackrel{~}{t}_2}(x)]+$$
$$+\mathrm{sin}^2\alpha [_0^1\mathrm{d}x\mathrm{ln}D_{m_H}^{\stackrel{~}{t}_1\stackrel{~}{t}_1}(x)+_0^1\mathrm{d}x\mathrm{ln}D_{m_H}^{\stackrel{~}{t}_2\stackrel{~}{t}_2}(x)]\}$$
(17)
The last thing to be done is to apply expansion similar to (14) on the $`D`$’s in the integrals. Summing up the particular results above it is straightforward to obtain the leading logarithmic term ($`N_c=3`$)
$$\mathrm{\Delta }=\frac{3g^2m_t^4}{16\pi ^2m_W^2\mathrm{sin}^2\beta }\mathrm{ln}\left(\frac{m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2}{m_t^4}\right)+\mathrm{}$$
(18)
This is the main result of the whole computation. It shows that including the leading logarithmic one-loop correction the original tree-level relation (1) can be recast in the form
$$m_h^2+m_H^2m_A^2m_Z^2=\frac{3g^2m_t^4}{16\pi ^2m_W^2\mathrm{sin}^2\beta }\mathrm{ln}\left(\frac{m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2}{m_t^4}\right)+\mathrm{}$$
This relation agrees with the results presented in the literature (, ).
Note that the relative error of the approximations used to derive the previous relation does not exceed $`10`$%. It is mainly due to neglecting all the terms of order $`m_t^2`$ and lower. Next, the form of the leading term (18) is invalid in case that any significant mixing in the stop sector occurs.
## Conclusion
The paper is devoted to one of the most important features of the SUSY-theories – the total cancellation of ultraviolet divergences in the relations originating from the supersymmetry. Although it is expected to be so in general, there is still no explicit proof based on the Ward identities. Therefore it is convenient to demonstrate this mechanism at least in a particular case of the MSSM tree-level relation $`m_h^2+m_H^2m_A^2m_Z^2=0`$. The rule is renormalised using the usual diagrammatic technique. The only considered one-loop graphs are those involving top and supertop loops because the expected magnitude of the correction is the largest (taking into account the observation of that no significant contribution descends from the chargino-neutralino sector). The original results of and are recovered performing the whole computation in the unitary gauge; the explicit mechanism of the divergence cancellation is shown. The finite part is discussed in detail in the case of no mixing in the supertop sector.
## Appendix A Relevant Feynman graphs
This Appendix contains all the graphs discussed in previous sections. They are divided into two main subgroups - diagrams with quarks in loops and diagrams with the corresponding SUSY-partners. Each group consists of several types of graphs; the notation is usual and (perhaps) self-explanatory. The special symbols are defined as follows:
* $`\mathrm{C}_{\mathrm{uv}}\epsilon ^1\gamma _\mathrm{e}+\mathrm{ln}4\pi `$ denotes the ”divergent” part of a graph; here $`2\epsilon =4d`$; $`d`$ is the noninteger dimension used in the dimensional regularisation procedure; $`\gamma _\mathrm{e}`$ is the Euler-Mascheroni constant.
* $`B_X^{f_1f_2..}`$, $`T_X^{f_1f_2..}`$ and $`S_X^{f_1f_2..}`$ denote self-energies descending from the Feynman graphs (usualy called blobs, tadpoles and seagulls) with external lines $`X`$ and internal $`f_1,f_2,..`$.
* $`D_q^{f_1f_2}(x)m_{f_1}^2(1x)+m_{f_2}^2xq^2x(1x)`$ is the common factor arising from the regularisation prescription (see ) for UV-divergent graphs; $`q`$ is the momentum of the incoming (and outgoing) particle; in the on-shell scheme $`q^2=m_X^2`$.
* The constants $`g_{f_1f_2f_3..}`$ denote the numerical parts of the corresponding vertices . The non-number parts are contained in the structure of integrands.
* For the sake of brevity the overall colour factor $`N_c=3`$ is suppressed but must be included to obtain the correct results.
### A.1 Graphs with top quark loops
### I. Scalar and pseudoscalar self-energies:
$$\text{}B_X^t(q);Xh,H,A$$
$`B_X^t(q)`$ $`=`$ $`ig_{Xtt}^2\mu ^{2\epsilon }{\displaystyle \frac{\mathrm{d}^dk}{(2\pi )^d}\mathrm{Tr}\frac{i}{k/m_t}\frac{i}{k/q/m_t}};X=h,H`$
$`B_A^t(q)`$ $`=`$ $`ig_{Xtt}^2\mu ^{2\epsilon }{\displaystyle \frac{\mathrm{d}^dk}{(2\pi )^d}\mathrm{Tr}\gamma _5\frac{i}{k/m_t}\gamma _5\frac{i}{k/q/m_t}}`$
Using the routine calculational procedure (see for instance ) the results are
$`B_X^t(q)`$ $`=`$ $`g_{Xtt}^2{\displaystyle \frac{1}{4\pi ^2}}\{\mathrm{C}_{\mathrm{uv}}(3m_t^2{\displaystyle \frac{q^2}{2}})+(m_t^2{\displaystyle \frac{q^2}{6}})+`$ (19)
$`+`$ $`{\displaystyle _0^1}\mathrm{d}x\mathrm{ln}{\displaystyle \frac{D_q^{tt}(x)}{\mu ^2}}[3m_t^2+3q^2x(1x)]\};X=h,H`$
$`B_A^t(q)`$ $`=`$ $`g_{Att}^2{\displaystyle \frac{1}{4\pi ^2}}\{\mathrm{C}_{\mathrm{uv}}(m_t^2+{\displaystyle \frac{q^2}{2}})(m_t^2{\displaystyle \frac{q^2}{6}})+`$
$`+`$ $`{\displaystyle _0^1}\mathrm{d}x\mathrm{ln}{\displaystyle \frac{D_q^{tt}(x)}{\mu ^2}}[m_t^23q^2x(1x)]\}`$
The corresponding couplings are
$$g_{htt}=\frac{igm_t\mathrm{cos}\alpha }{2m_W\mathrm{sin}\beta };g_{Htt}=\frac{igm_t\mathrm{sin}\alpha }{2m_W\mathrm{sin}\beta };g_{Att}=\frac{gm_t}{2m_W}\mathrm{cot}\beta $$
(20)
### II. Z-boson self-energy graph (vacuum polarisation tensor):
$$\text{}B_Z^t(q)^{\mu \nu }$$
$$B_Z^t(q)^{\mu \nu }=ig_{Ztt}^2\mu ^{2\epsilon }\frac{\mathrm{d}^dk}{(2\pi )^d}\mathrm{Tr}\frac{i\gamma ^\mu K(\gamma )}{k/m_t}\frac{i\gamma ^\nu K(\gamma )}{k/q/m_t};$$
where
$$K(\gamma )(\epsilon _L^t+\epsilon _R^t)I_4(\epsilon _L^t\epsilon _R^t)\gamma _5.$$
The constants $`\epsilon _L^t`$, $`\epsilon _R^t`$ are defined as follows (in general $`\epsilon ^f=T_3Q_f\mathrm{sin}^2\theta _W`$)
$$\epsilon _L^t\frac{1}{2}q_t\mathrm{sin}^2\theta _W;\epsilon _R^tq_t\mathrm{sin}^2\theta _W$$
(21)
The coupling $`g_{Ztt}`$ is
$$g_{Ztt}=\frac{ig}{2\mathrm{c}\mathrm{o}\mathrm{s}\theta _W}$$
(22)
Performing the usual steps the result becomes
$`B_Z^t(q)^{\mu \nu }`$ $`=`$ $`g_{Ztt}^2\{{\displaystyle \frac{1}{4\pi ^2}}g^{\mu \nu }{\displaystyle _0^1}\mathrm{d}x\mathrm{ln}{\displaystyle \frac{D_q^{tt}(x)}{\mu ^2}}({\displaystyle \frac{m_t^2}{2}})`$
$`+`$ $`{\displaystyle \frac{1}{4\pi ^2}}\mathrm{C}_{\mathrm{uv}}\left[{\displaystyle \frac{m_t^2}{2}}g^{\mu \nu }+{\displaystyle \frac{2}{3}}(\epsilon _{L}^{t}{}_{}{}^{2}+\epsilon _{R}^{t}{}_{}{}^{2})(g^{\mu \nu }q^2q^\mu q^\nu )\right]+`$
$``$ $`{\displaystyle \frac{1}{2\pi ^2}}{\displaystyle _0^1}\mathrm{d}x\mathrm{ln}{\displaystyle \frac{D_q^{tt}(x)}{\mu ^2}}2(\epsilon _{L}^{t}{}_{}{}^{2}+\epsilon _{R}^{t}{}_{}{}^{2})(q^\mu q^\nu g^{\mu \nu }q^2)x(1x)\}`$
### III. Tadpoles involving top quark loop:
In general there are eight graphs to be considered in this paragraph. They are the following (here $`X=h`$,$`H`$,$`A`$ and $`S=h`$,$`H`$):
$$\text{}T_X^{St}(q)\text{}T_Z^{St}(q)^{\mu \nu }$$
Fortunately many of them cancel because of a nice property of the corresponding couplings ”sitting” in the upper vertex
$`g_{hhh}+g_{HHh}+g_{ZZh}=0`$
$`g_{hhH}+g_{HHH}+g_{ZZH}=0`$ (24)
The remaining graphs are evaluated in the usual way:
$$T_A^{St}(q)=ig_{AAS}g_{ttS}\frac{i}{m_S^2}\mu ^{2\epsilon }\frac{\mathrm{d}^dk}{(2\pi )^d}\mathrm{Tr}\frac{i}{k/m_t};$$
This d-dimensional integration is already easy to handle; we get
$$T_A^{St}(q)=4g_{AAS}g_{ttS}\frac{m_t^3}{m_S^2}\frac{1}{16\pi ^2}\left(\mathrm{C}_{\mathrm{uv}}+1\mathrm{ln}\frac{m_t^2}{\mu ^2}\right);$$
(25)
The couplings $`g_{htt}`$ and $`g_{htt}`$ are already written in (20); the remaining constants are
$$g_{AAh}=\frac{igm_Z}{2\mathrm{c}\mathrm{o}\mathrm{s}\theta _W}\mathrm{cos2}\beta \mathrm{sin}(\alpha +\beta );g_{AAH}=\frac{igm_Z}{2\mathrm{c}\mathrm{o}\mathrm{s}\theta _W}\mathrm{cos2}\beta \mathrm{cos}(\alpha +\beta )$$
(26)
At the end of this subsection note that the coupling constants used in this article can be found for example in and . In the case of supertops one must transform the rules in from the $`LR`$ basis to the supertop mass-diagonal basis $`12`$; this procedure is well described in the cited paper.
### A.2 Graphs with supertop loops
### I. Diagrams of the first type:
There are 7 graphs to be investigated in this cathegory, namely
$$\text{}B_X^{\stackrel{~}{t}_i\stackrel{~}{t}_j}(q);Xh,H,A;i,j=1,2$$
$$B_X^{\stackrel{~}{t}_i\stackrel{~}{t}_j}(q)=ig_{X\stackrel{~}{t}_i\stackrel{~}{t}_j}g_{X\stackrel{~}{t}_j\stackrel{~}{t}_i}\mu ^{2\epsilon }\frac{\mathrm{d}^dk}{(2\pi )^d}\frac{i}{k^2m_{\stackrel{~}{t}_i}^2}\frac{i}{(kq)^2m_{\stackrel{~}{t}_j}^2}(2\delta _{ij})$$
Using again , the last expression can be simplified to the final form
$$B_X^{\stackrel{~}{t}_i\stackrel{~}{t}_j}(q)=g_{X\stackrel{~}{t}_i\stackrel{~}{t}_j}g_{X\stackrel{~}{t}_j\stackrel{~}{t}_i}(2\delta _{ij})\frac{1}{16\pi ^2}\left[\mathrm{C}_{\mathrm{uv}}_0^1dx\mathrm{ln}\frac{D_q^{\stackrel{~}{t}_i\stackrel{~}{t}_j}(x)}{\mu ^2}\right]$$
(27)
Note that the factor $`(2\delta _{ij})`$ counts the number of nonequivalent contractions. To finish this paragraph it is necessary to specify the couplings; the vertices involving scalars are symmetric with respect to $`ij`$
$`g_{h\stackrel{~}{t}_1\stackrel{~}{t}_1}`$ $`=`$ $`{\displaystyle \frac{igm_Z}{\mathrm{cos}\theta _W}}\mathrm{sin}(\alpha +\beta )\left(\epsilon _L^t\mathrm{cos}^2\theta _t\epsilon _R^t\mathrm{sin}^2\theta _t\right){\displaystyle \frac{igm_t^2\mathrm{cos}\alpha }{m_W\mathrm{sin}\beta }}`$
$`{\displaystyle \frac{igm_t\mathrm{sin2}\theta _t}{2m_W\mathrm{sin}\beta }}\left(A_tm_6\mathrm{cos}\alpha \mu \mathrm{sin}\alpha \right)`$
$`g_{h\stackrel{~}{t}_2\stackrel{~}{t}_2}`$ $`=`$ $`{\displaystyle \frac{igm_Z}{\mathrm{cos}\theta _W}}\mathrm{sin}(\alpha +\beta )\left(\epsilon _L^t\mathrm{sin}^2\theta _t\epsilon _R^t\mathrm{cos}^2\theta _t\right){\displaystyle \frac{igm_t^2\mathrm{cos}\alpha }{m_W\mathrm{sin}\beta }}`$
$`+{\displaystyle \frac{igm_t\mathrm{sin2}\theta _t}{2m_W\mathrm{sin}\beta }}\left(A_tm_6\mathrm{cos}\alpha \mu \mathrm{sin}\alpha \right)`$
$`g_{h\stackrel{~}{t}_1\stackrel{~}{t}_2}`$ $`=`$ $`{\displaystyle \frac{igm_Z}{\mathrm{cos}\theta _W}}\mathrm{sin}(\alpha +\beta )\left(\epsilon _L^t+\epsilon _R^t\right)\mathrm{sin}\theta _t\mathrm{cos}\theta _t`$ (28)
$`{\displaystyle \frac{igm_t\mathrm{cos2}\theta _t}{2m_W\mathrm{sin}\beta }}\left(A_tm_6\mathrm{cos}\alpha \mu \mathrm{sin}\alpha \right)`$
$`g_{H\stackrel{~}{t}_1\stackrel{~}{t}_1}`$ $`=`$ $`{\displaystyle \frac{igm_Z}{\mathrm{cos}\theta _W}}\mathrm{cos}(\alpha +\beta )\left(\epsilon _L^t\mathrm{cos}^2\theta _t\epsilon _R^t\mathrm{sin}^2\theta _t\right){\displaystyle \frac{igm_t^2\mathrm{sin}\alpha }{m_W\mathrm{sin}\beta }}`$
$`{\displaystyle \frac{igm_t\mathrm{sin2}\theta _t}{2m_W\mathrm{sin}\beta }}\left(A_tm_6\mathrm{sin}\alpha +\mu \mathrm{cos}\alpha \right)`$
$`g_{H\stackrel{~}{t}_2\stackrel{~}{t}_2}`$ $`=`$ $`{\displaystyle \frac{igm_Z}{\mathrm{cos}\theta _W}}\mathrm{cos}(\alpha +\beta )\left(\epsilon _L^t\mathrm{sin}^2\theta _t\epsilon _R^t\mathrm{cos}^2\theta _t\right){\displaystyle \frac{igm_t^2\mathrm{sin}\alpha }{m_W\mathrm{sin}\beta }}`$
$`+{\displaystyle \frac{igm_t\mathrm{sin2}\theta _t}{2m_W\mathrm{sin}\beta }}\left(A_tm_6\mathrm{sin}\alpha +\mu \mathrm{cos}\alpha \right)`$
$`g_{H\stackrel{~}{t}_1\stackrel{~}{t}_2}`$ $`=`$ $`{\displaystyle \frac{igm_Z}{\mathrm{cos}\theta _W}}\mathrm{cos}(\alpha +\beta )\left(\epsilon _L^t+\epsilon _R^t\right)\mathrm{sin}\theta _t\mathrm{cos}\theta _t`$
$`{\displaystyle \frac{igm_t\mathrm{sin2}\theta _t}{2m_W\mathrm{sin}\beta }}\left(A_tm_6\mathrm{sin}\alpha +\mu \mathrm{cos}\alpha \right)`$
while vertices with pseudoscalar $`A`$ are antisymmetric:
$`g_{A\stackrel{~}{t}_1\stackrel{~}{t}_2}=g_{A\stackrel{~}{t}_2\stackrel{~}{t}_1}`$ $`=`$ $`{\displaystyle \frac{gm_t}{2m_W\mathrm{sin}\beta }}\left(A_tm_6\mathrm{cos}\beta \mu \mathrm{sin}\beta \right)`$
$`g_{A\stackrel{~}{t}_1\stackrel{~}{t}_1}=g_{A\stackrel{~}{t}_2\stackrel{~}{t}_2}`$ $`=`$ $`0`$ (29)
### II. Z-boson self-energy graphs with looping superquarks:
The graphs relevant to this paragraph are all of the type
$$\text{}B_Z^{\stackrel{~}{t}_i\stackrel{~}{t}_j}(q)^{\mu \nu };j=1,2$$
$$B_Z^{\stackrel{~}{t}_i\stackrel{~}{t}_j}(q)^{\mu \nu }=ig_{Z\stackrel{~}{t}_i\stackrel{~}{t}_j}^2\mu ^{2\epsilon }\frac{\mathrm{d}^dk}{(2\pi )^d}\frac{i(2kp)^\mu }{k^2m_{\stackrel{~}{t}_1}^2}\frac{i(2kp)^\nu }{(kq)^2m_{\stackrel{~}{t}_1}^2}$$
As in the previous cases, after some algebra one obtains
$$B_Z^{\stackrel{~}{t}_i\stackrel{~}{t}_j}(q)^{\mu \nu }=g_{Z\stackrel{~}{t}_i\stackrel{~}{t}_j}^2\frac{1}{16\pi ^2}\{\mathrm{C}_{\mathrm{uv}}[\frac{1}{3}(q^\mu q^\nu g^{\mu \nu }q^2)+(m_{\stackrel{~}{t}_i}^2+m_{\stackrel{~}{t}_j}^2)g^{\mu \nu }]+$$
$$+g^{\mu \nu }(m_{\stackrel{~}{t}_i}^2+m_{\stackrel{~}{t}_j}^2\frac{q^2}{3})_0^1\mathrm{d}x[q^\mu q^\nu (12x)^2+2g^{\mu \nu }D_q^{\stackrel{~}{t}_i\stackrel{~}{t}_j}(x)]\mathrm{ln}\frac{D_q^{\stackrel{~}{t}_i\stackrel{~}{t}_j}(x)}{\mu ^2}\}$$
### III. Tadpoles involving supertop loop:
In general there is again many diagrams belonging to this paragraph; as in the previous section the relations (A)-(A) ensure that most of the graphs cancel. The remaining are (here $`S=h`$,$`H`$ and $`i=1,2`$):
$$\text{}T_A^{S\stackrel{~}{t}_i}(q)$$
The contributions coming from these graphs are
$$T_A^{S\stackrel{~}{t}_i}(q)=ig_{AAS}g_{\stackrel{~}{t}_i\stackrel{~}{t}_iS}\frac{i}{m_S^2}\mu ^{2\epsilon }\frac{\mathrm{d}^dk}{(2\pi )^d}\frac{i}{k^2m_{\stackrel{~}{t}_i}^2};$$
which gives after regularisation
$$T_A^{S\stackrel{~}{t}_i}(q)=g_{AAS}g_{\stackrel{~}{t}_i\stackrel{~}{t}_iS}\frac{m_{\stackrel{~}{t}_i}^2}{m_S^2}\frac{1}{16\pi ^2}\left(\mathrm{C}_{\mathrm{uv}}+1\mathrm{ln}\frac{m_{\stackrel{~}{t}_i}^2}{\mu ^2}\right);$$
(30)
The necessary coupling constants are written in (26) and (28)
### IV. Seagull graphs:
Due to the presence of quadrilinear vertices involving two Higgses and two superquarks there is an additional sort of graphs in this section – the so-called seagull graphs which look as ($`S=h`$,$`H`$,$`A`$ and $`i=1,2`$)
$$\text{}S_S^{\stackrel{~}{t}_i}(q);\text{}S_Z^{\stackrel{~}{t}_i}(q)^{\mu \nu }$$
Note that similar graphs can not appear in the fermion sector because the quadrilinear vertices involving fermions have mass dimensions $`>4`$ and they would cause nonrenormalisability of the theory. The contributions originating from these graphs can be written as
$$S_S^{\stackrel{~}{t}_i}(q)=ig_{SS\stackrel{~}{t}_i\stackrel{~}{t}_i}\mu ^{2\epsilon }\frac{\mathrm{d}^dk}{(2\pi )^d}\frac{i}{k^2m_{\stackrel{~}{t}_i}^2};$$
$$S_Z^{\stackrel{~}{t}_i}(q)^{\mu \nu }=ig_{ZZ\stackrel{~}{t}_i\stackrel{~}{t}_i}\mu ^{2\epsilon }g^{\mu \nu }\frac{\mathrm{d}^dk}{(2\pi )^d}\frac{i}{k^2m_{\stackrel{~}{t}_i}^2};$$
which after regularisation gives
$$S_S^{\stackrel{~}{t}_i}(q)=g_{SS\stackrel{~}{t}_i\stackrel{~}{t}_i}m_{\stackrel{~}{t}_i}^2\frac{1}{16\pi ^2}\left(\mathrm{C}_{\mathrm{uv}}+1\mathrm{ln}\frac{m_{\stackrel{~}{t}_i}^2}{\mu ^2}\right);$$
(31)
$$S_S^{\stackrel{~}{t}_i}(q)^{\mu \nu }=g_{ZZ\stackrel{~}{t}_i\stackrel{~}{t}_i}m_{\stackrel{~}{t}_i}^2g^{\mu \nu }\frac{1}{16\pi ^2}\left(\mathrm{C}_{\mathrm{uv}}+1\mathrm{ln}\frac{m_{\stackrel{~}{t}_i}^2}{\mu ^2}\right);$$
In general there are 8 graphs to deal with. Fortunately the coupling constants can be nicely summed up so that
$$g_{hh\stackrel{~}{t}_1\stackrel{~}{t}_1}+g_{HH\stackrel{~}{t}_1\stackrel{~}{t}_1}+g_{hh\stackrel{~}{t}_2\stackrel{~}{t}_2}+g_{HH\stackrel{~}{t}_2\stackrel{~}{t}_2}g_{AA\stackrel{~}{t}_1\stackrel{~}{t}_1}g_{AA\stackrel{~}{t}_2\stackrel{~}{t}_2}=g_{GG\stackrel{~}{t}_2\stackrel{~}{t}_2}+g_{GG\stackrel{~}{t}_1\stackrel{~}{t}_1}$$
(32)
where
$`g_{GG\stackrel{~}{t}_1\stackrel{~}{t}_1}`$ $`=`$ $`{\displaystyle \frac{ig^2}{2\mathrm{c}\mathrm{o}\mathrm{s}^2\theta _W}}\mathrm{cos2}\beta \left(\epsilon _L^t\mathrm{cos}^2\theta _t\epsilon _R^t\mathrm{sin}^2\theta _t\right){\displaystyle \frac{ig^2m_t^2}{2m_W^2}}`$
$`g_{GG\stackrel{~}{t}_2\stackrel{~}{t}_2}`$ $`=`$ $`{\displaystyle \frac{ig^2}{2\mathrm{c}\mathrm{o}\mathrm{s}^2\theta _W}}\mathrm{cos2}\beta \left(\epsilon _L^t\mathrm{sin}^2\theta _t\epsilon _R^t\mathrm{cos}^2\theta _t\right){\displaystyle \frac{ig^2m_t^2}{2m_W^2}}`$ (33)
Note that these constants are exactly the couplings of the would-be Goldstone boson (which is within U-gauge absent). The last unspecified parameters are the couplings $`g_{ZZ\stackrel{~}{t}_i\stackrel{~}{t}_i}`$ :
$`g_{ZZ\stackrel{~}{t}_1\stackrel{~}{t}_1}`$ $`=`$ $`{\displaystyle \frac{2ig^2}{\mathrm{cos}^2\theta _W}}\left(\epsilon _{L}^{t}{}_{}{}^{2}\mathrm{cos}^2\theta _t+\epsilon _{R}^{t}{}_{}{}^{2}\mathrm{sin}^2\theta _t\right)`$
$`g_{ZZ\stackrel{~}{t}_2\stackrel{~}{t}_2}`$ $`=`$ $`{\displaystyle \frac{2ig^2}{\mathrm{cos}^2\theta _W}}\left(\epsilon _{L}^{t}{}_{}{}^{2}\mathrm{sin}^2\theta _t+\epsilon _{R}^{t}{}_{}{}^{2}\mathrm{cos}^2\theta _t\right)`$ (34)
## Appendix B Some useful relations and comments
This Appendix is devoted to several clarifications necessary to make the text more self-consistent.
The first note refers to the MSSM itself. The full-range discussion of the relevant part of the MSSM classical lagrangian is obviously out of the scope of this paper. There are several comprehensive works in the literature that can be used for this purpose, namely , , or . The notation is similar to that used in .
The rest of the Apppendix contains some comments on technical details of the computation. To enable the reader follow the steps described above it is necessary to write down several not so well known tree-relations often used during the computation. First of them,
$$m_h^2\mathrm{cos}^2\alpha +m_H^2\mathrm{sin}^2\alpha m_A^2\mathrm{cos}^2\beta m_Z^2\mathrm{sin}^2\beta =0$$
(35)
can be derived by utilising (1) and relations
$$\mathrm{sin2}\alpha =\left(\frac{m_h^2+m_H^2}{m_h^2m_H^2}\right)\mathrm{sin2}\beta ;\mathrm{cos2}\alpha =\left(\frac{m_A^2m_Z^2}{m_h^2m_H^2}\right)\mathrm{cos2}\beta $$
(36)
that can be found for example in . Note only that the parameters $`\alpha `$ and $`\beta `$ are the mixing angles in the scalar and pseudoscalar parts of the Higgs sector. These relations are also very handy if we want to express (tree) Higgs masses in terms of $`m_Z^2`$, $`\alpha `$ and $`\beta `$ dealing with the factors $`m_h^2`$ or $`m_H^2`$ coming from the tadpoles in (A.1) and (A.2).
The next thing to be clarified is the role of the parameters $`A_tm_6`$ and $`\mu `$ in the supertop mass-squared matrix. This matrix in the $`LR`$ basis looks (see )
$$M_{\stackrel{~}{t}_{L,R}}^2=\left(\begin{array}{cc}A& B\\ B& C\end{array}\right)$$
(37)
where
$`A`$ $`=`$ $`M_\mathrm{q}^2+m_Z^2\mathrm{cos2}\beta \epsilon _L^t+m_t^2`$
$`B`$ $`=`$ $`m_t\left(A_tm_6+\mu \mathrm{cot}\beta \right)`$
$`C`$ $`=`$ $`M_\mathrm{u}^2+m_Z^2\mathrm{cos2}\beta \epsilon _R^t+m_t^2`$
is the usual parametrisation of its entries. (Note that the constants $`M_\mathrm{q}^2`$ and $`M_\mathrm{u}^2`$ are the so-called soft-SUSY breaking terms which in general split the masses of supertops and shift them high above $`m_t`$, .) The eigenvalues of this matrix can be easily derived in the form
$$m_{\stackrel{~}{t}_{1,2}}^2=\frac{1}{2}\left[A+C\pm \sqrt{(A+C)^24(ACB^2)}\right]$$
(38)
The mixing (diagonalising) angle $`\theta _t`$ is then defined by
$$\mathrm{tan}\theta _t=\frac{2B}{AC}$$
which can be rewritten in terms of the eigenvalues (38) as follows
$$\mathrm{sin2}\theta _t=\frac{2m_t(A_tm_6+\mu \mathrm{cot}\beta )}{m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2}$$
(39)
This relation connects the off-diagonal entries in the supertop mass-squared matrix with the magnitude of the supertop mass-split and the mixing angle $`\theta _t`$. Assuming now that the supertop mass-squared does not exceed the top-scale too much the assumption of a small mixing can be recast in the form $`(A_tm_6+\mu \mathrm{cot}\beta )m_t`$.
I would like to thank to my supervisor Prof. Jiří Hořejší for all the comments and helpful suggestions during the long period of writing this article. The work has been partially supported by the grant GAČR No. 202/98/0506.
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# OCHA-PP-158 A Comment on Non-Archimedean Character of Quantum Buoyancy
## I Introduction
We believe that black hole is a thermodynamical object and has entropy in a sense. This belief is based on analogy with ordinary thermodynamics; mathematical relationship between black hole mechanics and ordinary law of thermodynamics, and the existence of thermal radiation from a black hole.
In order to transmute the mathematical relationship into physical one, that is, black hole mechanics into black hole thermodynamics, understanding of statistical origin of black hole entropy has progressed from the various points of view . It is any theory of quantum gravity which controls a measure of the density of the states that we need to understand genuinely statistical black hole entropy without divergent quantities. On the other hand, black hole entropy can be derived in metrical theories of gravity by a classical method such as Noether charge method or by semiclassical methods such as Euclidean path integral method , though these methods by no means count quantum degrees of freedom that are responsible for black hole entropy. Since any successful quantum gravity theory comes to a corresponding metrical theory of gravity in suitable low energy limit, we expect that the value of black hole entropy obtained by classical or semiclassical methods should be also derived by counting quantum degrees of freedom. Thus, we may regard success of statistical derivation of black hole entropy as a benchmark test of a proposed quantum gravity theory.
On the other hand, it is also necessary to understand black hole entropy better, even at the level of thermodynamics. For instance, it is the second law that characterizes the peculiar property of entropy in ordinary thermodynamics because of its referring to a direction in time. Therefore it is very important for understanding of black hole entropy to establish the second law of black hole thermodynamics.
Since we cannot regard a black hole as an isolated system owing to the universal interaction with ordinary matter outside the black hole by gravity, any second law of black hole thermodynamics should refer to total entropy of self-gravitating system including black holes. Therefore we are led to the generalized second law (GSL) of black hole thermodynamics, which asserts that in any process, the generalized entropy
$$S_G:=S_{BH}+S_M=\frac{1}{T_{BH}}\left(\frac{\kappa A_{BH}}{8\pi }\right)+S_M$$
(1)
never decreases, where $`S_{BH}`$ and $`S_M`$ denote the entropy of the black holes and that of ordinary matter outside the black hole, and then $`\kappa `$, $`A_{BH}`$ and $`T_{BH}`$ are the surface gravity, area of the event horizon and temperature of the black hole, respectivelySomeone may doubt additivity of entropy of self-gravitating system, due to long-range nature of gravity. Instead of no conclusive argument about the additivity, we assume the validity of the additivity. See for argument validating the additivity. . The validity of the GSL is essential for the consistency of black hole thermodynamics and for the interpretation of the horizon area as representing the physical entropy of a black hole, because it is nothing but the ordinary second law for self-gravitating systems containing black holes. Thus, the GSL is a cornerstone of black hole thermodynamics.
Although an explicit general proof of the GSL has not been given until now, the various attempts for special cases have been performed . Considering a process which transfers an infinitesimal energy $`\delta E`$ and entropy $`\delta S`$ in the external region into the black hole adiabatically, we obtain the change in the total entropy $`\delta S_G=\delta E/T_{BH}\delta S`$.
In classical theory, we may argue as follows. Since a black hole can classically export nothing outside the horizon, it is natural to give zero temperature $`T_{BH}=0`$ to the black hole. Therefore, if it were so, by dominance of the first term in Eq.(1), the GSL in classical theory should be no more than the second law for the black hole entropy alone and then it would amount to the area increasing law in black hole mechanics $`\delta A_{BH}>0`$ which holds by energy condition $`\delta E>0`$ .
However, it is awkward to assign $`T_{BH}=0`$ to the black hole, since the black hole entropy or the change in it becomes divergent and ill-defined. Thus, the physical analogy appears end in classical theory. In order to have non-zero black hole temperature and well-defined black hole entropy, it is indeed essential to incorporate quantum effects even semiclassically. Due to the breakdown of the energy condition of quantum fields, black holes can radiate and acquire non-zero temperature, and then the thermodynamic quantities of a black hole can be fixed as $`T_{BH}=\kappa /2\pi `$ and $`S_{BH}=A_{BH}/4`$ . Therefore, it is important to investigate the validity of the GSL by consistent arguments with taking account of quantum effects.
An observer accelerating with acceleration $`a`$ detects isotropic thermal radiation with temperature $`T_U=\mathrm{}a/2\pi `$ by the Unruh radiation (acceleration radiation) . An object suspended near a black hole is accelerated by virtue of its being prevented from following a geodesic. Unruh and Wald suggested that this object will likewise see Unruh radiance. Since its acceleration (i.e. temperature) varies with distance from the horizon, they surmised that the object will be subject to a buoyant force by the acceleration radiation fluid and the buoyancy affects the energetics of a process which exchange entropy and energy between the black hole and outer matter. They concluded that quantum buoyancy is sufficient by itself to protect the GSL.
Recently, Bekenstein reconsidered the nature of acceleration radiation and its implication on the GSL . He pointed out that the wave nature (diffractive effect) of the acceleration radiation cannot be neglected in the case that the size of the object lowered toward the black hole is smaller than a typical wavelength of the acceleration radiation and that the fluid approximation of the acceleration radiation is invalid. For such a case, he estimated the buoyant force as a wave scattering process and found that the buoyant force as a wave scattering process is weaker than in the fluid approximation. Therefore, the diffractive effect alters energetics of exchange process of the entropy and energy compared with that in fluid picture, and then the quantum buoyancy is insufficient by itself to protect the GSL. A breakdown of the GSL in the existing physics leads us to a new physics, such as an entropy bound for matter, if we take granted that the GSL holds. Thus, the question of the validity of the GSL is still be opened even in a simple gedankenexperiment.
In this letter, we observe that if a massless scalar field exists, the quantum buoyancy is sufficient to protect the GSL, rather strengthen the validity of the GSL, even though we take account of the wave nature of the acceleration radiation.
## II a gedankenexperiment
In this section, we specify a gedankenexperiment investigated in this letter and review two independent reasonings for the GSL to hold.
We consider a static black hole which area of the event horizon is $`𝒜`$ and a box of proper height $`b`$ and geometrical cross-sectional area $`A`$. Far from the black hole, the box is filled with matter, so that the total energy of the box and contents is $`E_0`$ and its total entropy $`S_0`$. Subsequently, the box is lowered adiabatically toward the black hole by a weightless rope to some height $`l`$ that is the proper distance between the horizon and the center of mass of the box. And then, the box and contents are released and dropped into the black hole.
Because of the process to be adiabatic, the total entropy of the box and contents remains constant. Therefore, the change in the total entropy becomesHere we implicitly assume that processes after the box released preserves, the total energy and entropy of the box and contents.
$$\mathrm{\Delta }S_G=\mathrm{\Delta }S_{BH}S_0=\frac{\mathrm{\Delta }M(l)}{T_{BH}}S_0,$$
(2)
where $`T_{BH}`$ is the (non-zero) black hole temperature.
On the other hand, the energy $`\mathrm{\Delta }M(l)`$ delivered to the black hole decreases during the lowering process, because the gravitational energy of the box and contents is lost by the work against the tension of the rope. As denoting the redshift factor $`\xi (l)`$, we obtain the equation,
$`\mathrm{\Delta }M(l)`$ $`=`$ $`E_0+(\text{work done by the rope})=E_0+W_{\mathrm{}}(l)=E_0+{\displaystyle _{\mathrm{}}^l}\left(F_{\mathrm{}}^{}{}_{}{}^{G}\right)𝑑l`$ (3)
$`=`$ $`E_0+E_0\left[\xi (l)1\right]=E_0\xi (l),`$ (4)
where we use the relation $`F_{\mathrm{}}^{}{}_{}{}^{G}dl=dE_{\mathrm{}}=d\left(E_0\xi \right)`$, that is derived from $`E_{\mathrm{}}=E_0\xi `$. Thus, the energy delivered to the black hole is “redshifted away”, due to the negative gravitational potential.
Therefore, the change in the total entropy is
$$\mathrm{\Delta }S_G=\frac{E_0}{T_{BH}}\xi (l)S_0=\frac{E_0}{T(l)}S_0,$$
(5)
where $`T(l):=T_{BH}/\xi (l)`$ means the locally measured temperature of the black hole atmosphere. Because, if the box can be close to the horizon without limit, $`\xi `$ can be arbitrary small near the horizon, we can make the value of $`\mathrm{\Delta }S_G`$ negative at will.
If we take granted that the GSL holds, then we need any mechanism which prevents the box from the horizon. At present, there exist two reasonings: one is invoking to an entropy bound for matter and the other makes use of the buoyant force by the black hole atmosphere. It is essential to recognize that the box must have a finite size which is greater than its Compton wavelength.
The argument of the first reasoning invoking an entropy bound is as follows: The finiteness of the box size imposes a constraint, $`lb/2`$, that is,
$$\xi (l)\kappa l=2\pi T_{BH}l\pi T_{BH}b,$$
(6)
because the bottom of the box cannot touch the horizon. If we premise the validity of the GSL and the energetics Eq.(4), we need an entropy bound for matter
$$S\pi Eb.$$
(7)
Thus, the entropy of any matter in this case is bounded above by its energy and size. Since, obviously, the size $`b/2`$ is greater than its gravitational radius $`r_g=2E`$, we obtain
$$S2\pi Er_g\frac{4\pi r_g^2}{4}.$$
(8)
Thus the maximum entropy of any matter is bounded above by its gravitational radius and the saturated state is attained by the black hole state. This relation is called holographic bound, which the validity of Eq.(8) is open problem and has actively been discussed in the different viewpoint, holographic principle . Even though it is finally true that there exists the entropy bound for matter or the holographic bound, it is important to investigate to what degree the GSL is protected by the known physics and whether the validity of the GSL implies the entropy bound or the holographic bound.
Another reasoning invoking the known physics makes use of quantum effect of matter field outside black holes, that is, the buoyant force by the black hole atmosphere, which has been neglected in the argument of the first reasoning. We may start with two main working hypothesis ;
The black hole atmosphere is describable by radiation fluid of a unconstrained thermal matter which is defined to be the state of matter that maximizes entropy density at a fixed energy density and the radiation fluid has the locally measured temperature $`T(l)`$.
The buoyant force on the box exerted by the black hole atmosphere is equal to the pressure gradient of the radiation fluid of unconstrained thermal matter.
The assumption A1 means that the Gibbs-Duhem relation holds,
$$\{\begin{array}{cc}\rho _{rad}+P_{rad}T(l)s_{rad}=0\hfill & \\ d\rho _{rad}=T(l)ds_{rad}\hfill & \end{array},$$
(9)
and by Eqs.(9) and $`T(l)=T_{BH}/\xi (l)`$, we obtain balance equation between gravitational force and pressure gradient force of the radiation fluid
$$\frac{d}{dl}(\xi P_{rad})=\rho _{rad}(l)\frac{d\xi }{dl}.$$
(10)
Using Eq.(10) and the assumption A2, we obtain “Archimedean principle”,
$$F_{\mathrm{}}^{}{}_{}{}^{B}=A\left[(\xi P_{rad})_{bottom}(\xi P_{rad})_{top}\right]=V\frac{d}{dl}(\xi P_{rad})=V\rho _{rad}(l)\frac{d\xi }{dl}.$$
(11)
Therefore, the work done by the total force $`F_{\mathrm{}}=F_{\mathrm{}}^{}{}_{}{}^{G}+F_{\mathrm{}}^{}{}_{}{}^{B}`$ becomes
$$W_{\mathrm{}}=_{\mathrm{}}^l\left(F_{\mathrm{}}\right)𝑑l=E_0\left[\xi (l)1\right]+V\xi (l)P_{rad}(l).$$
(12)
And then, the energy delivered into the black hole is
$`\mathrm{\Delta }M`$ $`=`$ $`E_0+W_{\mathrm{}}=\left[E_0+VP_{rad}(l)\right]\xi (l)`$ (13)
$`=`$ $`V\left[\rho _0\rho _{rad}+T(l)s_{rad}\right]\xi (l),`$ (14)
where $`\rho _0:=E_0/V`$ is the average energy density of the box and the contents. The change in the total entropy becomes
$$\mathrm{\Delta }S_G=V\left(\frac{\rho _0\rho _{rad}}{T(l)}+s_{rad}s_0\right),$$
(15)
where $`s_0:=S_0/V`$ is the average entropy density of the box and the contents.
The critical situation for the positivity of $`\mathrm{\Delta }S_G`$ is the case of minimizing $`\mathrm{\Delta }M`$,
$$0=\frac{d}{dl}W_{\mathrm{}}=F_{\mathrm{}}^{}{}_{}{}^{G}+F_{\mathrm{}}^{}{}_{}{}^{B}=V\left(\rho _0\rho _{rad}\right)\frac{d\xi }{dl},$$
(16)
so that, it is the most dangerous for the validity of the GSL when the box is dropped into the black hole at the floating point $`\rho _0=\rho _{rad}(l)`$.
Nevertheless the positivity of $`\mathrm{\Delta }S_G`$ holds by the definition of the radiation fluid, that is, we can show the validity of the GSL without invoking a new physics,
$$\mathrm{\Delta }S_GV\left(s_{rad}s_0\right)0,$$
(17)
where the last inequality follows the definition A1 of the radiation fluid, because of $`\rho _0=\rho _{rad}(l)`$ at the floating point.
Now we should check the validity of our assumptions, especially, the validity of the assumption A2. It is natural to think that if a typical wavelength of the acceleration radiation $`\lambda `$ is much bigger than the box size $`b`$, the assumption A2 is invalid due to the breakdown of the fluid picture, such as diffractive effect. Therefore, it is doubtful to consider that the fluid picture is still valid far from the black hole, such as $`b<\lambda T^1(l)`$.
## III The buoyant force by long wavelength scattering
Recently, Bekenstein pointed out the breakdown of the fluid picture far from the horizon .
Strictly speaking, the true pressure exerted on the surface of the box is given by integrating true stress tensor over the surface and the true stress tensor must be obtained by inclusion of the boundary condition of the surface. However, in the previous section, we estimated the pressure by the fictitious stress tensor, which means that the stress tensor is estimated by neglecting the surface, exclusive of the boundary condition. In order to estimate the true pressure, it is often useful to calculate the change in momentum flux on the surface and it is essential for calculating the change in the momentum flux to estimate the reflection coefficient, that is, to include the boundary condition on the surface. For example, a perfectly transparent glass is not exerted by photons, even though the momentum flux across the glass does not vanish.
Thus, we need to estimate the scattering cross section of the box for the acceleration radiation. If $`bR_H`$, where $`R_H`$ is the curvature radius at the horizon, then we can acquire a large local (Lorentz) frame including the target (the box) in which the target is at rest. Therefore, we can approximate the scattering process in the black hole spacetime by the scattering process in flat spacetime and at first estimate quantities in interest, such as the momentum transfer, in the local frame. A remained task is to transform quantities obtained in the local frame into ones in the global frame, that is, quantities as measured at infinity .
We calculate physical quantities in the long wavelength limit $`b\lambda =:2\pi /k`$, because we are especially interested in scattering phenomena in the situation that the fluid picture of the acceleration radiation is suspicious. In this limit, the differential cross section is indifferent to details such as the shape of the target. Hereafter we assume that the the shape of the target is spherically symmetric.
### A the buoyant force by dipole scattering
According to the above procedure, Bekenstein estimated the buoyant force exerted by dipole scattering. In order to estimate the buoyant force, we calculate the differential cross-section of the target object with the size $`b`$ by the dipole scattering. For the dipole scattering which transfers the incident wave with the wave vector $`\stackrel{}{k}`$ into the scattered wave with $`\stackrel{}{k}^{}`$ and preserves the magnitude of the momentum, $`k:=|\stackrel{}{k}|=|\stackrel{}{k}^{}|`$, we have
$$\frac{d\sigma }{d\mathrm{\Omega }^{}}=b^2\left(kb\right)^4F(\stackrel{}{n},\stackrel{}{n}^{}),$$
(18)
where $`\stackrel{}{n}`$ and $`\stackrel{}{n}^{}`$ are a pair of the unit vectors denoting the incident and scattering directions, $`\stackrel{}{n}:=\stackrel{}{k}/k`$ and $`\stackrel{}{n}^{}:=\stackrel{}{k}^{}/k`$, respectively. And $`F`$ is some dimensionless function, which, for example, is given by $`F(\stackrel{}{n},\stackrel{}{n}^{})=\pi ^2\left\{\frac{5}{8}\left[1+\mathrm{cos}^2\left(\stackrel{}{n}\stackrel{}{n}^{}\right)\right]\mathrm{cos}\left(\stackrel{}{n}\stackrel{}{n}^{}\right)\right\}`$ for electromagnetic scattering from a conducting sphere. The fourth order dependence of the cross-section on the wave vector is attributed to the fact that the dipole part is dominant in the scattering of the electromagnetic wave.
Given a distribution function of the incident wave as $`f(k)=1/\left[\mathrm{exp}\left(\mathrm{}k/T\right)1\right]`$, the incident momentum flux carried in the acceleration radiation in the vicinity of the wave vector $`\stackrel{}{k}=k\stackrel{}{n}`$ becomes
$`\stackrel{}{n}I(k,\stackrel{}{n})dkd\stackrel{}{n}`$ $`:=`$ $`(\mathrm{}\stackrel{}{k})f(k){\displaystyle \frac{d^3k}{(2\pi )^3}}=\stackrel{}{n}\mathrm{}k^3f(k){\displaystyle \frac{dkd\stackrel{}{n}}{(2\pi )^3}},`$ (19)
where $`d^3k=k^2dkd\stackrel{}{n}`$.
Because the fraction $`d\sigma /d\mathrm{\Omega }^{}`$ among the incident flux $`Idkd\stackrel{}{n}`$ is scattered into the direction $`\stackrel{}{n}^{}`$, we obtain the momentum transfer of the box in the local frame,
$`{\displaystyle \frac{d\stackrel{}{P}}{d\tau }}={\displaystyle 𝑑k𝑑\stackrel{}{n}𝑑\stackrel{}{n}^{}I(k,\stackrel{}{n})\frac{d\sigma }{d\mathrm{\Omega }^{}}(\stackrel{}{n}\stackrel{}{n}^{})},`$ (20)
where $`\tau `$ is time measured in the local frame, that is, the proper time of the target.
Since the acceleration radiation has the temperature gradient, the radiation going to the direction $`\stackrel{}{n}`$ hits on the box with the local temperature $`T(l,\stackrel{}{n})=T_{BH}/\xi (l,\stackrel{}{n})=\mathrm{}/2\pi [l+(\stackrel{}{e_l}\stackrel{}{n})b]=T(l)/[1+(\stackrel{}{e_l}\stackrel{}{n})b/l]`$, where $`\stackrel{}{e_l}`$ is the unit vector that is directed to the center of mass of the box from the black hole. Therefore, the buoyant force $`F_{\mathrm{}}^{}{}_{}{}^{Scatt}`$ by the dipole scattering in the global frame becomes
$$F_{\mathrm{}}^{}{}_{}{}^{Scatt}=\xi (l)\left|\frac{d\stackrel{}{P}}{d\tau }\right|(2\pi lT_{BH})\left[\frac{T(l)}{\mathrm{}}\right]^8b^6𝑑\stackrel{}{n}\left[1+\frac{b}{l}(\stackrel{}{e_l}\stackrel{}{n})\right]^8\frac{T_{BH}}{b}\left(\frac{b}{l}\right)^8,$$
(21)
where we neglect numerical factor. Thus, the buoyant force by the dipole scattering is proportional to the seventh power of the size $`b`$, not to the volume (non-Archimedean) and proportional to the eighth inverse power of the proper distance $`l`$ from the event horizon.
On the other hand, the buoyant force in fluid picture is estimated by,
$$F_{\mathrm{}}^{}{}_{}{}^{B}=V\rho _{rad}(l)\frac{d\xi }{dl}b^3[T(l)]^4T_{BH}\frac{T_{BH}}{b}\left(\frac{b}{l}\right)^4.$$
(22)
Thus, in the case of the dipole scattering, Archimedean character of buoyant force that the force is proportional to the volume of the box does not work for $`b\lambda `$ and the force rapidly decreases with the distance from the horizon than in the fluid picture.
Since it is possible to saturate the inequality Eq.(17) by lowering the radiation matter, even in the case for the fluid picture to be valid, the fact that the buoyant force by dipole scattering is weaker than in the fluid picture suggests that buoyant force alone is not enough for the GSL to be valid. Indeed, we can show that there exists cases satisfying both of the breakdown of the GSL and the validity of the approximation used .
### B the buoyant force by S-wave scattering
The non-Archimedean character of buoyant force shown in the previous subsection is attributed to the dipole dominant scattering. If we assume that a massless<sup>§</sup><sup>§</sup>§Since we concentrate on the case that the lowering process goes on far from the black hole in order to make the used approximations valid, massive fields less contribute to buoyant force. The reason is that massive quanta far from the black hole are much less “excited” for a stationary observer than massless ones. scalar field exists in nature, the argument based on dipole scattering does not work and implication on the GSL by wave nature of the acceleration radiation is drastically changed, due to S-wave scattering.
By mode decomposition of the equation of motion of the scalar field $`\text{ }\text{ }\text{ }\text{ }\text{ }\varphi =0`$ with respect to the plane wave in the local frame, we obtain
$`\left[k^2+\mathrm{\Delta }\right]\mathrm{\Psi }_\stackrel{}{k}=0,\varphi ={\displaystyle \frac{\mathrm{exp}(ik\tau )}{\sqrt{2k}}}{\displaystyle \frac{\mathrm{\Psi }_\stackrel{}{k}(\stackrel{}{x})}{(2\pi )^{3/2}}}.`$ (23)
Since we would like to consider the case that the total entropy of the box and contents remains constant, we regard the surface of the box as infinite potential barrier. Therefore, we solve a scattering problem by the infinite potential barrier at the radius $`b`$ in quantum mechanics. We easily obtain the result ,
$`\mathrm{\Psi }_\stackrel{}{k}(\stackrel{}{x})\mathrm{exp}(i\stackrel{}{k}\stackrel{}{x})+g(\mathrm{\Omega }){\displaystyle \frac{\mathrm{exp}(ikr)}{r}};{\displaystyle \frac{d\sigma }{d\mathrm{\Omega }}}=|g(\mathrm{\Omega })|^2,`$ (24)
$`g(\mathrm{\Omega })={\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{2l+1}{2ik}}\left(1+{\displaystyle \frac{h_l^{(2)}(kb)}{h_l^{(1)}(kb)}}\right)P_l(\mathrm{cos}\theta ),`$ (25)
where $`h_l^{(n)}`$ and $`P_l`$ are the spherical Hankel function of the $`n`$-th kind and the Legendre one of the first kind, respectively. In the long wavelength limit, the reflection coefficient $`g(\mathrm{\Omega })`$ is approximated by
$`g(\mathrm{\Omega })b{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(kb)^{2l}}{[(2l1)!!]^2}}P_l(\mathrm{cos}\theta ).`$ (26)
If S-wave scattering occurs, we have the differential cross section independent of wavelength,
$$\left(\frac{d\sigma }{d\mathrm{\Omega }}\right)_{l=0}b^2\frac{\sigma _T}{4\pi },\sigma _T4\pi b^2=4A,$$
(27)
where $`\sigma _T`$ is the total scattering cross section of the target.
In this connection, if the dipole dominated scattering occurs, such as the electromagnetic field, we have the differential cross section
$$\left(\frac{d\sigma }{d\mathrm{\Omega }}\right)_{l=1}b^2(kb)^4\left[P_1(\mathrm{cos}\theta )\right]^2,$$
(28)
which depends on the fourth power of $`k`$ as Eq.(18).
The contribution of the S-wave scattering to the momentum transfer of the target in the local frame is
$`{\displaystyle \frac{d\stackrel{}{P}}{d\tau }}={\displaystyle 𝑑k𝑑\stackrel{}{n}𝑑\stackrel{}{n}^{}I(k,\stackrel{}{n})\left(\frac{d\sigma }{d\mathrm{\Omega }^{}}\right)_{l=0}(\stackrel{}{n}\stackrel{}{n}^{})}={\displaystyle 𝑑\stackrel{}{n}\sigma _T𝑑k\left[\stackrel{}{n}I(k,\stackrel{}{n})\right]},`$ (29)
where the last equality is due to the spherical symmetric scattering of S-wave. Since the quantity $`\sigma _T𝑑k\left[\stackrel{}{n}I(k,\stackrel{}{n})\right]`$ is nothing but the momentum flux across the surface with the area $`\sigma _T`$ into the direction $`\stackrel{}{n}`$, Eq.(29) gives momentum transfer four times larger than the total momentum flux across the surface of the box.
Therefore, the buoyant force in the global frame exerted on the box by S-wave scattering is four times larger than that in the fluid picture, because of the diffractive effect,
$`F_{\mathrm{}}^{}{}_{}{}^{Scatt}=4F_{\mathrm{}}^{}{}_{}{}^{B}(\text{fluid})=4V\rho _{rad}(l){\displaystyle \frac{d\xi }{dl}}.`$ (30)
Since the diffractive effect of S-wave scattering strengthens the buoyant force than in the fluid picture, the above fact suggests that the GSL is protected by the buoyant force alone.
Indeed, following the argument in Sec.II for this case, we obtain the inequality from Eq.(15)
$`{\displaystyle \frac{\mathrm{\Delta }S_G}{V}}`$ $``$ $`3{\displaystyle \frac{\rho _{rad}}{T}}+s_{rad}(\rho _{rad})s_0`$ (31)
$``$ $`3{\displaystyle \frac{\rho _{rad}}{T}}+s_{rad}(\rho _{rad})s_{rad}(4\rho _{rad}),`$ (32)
where we explicitly denote the dependency of $`s_{rad}`$ on $`\rho _{rad}`$. The first line is given by $`\rho _0=4\rho _{rad}`$ at the floating point and the second comes from the assumption A1, $`s_0s_{rad}(\rho _0)=s_{rad}(4\rho _{rad})`$. Using the equations,
$`s_{rad}(\rho _{rad})={\displaystyle \frac{4}{3}}{\displaystyle \frac{\rho _{rad}}{T}},`$ (33)
$`{\displaystyle \frac{s_{rad}(\rho _{rad})}{s_{rad}(\rho _{rad}^{})}}=\left({\displaystyle \frac{\rho _{rad}}{\rho _{rad}^{}}}\right)^{3/4},`$ (34)
we can show the validity of the GSL
$$\frac{\mathrm{\Delta }S_G}{V}s_{rad}(\rho _{rad})\left(\frac{13}{4}2^{\frac{3}{2}}\right)>0.$$
(35)
Thus, if some massless scalar field exists, then without invoking a new physics such as an entropy bound for matter, the GSL holds thanks to the buoyant force strengthened by the diffractive effect of S-wave scattering of black hole atmosphere.
For the completeness, we should check the validity of the fluid picture in short wavelength limit. In this limit, we obtain the differential cross section
$`{\displaystyle \frac{d\sigma }{d\mathrm{\Omega }}}{\displaystyle \frac{A}{4\pi }},`$ (36)
and finally obtain the momentum transfer in the local frame as
$`{\displaystyle \frac{d\stackrel{}{P}}{d\tau }}`$ $`=`$ $`{\displaystyle 𝑑k𝑑\stackrel{}{n}𝑑\stackrel{}{n}^{}I(k,\stackrel{}{n})\left(\frac{d\sigma }{d\mathrm{\Omega }^{}}\right)(\stackrel{}{n}\stackrel{}{n}^{})}={\displaystyle 𝑑k𝑑\stackrel{}{n}\left[\stackrel{}{n}I(k,\stackrel{}{n})\right]𝑑\stackrel{}{n}^{}\frac{A}{4\pi }}`$ (37)
$`=`$ $`{\displaystyle 𝑑\stackrel{}{n}A𝑑k\left[\stackrel{}{n}I(k,\stackrel{}{n})\right]}.`$ (38)
As expected, the buoyant force in the geometrical optics approximation is equal to that in the fluid picture.
## IV Summary
In this letter, we briefly reviewed a gedankenexperiment to test the validity of the GSL, which is any process composed of adiabatically lowering the object toward the black hole and dropping into. In the analysis of this gedankenexperiment, the buoyant force by the black hole atmosphere plays an important role and the buoyant force is usually estimated by the pressure gradient of the radiation fluid. However, since the pressure exerted on the target is given by the change in the momentum flux, it is necessary to estimate the reflection coefficient on the surface of the target, in order to get the correct buoyant force. In the case that the size of the target $`b`$ is larger than a typical wavelength of the black hole atmosphere $`\lambda `$, the pressure exerted on the surface of the box is well estimated by the fluid picture for the black hole atmosphere. On the other hand, in the case that the lowering process goes on with satisfying $`b<\lambda `$, we cannot complete the reasoning which makes the GSL to hold by the buoyant force estimated in based on the fluid picture, because the fluid picture breaks down by diffractive effect of wave scattering.
For buoyant force far from the black hole, massless fields dominate over massive ones, due to less acceleration of the quasi-stationary target compared with their masses. Furthermore, in the long wavelength limit, the dependence of the scattering cross section on wavelength much varies according to the spin of the scattered wave. Therefore, it much depends on the spin of massless fields in nature whether we need to invoke a new physics such as an entropy bound for matter, in order to hold the GSL, or not. If some massless scalar field exists in nature, then the GSL can hold, due to the buoyant force alone by black hole atmosphere. If not so, the validity of the GSL might suggest the existence of some new physics such as an entropy bound.
The above conclusion is based on the viewpoint of an accelerated observer who rest on the box lowering adiabatically. Although the energy-momentum tensor normalized by the accelerated observer is different from the true one, we can expect that the calculation of buoyant force by the viewpoint of the accelerated observer gives correct estimation. It is because the essential quantity in the calculation is gradient of the energy-momentum tensor, not value itself, and the difference between the energy-momentum tensor normalized by the accelerated observer and the true one is divergence free.
In the Ref. , it was shown that in two dimensional spacetime, the estimation of $`\mathrm{\Delta }M(l)`$ delivered to the black hole in an accelerating viewpoint with the fluid approximation is equivalent to that in an inertial point of view. Does this equivalence suggest that the estimation of buoyant force by wave scattering is different from that in an inertial point of view, that is, not physical? Since, in two dimensional spacetime, the reflection coefficient of wave scattering by infinite potential is unity, two estimations of buoyant force in an accelerating viewpoint with and without the fluid approximation are equal to one another. Therefore, three estimations, including in an inertial point of view, are consistent and this result is due to the peculiarity of two dimensional spacetime. For completeness, it is worthwhile to estimate energetics $`\mathrm{\Delta }M(l)`$ from an inertial point of view for higher dimensional spacetimes.
Furthermore, although we regard the mere sum of black hole entropy and matter one as the total entropy, we have not yet obtained the foundation. Since gravity is long range force, it may be doubtful to assume the additivity of entropies of individual systems in self-gravitating system. It is future work to reconsider the GSL without the assumption of the additivity of entropies .
###### Acknowledgements.
The author would like to thank Professor T. Mishima and Dr. T. Shimomura for a careful reading of the manuscript. He also appreciate Professor M. Morikawa for kind hospitality in Ochanomizu University.
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# 1 Introduction
## 1 Introduction
In the late eighties, Turbiner and Ushveridze discovered some cases where a finite number of eigenvalues (and the associated eigenfunctions) of the spectral problem for the Schrödinger operator
$$\overline{H}\psi =E\psi ,\overline{H}=\frac{d^2}{dy^2}+V(y),yRorR^+$$
(1)
can be found explicitely. The corresponding problems have been called quasi-exactly solvable (QES). Since that first step, QES equations have been classified according to their relation with the finite-dimensional representations of the Lie algebra $`sl2,R)`$. Indeed QES Schrödinger equations as given in (1) can be written as
$$H\varphi =E\varphi ,H=p_4(x)\frac{d^2}{dx^2}+p_3(x)\frac{d}{dx}+p_2(x),xRorR^+$$
(2)
through ad-hoc changes of variables and functions
$$x=x(y),\psi =exp(\chi )\varphi ,$$
(3)
if $`p_j(x)(j=2,3,4)`$ refer to polynomials of order $`j`$ in $`x`$. The Hamiltonian $`H`$ in (2) can also be expressed in terms of the first-order differential operators
$$j_+=x^2\frac{d}{dx}+2jx,$$
$$j_0=x\frac{d}{dx}j,$$
(4)
$$j_{}=\frac{d}{dx},j=0,\frac{1}{2},1,\mathrm{},$$
satisfying the $`sl(2,R)`$ commutation relations i.e.
$$[j_0,j_\pm ]=\pm j_\pm ,$$
(5)
$$[j_+,j_{}]=2j_0,$$
(6)
the Casimir operator of this structure being $`C=j_+j_{}+j_0^2j_0`$. The operator $`H`$ is then
$$H=\underset{\mu ,\nu =\pm ,0,\mu \nu }{}c^{\mu \nu }j_\mu j_\nu +\underset{\mu =\pm ,0}{}c^\mu j_\mu $$
(7)
where the coefficients $`c^{\mu \nu },c^\mu `$ are arbitrary real numbers.
The crucial point in order to relate the operator (7) to QES problems is the introduction of the nonnegative integer $`2j`$ in (4). Indeed, the generators of $`sl(2,R)`$ as written in (4) preserve the space of polynomials of order $`2j`$
$$P(2j)=\{1,x,x^2,\mathrm{},x^{2j}\}$$
(8)
and so do the Hamiltonian (7). Searching for the eigenvalues of (1) is thus limited to the diagonalization of (7) in the $`(2j+1)`$-dimensional space $`P(2j)`$. It is a straightforward problem leading to the knowledge of the eigenvalues $`E_k(k=0,1,\mathrm{},2j)`$ as well as the corresponding eigenfunctions $`\psi _k`$ (cf. (3) where $`\varphi _k`$ belongs to $`P(2j)`$) and ensuring the quasi-exactly solvability of the original equation.
We propose in this paper to take a new look at this problem through the consideration of the so-called polynomial deformations of $`sl(2,R)`$ i.e. of the structures characterized by the following commutation relations
$$[J_0,J_\pm ]=\pm J_\pm ,$$
(9)
$$[J_+,J_{}]=p_n(J_0),$$
(10)
where $`p_n(J_0)`$ stands for a polynomial of order $`n`$ in the operator $`J_0`$.
More precisely, in Section 2, we show how to introduce in a natural manner the polynomial deformations (10) inside the operator (7) (n will then be limited to 3). In Section 3, we study the finite-dimensional representations of these polynomial deformations while Section 4 is devoted to their finite-dimensional differential realizations. In Section 5, we analyze two examples i.e. the sextic oscillator (Subsection 5.1) and the second harmonic generation (SHG) problem (Subsection 5.2). Finally, we give some conclusions in Section 6.
## 2 Polynomial deformations of $`sl(2,R)`$ inside QES problems
Instead of considering the operators (4) as expressed in (7), let us introduce the following operators (so defined for natural reasons in connection with their respective raising, diagonal or lowering characteristics)
$$J_+c^{++}j_+^2+c^{+0}j_+j_0+c^+j_+,$$
(11)
$$J_0c^+j_+j_{}+c^{00}j_0^2+c^0j_0,$$
(12)
$$J_{}c^0j_0j_{}+c^{}j_{}^2+c^{}j_{},$$
(13)
so that $`H`$ simply writes
$$H=J_++J_0+J_{}.$$
(14)
As we will see, restoring the linearity inside (7) such that it becomes the combination (14) will have the consequence of introducing nonlinearity inside (6) so that we will be concerned with the algebra (9)-(10). Indeed asking for the relations (9) to be satisfied with the operators (11)-(13) and the relations (5)-(6) lead to two cases only i.e. either
$$c^{++}=c^{}=0,c^+=c^{00},c^0+c^{00}=1$$
(15)
or
$$c^{++}0,c^{}0,c^{+0}=c^+=c^0=c^{}=0,c^+=c^{00},c^0+c^{00}=\frac{1}{2}.$$
(16)
The relation (10) is then ensured with the respective polynomials $`p_{n=3}(J_0)`$
$`p_3(J_0)=4c^{+0}c^0J_0^3+3(c^0c^++c^{+0}c^{}c^{+0}c^0)J_0^2`$
$`+[2c^+c^{}c^{+0}c^{}c^0c^++c^{+0}c^0(12j(j+1))]J_0`$ (17)
$`+j(j+1)(c^{+0}c^0c^{+0}c^{}c^0c^+)`$
or
$$p_3(J_0)=8c^{++}c^{}((2j^2+2j1)J_08J_0^3).$$
(18)
Notice that in these expressions, $`c^+(=c^{00})`$ has been put equal to zero without loosing generality (cf. the Casimir operator of $`sl(2,R)`$). Moreover, the relation (18) refers to the Higgs algebra already intensively visited .
We can thus consider the QES Hamiltonians as linear combinations of operators generating the following polynomial deformation of $`sl(2,R)`$
$$[J_0,J_\pm ]=\pm J_\pm ,$$
(19)
$$[J_+,J_{}]=\alpha J_0^3+\beta J_0^2+\gamma J_0+\delta ,\alpha ,\beta ,\gamma ,\delta R$$
(20)
taking account of the two possibilities (17) and (18). Notice that the Casimir operator of this deformed algebra is
$$C=J_+J_{}+\frac{\alpha }{4}J_0^4+(\frac{\beta }{3}\frac{\alpha }{2})J_0^2+(\frac{\alpha }{4}\frac{\beta }{2}+\frac{\gamma }{2})J_0^2+(\frac{\beta }{6}\frac{\gamma }{2}+\delta )J_0$$
.
The next step will be the determination of the finite-dimensional representations of the algebra (19)-(20) denoted in the following by $`sl^{(3)}(2,R)`$, the upper index referring to the highest power of the diagonal operator.
## 3 Finite-dimensional representations of $`sl^{(3)}(2,R)`$
As stated in the Introduction and in relation with the possible diagonalization of $`H`$, we are interested in the finite-dimensional ($`=2J+1,J=0,\frac{1}{2},1,\mathrm{}`$) representations of $`sl^{(3)}(2,R)`$, only. We thus consider kets of type $`J,M>`$ with $`M`$ running from $`J`$ to $`J`$ and such that
$$J_0J,M>=(\frac{M}{q}+c)J,M>,$$
(21)
$$J_+J,M>=f(M)J,M+q>,$$
(22)
$$J_{}J,M>=g(M)J,Mq>,$$
(23)
where $`q`$ is a positive integer and $`c`$ a real number. The relations (21)-(23) are such that (19) is satisfied. In order to ensure (20), we have to impose
$`f(Mq)g(M)f(M)g(M+q)=\alpha ({\displaystyle \frac{M}{q}}+c)^3+\beta ({\displaystyle \frac{M}{q}}+c)^2`$
$`+\gamma ({\displaystyle \frac{M}{q}}+c)+\delta ,M=J,\mathrm{},J.`$ (24)
Moreover, we have to take account of the dimension of the representations, leading to the constraints
$$f(J)=f(J1)=\mathrm{}=f(Jq+1)=0$$
(25)
and
$$g(J)=g(J+1)=\mathrm{}=g(J+q1)=0.$$
(26)
Being interested in the highest weight representations, we obtain from (24) and (25) the following result
$`f(J(k+1)ql)g(Jkql)=(k+1)\{\alpha ({\displaystyle \frac{Jl}{q}}+c)^3+\beta ({\displaystyle \frac{Jl}{q}}+c)^2`$
$`+\gamma ({\displaystyle \frac{Jl}{q}}+c)+\delta {\displaystyle \frac{1}{2}}[3\alpha ({\displaystyle \frac{Jl}{q}}+c)^2+2\beta ({\displaystyle \frac{Jl}{q}}+c)+\gamma ]k`$ (27)
$`+{\displaystyle \frac{1}{6}}[3\alpha ({\displaystyle \frac{Jl}{q}}+c)+\beta ]k(2k+1){\displaystyle \frac{\alpha }{4}}k^2(k+1)\}`$
where $`l=0,1,\mathrm{},q1`$ and $`k=0,1,\mathrm{},\frac{2Jdl}{q}`$. The nonnegative integer $`d`$ introduced in the last formula has to take specific values according to $`l`$ but also to $`J`$. These values are summarized in the following table, $`n`$ denoting a nonnegative integer.
| Table | $`l=0`$ | $`l=1`$ | $`l=2`$ | $`\mathrm{}`$ | $`l=q1`$ |
| --- | --- | --- | --- | --- | --- |
| $`J=(qn)/2`$ | $`d=0`$ | $`d=q1`$ | $`d=q2`$ | $`\mathrm{}`$ | $`d=1`$ |
| $`J=(qn+1)/2`$ | $`d=1`$ | $`d=0`$ | $`d=q1`$ | $`\mathrm{}`$ | $`d=2`$ |
| $`J=(qn+2)/2`$ | $`d=2`$ | $`d=1`$ | $`d=0`$ | $`\mathrm{}`$ | $`d=3`$ |
| $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ |
| $`J=(qn+q1)/2`$ | $`d=q1`$ | $`d=q2`$ | $`d=q3`$ | $`\mathrm{}`$ | $`d=0`$ |
Taking care of these values, we also have to constrain the real $`c`$ in (21) and (27) through the conditions (26). This leads to equal to zero the expression inside the brackets { . } in (27) for $`k=\frac{2Jdl}{q}`$ or, in other words, to consider the $`q`$ equations on $`c`$
$`\alpha [c^3+{\displaystyle \frac{3}{2q}}(dl)c^2+{\displaystyle \frac{1}{q^2}}(J^2J(d+l)+l^2dl+d^2)c+{\displaystyle \frac{1}{2q}}(2Jdl)c`$
$`+{\displaystyle \frac{1}{4q^3}}(2J^2(dl)2J(d^2l^2)+d^3d^2l+dl^2l^3)+{\displaystyle \frac{1}{4q^2}}(l^2d^2+2J(dl))]`$
$`+\beta [c^2+{\displaystyle \frac{1}{q}}(dl)c+{\displaystyle \frac{1}{3q^2}}(J^2J(d+l)+d^2dl+l^2)+{\displaystyle \frac{1}{6q}}(2Jdl)]`$
$`+\gamma (c+{\displaystyle \frac{1}{2q}}(dl))+\delta =0,l=0,1,\mathrm{},q1.`$ (28)
Excluding the trivial case $`(\alpha ,\beta ,\gamma ,\delta )=(0,0,0,0)`$, we notice that these equations reduce to
$$\alpha c(c^2+J(J+1))+\beta (c^2+\frac{1}{3}J(J+1))+\gamma c+\delta =0,$$
(29)
for $`q=1`$ while for $`q=2`$, we are led to either
$$\alpha =0\beta =0,c=\frac{\delta }{\gamma }$$
(30)
or
$$\alpha 0\delta =\frac{\beta \gamma }{3\alpha }\frac{2\beta ^3}{27\alpha ^2},c=\frac{\beta }{3\alpha }$$
(31)
if $`J=n`$ and to
$$\alpha (3c^2+\frac{1}{4}J(J+1)\frac{1}{8})+2\beta c+\gamma =0,$$
(32)
$$\alpha (c^3+\frac{1}{4}J(J+1)c)+\beta (c^2+\frac{1}{12}J(J+1))+\gamma c+\delta =0,$$
(33)
if $`J=n+\frac{1}{2}`$. Two other values of $`q`$ are also available namely $`q=3`$ and $`q=4`$. We respectively obtain
$$\alpha 0,\gamma =\frac{\beta ^2}{3\alpha }\frac{1}{9}\alpha J^2+\frac{2}{9}\alpha ,\delta =\frac{\beta \gamma }{3\alpha }\frac{2\beta ^3}{27\alpha ^2},c=\frac{\beta }{3\alpha }$$
(34)
if $`J=\frac{3n}{2}`$,
$$\alpha 0,\gamma =\frac{\beta ^2}{3\alpha }\frac{1}{9}\alpha J^2\frac{2}{9}\alpha J+\frac{1}{9}\alpha ,\delta =\frac{\beta \gamma }{3\alpha }\frac{2\beta ^3}{27\alpha ^2},c=\frac{\beta }{3\alpha }$$
(35)
if $`J=\frac{3n+1}{2}`$ and
$$\alpha 0,\gamma =\frac{\beta ^2}{3\alpha }\frac{1}{9}\alpha J^2\frac{1}{9}\alpha J+\frac{1}{9}\alpha ,\delta =\frac{\beta \gamma }{3\alpha }\frac{2\beta ^3}{27\alpha ^2},c=\frac{\beta }{3\alpha }$$
(36)
if $`J=\frac{3n+2}{2}`$, these three contexts being associated with $`q=3`$. We also have
$$\alpha 0,\gamma =\frac{\beta ^2}{3\alpha }\frac{1}{16}\alpha J^2+\frac{3}{16}\alpha ,\delta =\frac{\beta \gamma }{3\alpha }\frac{2\beta ^3}{27\alpha ^2},c=\frac{\beta }{3\alpha }$$
(37)
if $`J=2n`$ and
$$\alpha 0,\gamma =\frac{\beta ^2}{3\alpha }\frac{1}{16}\alpha J^2\frac{1}{8}\alpha J+\frac{1}{8}\alpha ,\delta =\frac{\beta \gamma }{3\alpha }\frac{2\beta ^3}{27\alpha ^2},c=\frac{\beta }{3\alpha }$$
(38)
if $`J=2n+1`$, these two cases being related with $`q=4`$. The two other systems related to $`q=4`$ i.e. those corresponding to $`J=2n+\frac{1}{2}`$ and $`J=2n+\frac{3}{2}`$ are incompatible ones as are those related to $`q>4`$.
Let us also notice that some of these representations are reducible. For instance, if we consider the case of the usual $`sl(2,R)`$ algebra, corresponding to $`\alpha =0,\beta =0,\gamma =2,\delta =0`$, we obtain
$$q=1c=0$$
(39)
and
$$q=2c=0,J=n.$$
(40)
The first case (39) is associated to (see (27))
$$f(Jk1)g(Jk)=(k+1)\{2Jk\},k=0,1,\mathrm{},2J$$
(41)
or, in other words, to
$$f(M1)g(M)=(JM+1)(J+M),M=J,\mathrm{},J.$$
(42)
We recognize in (42) the well known result of the angular momentum theory subtended by this $`sl(2,R)`$ algebra. The second case (40) corresponds to
$$f(J2k2)g(J2k)=(k+1)\{Jk\},k=0,1,\mathrm{},J,$$
(43)
$$f(J2k3)g(J2k1)=(k+1)\{Jk1\},k=0,1,\mathrm{},J1,$$
(44)
or, in other words, to
$$f(M2)g(M)=\frac{1}{4}(JM+2)(J+M),M=J,J+2,\mathrm{},J,$$
(45)
$$f(M2)g(M)=\frac{1}{4}(JM+1)(J+M1),M=J+1,J+3,\mathrm{},J1.$$
(46)
It is then clear that the representation $`(J=n,q=2)`$ is in fact the direct sum of the two (irreducible) representations $`(J=\frac{n1}{2},q=1)`$ and $`(J=\frac{n}{2},q=1)`$ (the eigenvalues of the Casimir being equal).
## 4 Finite-dimensional differential realizations of $`sl^{(3)}(2,R)`$
We now turn to the construction of the differential realizations (expressed in terms of the real variable $`x`$) of the algebra (19)-(20). In correspondence with (21)-(23), the generators of $`sl^{(3)}(2,R)`$ have the following forms
$$J_+x^qF(D),$$
(47)
$$J_0\frac{1}{q}(DJ)+c,$$
(48)
$$J_{}G(D)\frac{d^q}{dx^q}$$
(49)
and the basis $`\{J,M>,M=J,\mathrm{},J\}`$ stands for the space $`P(2J)`$ (cf. (8)) of monomials $`\{x^{J+M},M=J,\mathrm{},J\}`$. Moreover, $`D`$ is the dilation operator
$$Dx\frac{d}{dx}.$$
(50)
By remembering that
$$\frac{d^q}{dx^q}x^q=\underset{k=1}{\overset{q}{}}(D+k)\frac{(D+q)!}{D!}$$
(51)
and
$$x^q\frac{d^q}{dx^q}=\underset{k=0}{\overset{q1}{}}(Dk)\frac{D!}{(Dq)!},$$
(52)
the relation (20) gives
$`F(Dq)G(Dq){\displaystyle \frac{D!}{(Dq)!}}F(D)G(D){\displaystyle \frac{(D+q)!}{D!}}=`$
$`\alpha [{\displaystyle \frac{1}{q}}(DJ)+c]^3+\beta [{\displaystyle \frac{1}{q}}(DJ)+c]^2`$
$`+\gamma [{\displaystyle \frac{1}{q}}(DJ)+c]+\delta .`$ (53)
Let us discuss this condition within the $`(\alpha 0)`$-case first and the $`(\alpha =0)`$-case second. a) In order to avoid singularities, we thus impose, when $`\alpha 0`$,
$$F(D)G(D)=\frac{\alpha }{4q^4}\frac{D!}{(D+q)!}(D+\lambda _1)(D+\lambda _2)(D+\lambda _3)(D+\lambda _4)$$
(54)
and obtain the following system on the real unknows $`\lambda _1,\mathrm{},\lambda _4`$
$$\lambda _1+\lambda _2+\lambda _3+\lambda _4=4qc+2q4J+\frac{4}{3}\frac{\beta }{\alpha }q,$$
(55)
$`\lambda _1\lambda _2+\lambda _1\lambda _3+\lambda _1\lambda _4+\lambda _2\lambda _3+\lambda _2\lambda _4+\lambda _3\lambda _4=6q^2c^2+6q^2c12qJc`$
$`+4{\displaystyle \frac{\beta }{\alpha }}q^2c+6J^26qJ4{\displaystyle \frac{\beta }{\alpha }}qJ+q^2+2{\displaystyle \frac{\beta }{\alpha }}q^2+2{\displaystyle \frac{\gamma }{\alpha }}q^2,`$ (56)
$`\lambda _1\lambda _2\lambda _3+\lambda _1\lambda _2\lambda _4+\lambda _1\lambda _3\lambda _4+\lambda _2\lambda _3\lambda _4=`$
$`4q^3c^312q^2Jc^2+6q^3c^2+4{\displaystyle \frac{\beta }{\alpha }}q^3c^2+12qJ^2c12q^2Jc`$
$`+2q^3c8{\displaystyle \frac{\beta }{\alpha }}q^2Jc+4{\displaystyle \frac{\beta }{\alpha }}q^3c+4{\displaystyle \frac{\gamma }{\alpha }}q^3c4J^3+6qJ^22q^2J`$
$`+4{\displaystyle \frac{\beta }{\alpha }}qJ^24{\displaystyle \frac{\beta }{\alpha }}q^2J+{\displaystyle \frac{2}{3}}{\displaystyle \frac{\beta }{\alpha }}q^34{\displaystyle \frac{\gamma }{\alpha }}q^2J+2{\displaystyle \frac{\gamma }{\alpha }}q^3+4{\displaystyle \frac{\delta }{\alpha }}q^3.`$ (57)
Let us recall that we are interested in finite-dimensional (=$`2J+1`$) realizations only. This means that the cases $`q=1`$ and $`q=2`$ are the only possibilities in accordance with
$$q=1\lambda _1=1,\lambda _2=2J$$
(58)
and
$$q=2\lambda _1=1,\lambda _2=2,\lambda _3=2J,\lambda _4=2J+1.$$
(59)
In the first context ($`q=1`$), the equations (55) and (56) fix $`\lambda _3`$ and $`\lambda _4`$ as follows
$$\lambda _3=2c+\frac{1}{2}J+\frac{2}{3}\frac{\beta }{\alpha }+\frac{ϵ}{2},\lambda _4=2c+\frac{1}{2}J+\frac{2}{3}\frac{\beta }{\alpha }\frac{ϵ}{2}$$
(60)
with
$$ϵ^2=1+\frac{16}{9}\frac{\beta ^2}{\alpha ^2}4J(J+1)8c^2\frac{16}{3}\frac{\beta }{\alpha }c8\frac{\gamma }{\alpha }$$
(61)
while the equation (57) coincides with (29). In the second context ($`q=2`$), we are led to (31) supplemented by
$$\gamma =\frac{\beta ^2}{3\alpha }\frac{\alpha }{4}J^2\frac{\alpha }{4}J+\frac{\alpha }{8}.$$
(62)
b) The case $`\alpha =0`$ is more simple and has already been analyzed . For self-consistency, we recall the main results i.e.
$$F(D)G(D)=\frac{\beta }{3q^3}\frac{D!}{(D+q)!}(D+\lambda _1)(D+\lambda _2)(D+\lambda _3)$$
(63)
where the only possible finite-dimensional (=$`2J+1`$) realization is associated with $`q=1`$ and corresponds to
$$\lambda _1=1,\lambda _2=2J,\lambda _3=J+3c+\frac{3\gamma }{2\beta }+\frac{1}{2},$$
(64)
the real $`c`$ being fixed through
$$\beta c^2+\gamma c+\delta +\frac{\beta }{3}J(J+1)=0.$$
(65)
Now that both cases have been considered, let us conclude this Section by noticing that some realizations (namely the ones corresponding to $`q=2`$ without the condition (62) and the ones associated with $`q=3,4`$) are missing with respect to the representations developed in the previous Section. Indeed, the relation (53) we have imposed is more constraining because it is a relation between operators independently of the basis on which they are supposed to act. If we take account of this basis i.e. in this case $`P(2J)=\{x^{J+M},M=J,\mathrm{},J\}`$, we can recover all the cases previously discussed. For example, in the context $`q=3`$, $`J=\frac{3}{2}`$, we can consider
$$J_+=\frac{1}{6}f(\frac{3}{2})x^3(D1)(D2)(D3),$$
(66)
$$J_0=\frac{1}{3}D\frac{1}{2}\frac{\beta }{3\alpha },$$
(67)
$$J_{}=\frac{1}{6}g(\frac{3}{2})\frac{d^3}{dx^3}$$
(68)
with
$$f(\frac{3}{2})g(\frac{3}{2})=\frac{\alpha }{9}.$$
(69)
It is then easy to convince ourselves that these operators generate $`sl^{(3)}(2,R)`$ with $`\gamma =\frac{\beta ^2}{3\alpha }\frac{\alpha }{36}`$ and $`\delta =\frac{\beta ^3}{27\alpha ^2}\frac{\beta }{108}`$ but on the space $`P(3)`$ only (the relation (53) being trivially not satisfied except on this space). However it has to be stressed that the relations corresponding to (66)-(68) but in the general context become really heavy when the value of $`J`$ increases.
## 5 Two examples
We first consider the prototype of QES systems i.e. the so-called sextic oscillator and then turn to a more physical example: the SHG problem.
### 5.1 The sextic oscillator
This system is characterized by the following potential
$$V(y)=a^2y^6+2aby^4+(b^22ap8aj3a)y^2$$
(70)
with $`a(0),bR`$ and $`p=0,1`$ while $`j`$ is the quantum number appearing in (4). With
$$x=y^2$$
(71)
and
$$\chi =(\frac{a}{2}x+\frac{b}{2}\frac{p}{2x})𝑑x,$$
(72)
we can be convinced of its QES character via the form (7) (up to a translation)
$$H=J_++J_0+J_{}+2bp+4bj+b$$
(73)
and
$$c^0=4,c^+=4a,c^0=4b,c^{}=(4j+2+4p),$$
(74)
the other $`c^{}s`$ being equal to zero. Without loss of generality, we can put $`b=\frac{1}{4}`$ (in order to recover (19)) and obtain, through (11)-(13), the relation (20) with
$$\alpha =0,\beta =48a,\gamma =32a(p+j),\delta =16aj(j+1).$$
(75)
The case $`q=1`$ is thus the only one to be available. We actually have
$$J_0J,M>=(M+c)J,M>,$$
(76)
$$J_+J,M>=f(M)J,M+1>,$$
(77)
$$J_{}J,M>=g(M)J,M1>,$$
(78)
with
$$f(M1)g(M)=(JM+1)(J+M)(48ac+16aM+16ap+16aj8a).$$
(79)
Moreover, the parameter $`c`$ is fixed according to
$$c=\frac{1}{3}(p+j)\pm \frac{1}{3}\sqrt{(p+j)^23J(J+1)+3j(j+1)}$$
(80)
leading to constrain $`J`$ through
$$J\frac{1}{2}+\frac{1}{6}\sqrt{36j(j+1)+12(p+j)^2+9}$$
(81)
in order to ensure the reality of $`c`$. Because the space $`P(2J)`$ is preserved, we just have to equal to zero the determinant of the following matrix
$`\mathrm{M}={\displaystyle \underset{\mathrm{k}=1}{\overset{2\mathrm{J}+1}{}}}(\mathrm{E}+\mathrm{J}\mathrm{k}\mathrm{c}{\displaystyle \frac{\mathrm{p}}{2}}\mathrm{j}+{\displaystyle \frac{3}{4}})\mathrm{e}_{\mathrm{k},\mathrm{k}}`$
$`{\displaystyle \underset{k=1}{\overset{2J}{}}}g(J+k)e_{k,k+1}{\displaystyle \underset{k=1}{\overset{2J}{}}}f(J+k1)e_{k+1,k}`$ (82)
in order to find the energies. In the matrix (82), the notation $`e_{k,l}`$ stands for a $`(2J+1)`$-dimensional matrix where 1 is at the intersection of the $`k^{th}`$ row and the $`l^{th}`$ column, the other elements being 0. For example, when $`j=\frac{1}{2}`$, we have
$$J=0,\frac{1}{2},$$
(83)
according to (81) while the relation (80) gives
$$c=\frac{p}{3}\frac{1}{6}\pm \frac{1}{3}\sqrt{(p+\frac{1}{2})^2+\frac{9}{4}}$$
(84)
if $`J=0`$ and
$$c=0orc=\frac{1}{3}(2p+1)$$
(85)
if $`J=\frac{1}{2}`$. In the case (84), the energies are
$$E=c+\frac{p}{2}+\frac{3}{4}E=0.0428932;0.0562872;1.1103796;1.4571067$$
(86)
and in the case (85), the resolution of the vanishing determinant associated with (82) leads to
$$E=\frac{3}{4}\pm \frac{1}{2}\sqrt{1+32a};E=\frac{5}{4}\pm \frac{1}{2}\sqrt{1+96a}$$
(87)
if $`c=0`$ and
$$E=\frac{5}{12}\pm \frac{1}{2}\sqrt{132a};E=\frac{1}{4}\pm \frac{1}{2}\sqrt{196a}$$
(88)
if $`c=\frac{1}{3}`$; $`c=1`$. Only the values given in (87) correspond to the previously obtained ones . This is indeed a general result that the known energies are recovered through our approach when $`c=0,J=j`$. The other contexts ($`J<j`$) or ($`J=j,c=\frac{2}{3}(p+j)`$) lead to supplementary new values of the energy.
Let us analyze more deeply this result by going to the differential realization (47)-(49) i.e.
$$J_+=xF(D),$$
(89)
$$J_0=DJ+c,$$
(90)
$$J_{}=G(D)\frac{d}{dx},$$
(91)
where (cf. (63) and (64))
$$F(D)G(D)=16a(D2J)(DJ+3c+\frac{1}{2}+p+j).$$
(92)
In order to preserve the space $`P(2J)`$, let us make the choice (without loss of generality, this freedom being due to the fact that $`sl^{(3)}(2,R)`$ is defined up to an automorphism )
$$G(D)=4(DJ+3c+\frac{1}{2}+p+j).$$
(93)
In that case, the Hamiltonian (73) is realized as
$`H=4x{\displaystyle \frac{d^2}{dx^2}}+[4ax^2+x4(J+3c+{\displaystyle \frac{1}{2}}+p+j)]{\displaystyle \frac{d}{dx}}`$
$`8aJx+{\displaystyle \frac{p}{2}}+{\displaystyle \frac{1}{4}}+jJ+c.`$ (94)
This form is analog to (2) and we propose to write it in the Schrödinger form (1) through the changes (3) i.e.
$$x=y^2,$$
(95)
$$\psi =exp((\frac{1}{8}2ax+\frac{1}{2}(J3cpj)\frac{1}{x})𝑑x)\varphi .$$
(96)
The potential obtained in this manner is given by
$`V(y)=a^2y^6+{\displaystyle \frac{1}{2}}ay^4+({\displaystyle \frac{1}{16}}6aJ2aj6ac2ap3a)y^2`$
$`+{\displaystyle \frac{1}{2}}(jJc)+(J3cpj)(J3cpj+1){\displaystyle \frac{1}{y^2}}.`$ (97)
Compared with (70), this expression actually reduces to the sextic oscillator potential iff $`c=0`$ and $`J=j`$. For other values of the parameters (i.e. the already cited $`(J<j)`$ and $`(J=j,c=\frac{2}{3}(p+j))`$, the new eigenvalues of the problem (appearing for $`j=\frac{1}{2}`$ in (86) and (88)) do correspond to another model, namely the radial sextic oscillator as shown by (97).
### 5.2 The second harmonic generation
This nonlinear optical process as well as others such as coherent spontaneous emission and down conversion can be described by the following effective Hamiltonian
$$H=a_1^{}a_1+2a_2^{}a_2+g(a_2^{}a_1^2+(a_1^{})^2a_2)$$
(98)
with cubic terms in the (independent) bosonic creation and annihilation operators. The Hamiltonian (98) has already been recognized as a QES model. We confirm such a result by using the technique developed in Section 2. Indeed following Section 2, we propose to define
$$J_+=a_2^{}a_1^2,$$
(99)
$$J_0=\frac{1}{3}(a_2^{}a_2a_1^{}a_1),$$
(100)
$$J_{}=(a_1^{})^2a_2$$
(101)
such that the algebra (19)-(20) is satisfied with
$$\alpha =0,\beta =12,\gamma =0,\delta =\frac{1}{3}N^2+N.$$
(102)
In the last expression, $`N`$ is the invariant
$$N=a_1^{}a_1+2a_2^{}a_2$$
(103)
satisfying
$$[N,J_0]=[N,J_\pm ]=0.$$
(104)
Once again, the values (102) are typical of the $`q=1`$-representation only so that the relations (76)-(78) are the ones to be taken care of but with
$$f(M1)g(M)=(JM+1)(J+M)(12c4M+2).$$
(105)
The parameter $`c`$ is fixed according to
$$c^2=\frac{1}{3}J(J+1)+\frac{1}{36}N^2+\frac{1}{12}N$$
(106)
and its reality is ensured if
$$J\frac{1}{2}+\frac{1}{2}\sqrt{\frac{1}{3}N^2+N+1}.$$
(107)
It is then possible to determine the energies by putting to zero the determinant of a $`(2J+1)`$ by $`(2J+1)`$ matrix analog to (82) as well as it is possible to determine them through the differential realization (89)-(91). In this case, we have
$$F(D)G(D)=4(D2J)(DJ+3c+\frac{1}{2})$$
(108)
and choosing
$$G(D)=1$$
(109)
the Hamiltonian (98) becomes
$`H=N+4gx^3{\displaystyle \frac{d^2}{dx^2}}+g(1+12(J+c+{\displaystyle \frac{1}{2}})x^2){\displaystyle \frac{d}{dx}}`$
$`+4g(J+2J^26Jc)x.`$ (110)
With the respective changes of variables and wavefunctions
$$x=\frac{1}{gy^2},\psi =exp((\frac{1}{8x^3}+\frac{3}{2}(cJ)\frac{1}{x})𝑑x)\varphi $$
(111)
we can put (110) on the Schrödingerlike form (1) with
$$V(y)=\frac{g^4}{16}y^6+\frac{3}{2}g^2(cJ\frac{1}{2})y^2+(J+3c)(J+3c+1)\frac{1}{y^2}+N.$$
(112)
Once again this is typical of a radial sextic oscillator and the QES characteristics of the second harmonic generation is thus proved. The determination of the energies is then possible without any difficulty . For example, when $`N=4`$, we have
$$J=0E=4,$$
(113)
$$J=\frac{1}{2}E=4\pm g\sqrt{2\sqrt{19}},$$
(114)
$$J=1E=4,4\pm 4g.$$
(115)
Only the values (115) correspond to known ones, the values (113), (114) coming from other models. This is a general result in the sense that the energies of SHG are the ones of the Schrödinger potential (112) with
$$c=J\frac{N}{3}$$
(116)
and
$$J=\frac{N}{4}orJ=\frac{N1}{4}$$
(117)
according to even or odd values of $`N`$. The SHG potential thus writes
$$V(y)=\frac{g^4}{16}y^6\frac{g^2}{4}(2N+3)y^2+N$$
(118)
and exactly coincides with the potential (42) of Ref. . In terms of the operators (4) (with $`j=J`$), this gives
$$H=gj_+(j_0+\frac{1}{2}\frac{N}{4})4gj_{}+N$$
(119)
if $`N`$ is even and
$$H=gj_+(j_0\frac{1}{4}\frac{N}{4})4gj_{}+N$$
(120)
if $`N`$ is odd.
## 6 Conclusions
We have developed a general method based on the polynomial deformations of the Lie algebra $`sl(2,R)`$ in order to exhibit the QES characteristics of a Hamiltonian. We have applied this method to two examples: one more theoretical -the sextic oscillator- and one more physical -the second harmonic generation-. In both cases, a finite number of energies as well as eigenfunctions are determined through the finite-dimensional representations -or, in an equivalent way, through the realizations- of these polynomial deformations. Some of these energies (and eigenfunctions) were previously known, not the others. It seems, through the analysis of the two examples, that these previously unknown energies do not correspond to the same model but to another one being closely related to the first one. In some cases, the comparison between these new and old models could be interesting, with respect to experimental data in particular.
The main advantage of the method we have proposed is that it can be systematically applied to any Hamiltonian written in terms of a raising and a lowering operator. Numerous physical Hamiltonians are of that type. In particular, we plan to analyze one of them: the so-called Lipkin-Meshkov-Glick Hamiltonian of specific interest in nuclear physics. Being based on a $`(\alpha 0)`$ polynomial deformation, its analysis is more delicate but also richer in the number of available representations .
The main drawback is that it is limited to $`sl(2,R)`$ and its deformations when we know that some QES models need more extended algebras. However, the method we have presented here can also be generalized to these extended algebras as well as superalgebras. We also plan to come back on these points in the near future.
ACKNOWLEDGMENTS
I would like to warmly thank Prof. J. Beckers (University of Liège, Belgium) for fruitful discussions. Thanks are also due to Prof. A. Klimov (University of Guadalajara, Mexico) for useful information on the SHG model and for pointing out the Ref. in particular.
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# 1 Results of the simulation[] of “𝑍 bursts” homogeneously distributed up to a redshift of 4. The points and error bars are the AGASA data. The point labelled with 𝜈_{𝑈𝐻𝐸} shows the needed flux of ultrahigh energy neutrinos at the Z-resonance energy, with no relic lepton asymmetry (it could be reduced by at most 1/20 with large lepton asymmetries).
UCLA/00/TEP/16
May 2000
Super-Kamiokande 0.07 eV Neutrinos in Cosmology: Hot Dark Matter and the Highest Energy Cosmic Rays<sup>*</sup><sup>*</sup>*Talk given at the “4th International Symposium on Sources and Detection of Dark Matter in the Universe”, February 23-25, 2000, Marina del Rey, CA (to appear in its proceedings) and at the “Cosmic Genesis and Fundamental Physics” workshop, October 28-30, 1999, Sonoma State University, Santa Rosa, CA.
Graciela Gelmini<sup>1</sup>
<sup>1</sup>Dept. of Physics and Astronomy, UCLA, Los Angeles, CA 90095-1547
## Abstract
Relic neutrinos with mass in the range indicated by Super-Kamiokande results if neutrino masses are hierarchial (about 0.07 eV) are many times deemed too light to be cosmologically relevant. Here we remark that these neutrinos may significantly contribute to the dark matter of the Universe (with a large lepton asymmetry $`L`$) and that their existence might be revealed by the spectrum of ultra high energy cosmic rays (maybe even in the absence of a large $`L`$).
Super-Kamiokande has provided a strong evidence for the oscillation in atmospheric showers of two neutrino species with masses $`m_1`$, $`m_2`$ and $`\delta m^2`$ = $`m_1^2m_2^2`$ = $`(18)\times 10^3`$ eV consisting mostly of about equal amounts of $`\nu _\mu `$ and another flavor eigenstate neutrino, $`\nu _\tau `$ or a sterile neutrino. If neutrino masses are hierarchial, as those of the other leptons and quarks, then the heavier of the two oscillating neutrinos, call it $`\nu _{\mathrm{SK}}`$, has a mass $`m_{\mathrm{SK}}=\sqrt{\delta m^2}0.07`$ eV.
The possibility of $`m_1`$ and $`m_2`$ being much larger than $`m_{\mathrm{SK}}`$ has also been invoked, in part with the motivation to have $`m_1m_2`$ of the order of eV, in the range previously considered necessary for relic neutrinos to constitute a cosmologically relevant component of the dark matter in the Universe. In fact, with no lepton asymmetry, i.e. with $`L_\nu [(n_\nu n_{\overline{\nu }})/n_\gamma ]=0`$, the number density of relic neutrinos and antineutrinos of each species is $`n_\nu =n_{\overline{\nu }}=3n_\gamma /22=56cm^3`$. With this number density, the contribution of relic neutrinos to the energy density of the Universe (in units of the critical density) is $`\mathrm{\Omega }_\nu h^2=(_im_{\nu _i}/92`$ eV) (here $`h0.7`$ is the Hubble constant in units of 100 km/Mpc sec). This amounts to only to $`\mathrm{\Omega }_\nu h^2=0.8\times 10^3`$, while a value about 10 times larger was considered necessary in the context of Cold-Hot Dark Matter (CHDM) models .
In Ref. A. Kusenko and I pointed out that if the lepton asymmetry of $`\nu _{\mathrm{SK}}`$ in the Universe is of order one the neutrinos with $`m_{\mathrm{SK}}`$ can make a significant contribution to the energy content of the Universe This was also pointed out in Ref.. However, in this ref. the neutrino decoupling temperature $`T_d`$, which increases with increasing values of $`|L_\nu |`$ due to the effect of Fermi blocking factors, was instead taken to decrease with $`|L_\nu |`$, which lead to an incorrect relation between lepton number and density.
$$\mathrm{\Omega }_\nu h^20.01\left(\frac{|L_\nu |}{3.6}\right)\left(\frac{m_\nu }{0.07eV}\right),$$
(1)
In the same reference A. Kusenko and I also pointed out that, not only Fermi-degenerate relic neutrinos with $`m_{_{SK}}`$ could be revealed in the highest energy cosmic rays (if these are due to “Z-bursts”) but, based on the results of Adams and Sarkar, these neutrinos could be a new form of Hot Dark Matter. Adams and Sarkar had found that a massless relic neutrino species with chemical potential $`\mu _\nu =3.4T_\nu `$ added to a ‘standard’ Cold Dark Matter (CDM) model (flat matter dominated Universe with no cosmological constant) provides a good fit to the large scale structure (LSS) and anisotropy of the cosmic microwave background radiation (CMBR) data.
Neutrinos with mass $`m_{\mathrm{SK}}`$, contrary to the neutrinos studied by Adams and Sarkar, are non-relativistic at present, however they are still relativistic at the time of radiation-matter equality (the time at which the Universe becomes matter dominated), when their effect is most important.
In fact, the results of Adams and Sarkar were later confirmed and expanded by Lesgourgues and Pastor who studied the impact of both massless and massive relic neutrinos with $`m_{\mathrm{SK}}`$, with large lepton asymmetries on structure formation and the CMBR anisotropy. They studied CDM models with and without cosmological constant $`\mathrm{\Lambda }`$. I propose to call these models LCDM and L$`𝚲`$CDM, where L stands for the addition of a large lepton asymmetry. The major effect of the lighter and more abundant relic neutrinos is to delay the onset of matter domination in the Universe (which increases the amplitudes of the acoustic peaks in the angular spectrum of CMBR anisotropies). The neutrino mass is thus largely irrelevant, while neutrinos are relativistic at the moment of radiation-matter equality.
For large $`L_\nu `$ the relation between this asymmetry and the chemical potential $`\xi \mu _\nu /T_\nu `$ (which is constant after neutrinos decouple) is
$$L_\nu =\frac{1}{12\zeta (3)}\left(\frac{T_\nu }{T_\gamma }\right)^3[\pi ^2\xi +\xi ^3]=0.0252(9.87\xi +\xi ^3).$$
(2)
Here, $`\zeta (3)=1.202`$ and $`(T_\nu /T_\gamma )^3=4/11`$. This value for the temperature ratio is valid as long as the neutrino decoupling temperature $`T_d`$ is lower than the muon mass, which translates into the upper bound $`\xi <12`$ .
Neutrinos with chemical potentials $`\xi 1`$ ($`L_\nu 0.27)`$ are Fermi-degenerate. Only for these values of $`\xi `$ the number density of relic neutrinos becomes considerably different than in the usual case with no asymmetry.
After neutrino-antineutrino annihilation ceases in the early Universe, only the particles in excess remain and $`|L_\nu |`$ is just the ratio of the number density of these particles, say $`n_\nu `$, over the photon density, $`|L_\nu |=n_\nu /n_\gamma `$. For $`\xi =5`$, for example, one obtains $`L_\nu =4`$, which means that there are 4 background neutrinos for every background photon (thus, neutrinos dominate the entropy of the Universe) and consequently, $`n_\nu 1700\mathrm{cm}^3`$. This would make relic neutrinos 30 times more abundant than standard neutrinos or antineutrinos of every species with no lepton asymmetry and would make the relic density $`\mathrm{\Omega }_\nu h^2(m_\nu /3`$ eV).
Lesgourgues and Pastor found that, even with no cosmological constant (an assumption disfavored at present by type IA Supernovae data), i.e. in a LCDM model, $`\nu _{\mathrm{SK}}`$ with chemical potentials $`\xi `$ between 3 and 6 ($`L_\nu `$ between 1.4 and 6.9) added to CDM provide a good agreement with LSS and CMBR observations. In L$`𝚲`$CDM models instead, with a cosmological constant contribution of $`\mathrm{\Omega }_\mathrm{\Lambda }=0.5`$ to the energy density, $`\xi `$ could be between 0 to 4 ($`L_\nu `$ between 0 and 2.6). As the cosmological constant increases there is less room for neutrinos and for $`\mathrm{\Omega }_\mathrm{\Lambda }>0.7`$ no-lepton asymmetry is allowed, i.e. $`\xi =0`$. Kinney and Riotto also studied the CMB anisotropy in the presence of large lepton asymmetries, with similar results.
The large lepton asymmetries invoked here may seem odd. However, they have been studied several times from 1967 onwards, in the context of nucleosynthesis. In the presence of neutrino degeneracy nucleosynthesis becomes severely non-standard. A large number of electron neutrinos, $`\nu _e`$, present during nucleosynthesis yields a reduction of the neutron to proton ratio, $`n/p`$, through the reaction $`n\nu _epe`$. This in turn lowers the <sup>4</sup>He abundance, since when nucleosynthesis takes place essentially all neutrons end up in <sup>4</sup>He nuclei. Extra neutrinos of any flavor increase the energy density of the Universe, leading to an earlier decoupling of weak interactions and consequent increase of the $`n/p`$ ratio (and <sup>4</sup>He). This last effect is less important than the former one in the case of $`\nu _e`$, but it is the only effect of an excess of $`\nu _\mu `$ and/or $`\nu _\tau `$ (or their antinuetrinos). Thus, when both the chemical potentials of $`\nu _e`$ and of $`\nu _\mu `$ or $`\nu _\tau `$ are large, their effects largely compensate each other.
Combining nucleosynthesis bounds with the requirement of neutrinos becoming subdominant before the recombination epoch, in 1992 Kang and Steigman found
$$0.06\xi _{\nu _e}1.1,|\xi _{\nu _\mu ,\nu _\tau }|6.9,\mathrm{\Omega }_Bh^20.069.$$
(3)
The last is a bound on the baryon density $`\mathrm{\Omega }_B`$ almost an order of magnitude larger than obtained in conventional nucleosynthesis. Newer data on primordial element abundances seem to impose similar bounds on neutrino chemical potentials . Notice that it is the upper bound on $`\xi _{\nu _\mu }`$ which is relevant for $`\nu _{SK}`$.
The existence of ultra high energy cosmic rays (UHECR) with energies above the Greisen-Zatsepin-Kuzmin (GZK) cutoff of about $`5\times 10^{19}`$ eV, presents a problem. Photons and nucleons with those energies have short attenuation lengths and could only come from distances of 100 Mpc or less, while possible astrophysical sources for those energetic particles are much farther away. An elegant and economical solution to this problem, proposed by T. Weiler consists of the production of the necessary photons and nucleons close to Earth, in the annihilation at the $`Z`$-resonance of ultra-high energy neutrinos, $`\nu _{\mathrm{UHE}}`$, coming from remote sources, and relic background neutrinos. These events were named “$`Z`$-bursts” by T. Weiler. The $`Z`$-resonance occurs when the energy of the incoming $`\nu _{\mathrm{UHE}}`$ neutrinos is $`E_{\nu _{\mathrm{UHE}}}=E_{Res}`$,
$$E_{Res}=\frac{M_Z^2}{2m_\nu },$$
(4)
where $`m_\nu `$ is the mass of the relic neutrinos. Since galaxy formation arguments show $`m_\nu <`$ few eV, then $`E_{\nu _{\mathrm{UHE}}}>10^{21}`$ eV, precisely above the GZK cutoff, as needed. As we see in the equation, in this mechanism the energy cutoff $`E_{Res}`$ of the UHECR is related to the mass of the relic neutrinos, and this cutoff should be $`E_{Res}0.610^{23}`$ eV for $`m_\nu =m_{\mathrm{SK}}`$.
Depending on the assumed spectrum of $`\nu _{\mathrm{UHE}}`$ used, upper bounds on the intensity of the $`\nu _{\mathrm{UHE}}`$ flux can be obtained. These bounds determine the need for an enhancement in the density of relic neutrinos above the standard density of 56 cm<sup>-3</sup>, to account for the observed flux of UHECR. In the case of eV mass neutrinos the enhancement could come from gravitational clustering. Neutrinos with $`m_{\mathrm{SK}}`$ are too light to cluster significantly, but, if needed, the density enhancement could come from a large lepton asymmetry.
Most bounds on “$`Z`$-burst” models (see for example ) assume that the $`\nu _{\mathrm{UHE}}`$ have a typical astrophysical spectrum, decreasing with energy as $`E^\gamma `$, with $`\gamma `$ a number of order one. These bounds would not hold if the $`\nu _{\mathrm{UHE}}`$ spectrum had a very different energy dependence, as, for example if the sources would be unstable superheavy relic particles, which form part of the cold dark matter, decaying mostly into neutrinos, . In this case the spectrum of $`\nu _{\mathrm{UHE}}`$ is opposite to an astrophysical spectrum, it grows rapidly with energy, up to a sharp cutoff at an energy of the order of the parent particle mass. A model for these parent particles is arguably difficult to obtain, but the consequences of this idea make it worth considering. Besides producing, as already mentioned, a spectrum of $`\nu _{\mathrm{UHE}}`$ very different than those of astrophysical origin, with almost no neutrinos at low energy where bounds exist, this idea implies that the directions of UHECR could map the distribution of parent particles (which should coincide with the distribution of cold dark matter) at large red shifts. This is because the initial energy of the $`\nu _{\mathrm{UHE}}`$ decay product needs to be redshifted to the energy of the “$`Z`$-burst” in its way to the Earth. This idea of unstable superheavy relic particles could be constrained by the EGRET bound on the diffuse low-energy gamma ray flux resulting from the $`Z`$-bursts . This question deserves further study.
In an upcoming paper simulations will be presented for the photon, nucleon and neutrino fluxes coming from “$`Z`$-bursts” of $`0.6\times 10^{23}`$ eV, as would be if due to relic neutrinos of mass $`m_{\mathrm{SK}}`$. The “$`Z`$-bursts” were simulated using PYTHIA and the absorption of photons and nucleons was modelled using energy attenuation lengths provided by Bhattacharjee and Sigl supplemented by radio-background models by Protheroe and Biermann . Fig. 1
shows the result of this simulation with an approximate fit to the AGASA cosmic ray data. The photon lines labelled 1, 2 and 3 correspond to three different radio background models, which we take to well represent the uncertainty in these models. The fit to the AGASA data provides the normalization of the photon and nucleon fluxes that allows us to determine the (assumed homogeneous and isotropic) flux of ultra-high energy neutrinos at the $`Z`$-resonance energy ($`0.6\times 10^{23}`$ eV) to be
$$F_{\nu _{\mathrm{UHE}}}=1.510^{36}\frac{1}{eVm^2srsec},$$
(5)
if no lepton asymmetry is assumed. With a large lepton asymmetry this flux could be reduced by up to a factor of about 1/20.
Without making an assumption on the spectrum of $`\nu _{\mathrm{UHE}}`$, the “$`Z`$-burst” mechanism determines the $`\nu _{\mathrm{UHE}}`$ flux only at one point, the $`Z`$-resonance energy $`E_{Res}`$. It is interesting to point out that just on the basis of that one point, at $`0.6\times 10^{23}`$ eV, taking the results of the “Goldstone Experiment” at face value, the model of “$`Z`$-bursts” with relic neutrinos of mass $`m_{\mathrm{SK}}`$ would be already rejected, except for large lepton asymmetries.
The “Goldstone Experiment” searches for Lunar radio emissions from interactions of neutrinos (and cosmic rays) above $`10^{19}`$ eV of energy. The published preliminary results correspond to only 12 hours of observation and systematic errors affecting the bounds are not yet well understood. However, the present bounds show that the “Goldstone Experiment” will provide important constraints on “$`Z`$-burst” models in the near future.
Finally, let us address the issue of how the huge lepton asymmetries mentioned here could arise in the early Universe. Let us recall that, while charge neutrality imposes a lepton number asymmetry in electrons as large as the baryon asymmetry in protons, i.e. $`(n_en_{\overline{e}}/n_\gamma )=(n_Bn_{\overline{B}}/n_\gamma )10^{10}`$, no such restrictive bound operates on neutrinos.
A realistic model to produce very large lepton asymmetries without producing a large baryon asymmetry was presented in Ref.. The model requires the existence of bosons carrying lepton number (as in supersymmetric models), a period of inflation ending at a relatively low temperature, and a lepton asymmetry large enough for the electroweak symmetry to be spontaneously broken at all temperatures after inflation, which suppresses sphaleron transitions. Sphaleron transitions violate baryon plus lepton number ($`B+L`$) while preserving ($`BL`$), with the effect of producing $`B=L`$ if they are in equilibrium. Thus, in the presence of a very large lepton number $`L`$, sphaleron transitions in equilibrium would produce an equally large baryon number $`B`$ and this needs to be avoided in these models.
A model along similar lines is in Ref.. Lepton asymmetries can also be generated in neutrino oscillations after the electroweak phase transition, but it is unclear if an asymmetry of order one can arise in this fashion, .
Acknowledgements- This work was supported in part by the U.S. Department of Energy Grant No. DE-FG03-91ER40662, Task C.
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# References
Most of supersymmetry (SUSY) breaking models in supergravity (even including gauge mediation models) assume a separation of the SUSY-breaking and the SUSY standard-model sectors . However, the origin of the separation is not well understood, although such a separation is crucial to obtain phenomenologically consistent spectra for SUSY particles.
The brane world proposed by Randall and Sundrum provides a beautiful geometric explanation for the separation. That is, the hidden and observable sectors live on different three-dimensional branes separated by a gravitational bulk in higher dimensional spacetime. It has been, recently, claimed that this brane separation produces the hidden and observable separation in the “conformal” frame in supergravity, which was proposed long time ago from a phenomenological ground .
It is a crucial observation in Ref. that the above separation in the “conformal” frame induces a no-scale type Kähler potential<sup>1</sup><sup>1</sup>1 The no-scale supergravity adopts a specific form $`f_H=Z+Z^{}`$, where $`Z`$ is the superfield responsible for the SUSY breaking. We assume the Kähler potential for the $`Z`$ field to be of the form $`f_H=Z^{}Z+\mathrm{}`$, where the ellipsis denotes higher order terms. in the Einstein frame,
$`K(\mathrm{\Phi }_{\mathrm{obs}},\mathrm{\Phi }_{\mathrm{obs}}^{},\mathrm{\Phi }_{\mathrm{hid}},\mathrm{\Phi }_{\mathrm{hid}}^{})=3\mathrm{log}\left(1{\displaystyle \frac{1}{3}}f_O(\mathrm{\Phi }_{\mathrm{obs}},\mathrm{\Phi }_{\mathrm{obs}}^{}){\displaystyle \frac{1}{3}}f_H(\mathrm{\Phi }_{\mathrm{hid}},\mathrm{\Phi }_{\mathrm{hid}}^{})\right).`$ (1)
Here, $`\mathrm{\Phi }_{\mathrm{obs}}`$ and $`\mathrm{\Phi }_{\mathrm{hid}}`$ denote superfields in the observable and hidden sectors, respectively. With the above Kähler potential Eq. (1) we easily show that all soft SUSY-breaking masses and $`A`$ terms in the observable sector vanish in the limit of the zero cosmological constant. All gaugino masses in the observable sector also vanish because of the decoupling of the hidden superfield $`Z`$ from the gauge kinetic function .<sup>2</sup><sup>2</sup>2 The gaugino-mediated SUSY breaking was proposed in Ref. . See also recent works .
On the contrary, the SUSY-invariant $`\mu `$ term ($`\mu H\overline{H}`$) naturally arises from the Kähler potential if $`f_O`$ contains $`f_OH\overline{H}`$, where $`H`$ and $`\overline{H}`$ are chiral superfields of Higgs doublets. This mechanism produces the SUSY-invariant mass $`\mu `$ of the order of the gravitino mass, i.e. $`\mu m_{3/2}`$. This requires the gravitino mass $`m_{3/2}100\mathrm{GeV}1\mathrm{TeV}`$ for the correct electroweak symmetry breaking. With these gravitino masses the anomaly mediation generates too small SUSY-breaking masses in the observable sector and hence all SUSY particles in the observable sector except for the Higgsinos and the gravitino remain almost massless.
In this letter we introduce a new U(1) gauge interaction in the bulk to solve the above problem.<sup>3</sup><sup>3</sup>3 A similar model in the brane world has been also discussed in Ref. . The U(1)<sub>bulk</sub> gauge superfield plays a role of messenger between the hidden and the observable branes. SUSY-breaking effects on the hidden brane are transmitted to the observable brane through the bulk U(1)<sub>bulk</sub> gauge interaction and all of the gauginos, squarks and sleptons in the observable sector acquire suitable SUSY-breaking masses.
It is well known that a similar U(1) gauge interaction is also used as a messenger between the SUSY- breaking and observable sectors in a class of gauge-mediated SUSY breaking models . Thus, it is quite natural to identify the above bulk U(1)<sub>bulk</sub> gauge interaction with the messenger U(1)<sub>m</sub> gauge interaction in the gauge mediation models. We adopt a model proposed in Ref. , and interpret it as a low-energy effective theory of the brane world. We then show that rather small gauge coupling $`\alpha _{\mathrm{bulk}}5\times 10^4`$ of the U(1)<sub>bulk</sub> is required for a successful phenomenology. This small coupling is regarded as a consequence of a large volume of extra dimension. In fact, the result implies the compactification length $`L`$ of the extra dimension to be $`L^12\times 10^{15}\mathrm{GeV}`$ for $`(4+1)`$-dimensional spacetime and the fundamental scale $`M_{}`$ is determined as $`M_{}2\times 10^{17}\mathrm{GeV}`$ to reproduce the gravitational scale $`M_G(2M_{}^3L)^{1/2}2\times 10^{18}\mathrm{GeV}`$. Such a large compactification length is a crucial ingredient to suppress the flavor-changing neutral currents (FCNC’s) and hence our proposal is very consistent with the brane-world scenario of Randall and Sundrum .
Let us first briefly review the gauge mediation model proposed in Ref. . The model consists of three sectors: dynamical SUSY breaking (DSB) sector, messenger sector, and the minimal SUSY standard model (MSSM) sector. The DSB sector is based on a SUSY SU(2) gauge theory with four doublet chiral superfields $`Q_i`$ where $`i`$ is a flavor index ($`i=1,\mathrm{},4`$). Here, we have a global flavor SU(4)<sub>F</sub>. We assume the following tree-level superpotential introducing six singlet chiral superfields $`Z`$ and $`Z^a`$ ($`a=1,\mathrm{},5`$):
$`W_{\mathrm{tree}}=\lambda Z(QQ)+\lambda _ZZ^a(QQ)_a,`$ (2)
where $`(QQ)`$ and $`Z`$ are singlets of the SP(4)<sub>F</sub> subgroup of the flavor SU(4)<sub>F</sub> and $`(QQ)_a`$ and $`Z^a`$ are five-dimensional representations of the SP(4)<sub>F</sub>. As shown in Ref. , integration of the SU(2) gauge fields together with $`Q_i`$ and $`Z^a`$ leads to the low-energy effective superpotential
$`W_{\mathrm{eff}}{\displaystyle \frac{\lambda }{(4\pi )^2}}\mathrm{\Lambda }^2Z,`$ (3)
for $`\lambda _Z>\lambda `$, where $`\mathrm{\Lambda }`$ is a dynamical scale of the SU(2) gauge interaction.<sup>4</sup><sup>4</sup>4 The factor of $`4\pi `$ is determined by the naïve dimensional analysis . We have nonvanishing $`F`$ term, $`F_Z\lambda \mathrm{\Lambda }^2/(4\pi )^20`$, and hence SUSY is broken . We also assume that the fields in the DSB sector are charged under the U(1)<sub>m</sub> gauge interaction. The charge assignments of chiral superfields are given by
$$Q_1(+1),Q_2(1),Q_3(0),Q_4(0).$$
(4)
Here, the numbers in each parentheses denote the U(1)<sub>m</sub> charges. The U(1)<sub>m</sub> charges for $`Z`$ and $`Z^a`$ are determined such that the superpotential Eq. (2) is invariant under the U(1)<sub>m</sub>.
We now turn to the messenger sector. It consists of three chiral superfields,
$`E(+1),\overline{E}(1),S(0),`$ (5)
and vector-like messenger quark and lepton superfields, $`d`$, $`\overline{d}`$, $`l`$, $`\overline{l}`$. Here, the messenger quark multiplets $`d`$, $`\overline{d}`$ and lepton multiplets $`l`$, $`\overline{l}`$ are all neutral under the U(1)<sub>m</sub>. The $`d`$ and $`\overline{d}`$ ($`l`$ and $`\overline{l}`$) transform as the right-handed down quark and its antiparticle (the left-handed lepton doublet and its antiparticle) under the standard-model gauge group, respectively. The superpotential for the messenger sector is given by
$`W_{\mathrm{mess}}=k_ESE\overline{E}+{\displaystyle \frac{f}{3}}S^3+k_dSd\overline{d}+k_lSl\overline{l}.`$ (6)
The SUSY-breaking effects in the DSB sector are transmitted to the messenger sector through the U(1)<sub>m</sub> gauge interaction. As a result, the $`E`$ and $`\overline{E}`$ fields obtain positive soft SUSY-breaking squared masses $`m_E^2`$ and $`m_{\overline{E}}^2`$,<sup>5</sup><sup>5</sup>5 If $`Z=0`$, there is an unbroken U(1) $`R`$-symmetry. In this case, the U(1)<sub>m</sub> gaugino remains massless while $`E`$ and $`\overline{E}`$ fields obtain soft SUSY-breaking masses in Eq. (7).
$`m_E=m_{\overline{E}}{\displaystyle \frac{\alpha _m}{4\pi }}{\displaystyle \frac{\lambda F_Z}{\mathrm{\Lambda }}}{\displaystyle \frac{\alpha _m}{4\pi }}{\displaystyle \frac{\lambda ^2}{16\pi ^2}}\mathrm{\Lambda },`$ (7)
where $`\alpha _m=g_m^2/4\pi `$ is the U(1)<sub>m</sub> gauge coupling constant. They generate the following negative soft SUSY-breaking mass squared $`m_S^2`$ for the $`S`$ field through the Yukawa coupling $`k_ESE\overline{E}`$ in Eq. (6) at the one-loop level:
$`m_S^2{\displaystyle \frac{4}{(4\pi )^2}}k_E^2m_E^2\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }}{m_E}}.`$ (8)
Altogether, the resulting scalar potential is
$`V_{\mathrm{mess}}={\displaystyle \underset{\eta }{}}|{\displaystyle \frac{W_{\mathrm{mess}}}{\eta }}|^2+m_E^2|E|^2+m_{\overline{E}}^2|\overline{E}|^2m_S^2|S|^2,`$ (9)
where $`\eta `$ denotes chiral superfields $`E,\overline{E},S,d,\overline{d},l`$ and $`\overline{l}`$. This potential has a global minimum at
$`S^{}S={\displaystyle \frac{m_S^2}{2f^2}},|F_S|={\displaystyle \frac{m_S^2}{2f}},`$
$`E=\overline{E}=d=\overline{d}=l=\overline{l}=0,`$ (10)
in a certain parameter region . Thus, all the standard-model gauge symmetries are preserved and the SUSY-breaking effects are transmitted to the messenger quark and lepton multiplets through $`F_S`$.
The soft SUSY-breaking masses for the gauginos $`\stackrel{~}{g}_i`$ ($`i=1,\mathrm{},3`$) and the squarks, sleptons, and Higgses $`\stackrel{~}{f}`$ in the MSSM sector are generated by integrating out the messenger quarks and leptons as
$`m_{\stackrel{~}{g}_i}`$ $`=`$ $`c_i{\displaystyle \frac{\alpha _i}{4\pi }}\mathrm{\Lambda }_{\mathrm{mess}},`$ (11)
$`m_{\stackrel{~}{f}}^2`$ $`=`$ $`2\mathrm{\Lambda }_{\mathrm{mess}}^2\left[C_3({\displaystyle \frac{\alpha _3}{4\pi }})^2+C_2({\displaystyle \frac{\alpha _2}{4\pi }})^2+{\displaystyle \frac{5}{3}}Y^2({\displaystyle \frac{\alpha _1}{4\pi }})^2\right],`$ (12)
where $`c_1=5/3`$, $`c_2=c_3=1`$; $`C_3=4/3`$ for color triplets and zero for singlets, $`C_2=3/4`$ for weak doublets and zero for singlets, and $`Y`$ is the hypercharge ($`Y=Q_{\mathrm{em}}T_3`$). Here, $`\mathrm{\Lambda }_{\mathrm{mess}}`$ is an effective messenger scale defined as
$`\mathrm{\Lambda }_{\mathrm{mess}}{\displaystyle \frac{|F_S|}{|S|}}={\displaystyle \frac{m_S}{\sqrt{2}}},`$ (13)
which can be written in terms of the SUSY-breaking scale $`\sqrt{F_Z}`$ as
$`\mathrm{\Lambda }_{\mathrm{mess}}`$ $``$ $`{\displaystyle \frac{\sqrt{2}}{(4\pi )^4}}\alpha _m\lambda ^2k_E\sqrt{\mathrm{ln}{\displaystyle \frac{(4\pi )^3}{\alpha _m\lambda ^2}}}\mathrm{\Lambda }`$ (14)
$`=`$ $`{\displaystyle \frac{\sqrt{2}}{(4\pi )^3}}\alpha _m\lambda \sqrt{\lambda }k_E\sqrt{\mathrm{ln}{\displaystyle \frac{(4\pi )^3}{\alpha _m\lambda ^2}}}\sqrt{F_Z}.`$
We are now at the point of this letter. We consider that the above model is the low-energy effective theory of the brane world, in which all fields in the DSB sector reside on the hidden brane while the messenger and MSSM sectors are localized on the observable brane. Then, the U(1)<sub>m</sub> gauge multiplet should necessarily live in the bulk. Thus, we identify U(1)<sub>m</sub> with the bulk U(1)<sub>bulk</sub>.<sup>6</sup><sup>6</sup>6 We assume that the compactification scale $`L^1`$ is sufficiently high as $`L^1\mathrm{\Lambda }`$. The SUSY breaking on the hidden brane is transmitted to the observable brane by the U(1)<sub>bulk</sub> gauge and gravitational interactions across the bulk between two branes.
The effective messenger scale $`\mathrm{\Lambda }_{\mathrm{mess}}`$ should be taken at $`(10^410^5)\mathrm{GeV}`$ to induce the MSSM gaugino and sfermion masses of the electroweak scale. On the other hand, the supergravity effects generate simultaneously the $`\mu `$ term and the gravitino mass as discussed in the introduction ,
$`\mu m_{3/2}={\displaystyle \frac{F_Z}{\sqrt{3}M_G}}.`$ (15)
We should set $`\sqrt{F_Z}(26)\times 10^{10}\mathrm{GeV}`$ to reproduce correctly the electroweak symmetry breaking (i.e. $`\mu 100\mathrm{GeV}1\mathrm{TeV}`$). It is a crucial point that the above two conditions determine the U(1)<sub>bulk</sub> gauge coupling constant through Eq. (14). Assuming the Yukawa coupling constants $`\lambda `$ and $`k_E`$ connecting fields on the same brane to be of order one, we obtain the U(1)<sub>bulk</sub> gauge coupling $`\alpha _{\mathrm{bulk}}5\times 10^4`$ at the scale $`\mathrm{\Lambda }`$. Since the running effect of the U(1)<sub>bulk</sub> gauge coupling is negligible for the matter content in the present model, we find that $`\alpha _{\mathrm{bulk}}5\times 10^4`$ at the compactification scale of the extra dimension.
For a definiteness, we here assume the $`(4+1)`$-dimensional spacetime with one extra dimension compactified on the orbifold $`S^1/𝐙_2`$. The hidden and observable branes are located at two different fixed points of the $`S^1/𝐙_2`$ separated by a distance $`L`$. In this case, the 5-dimensional U(1)<sub>bulk</sub> multiplet is composed of a vector field, a Dirac spinor field and a real scalar field, which corresponds to a $`𝒩=2`$ vector multiplet in 4 dimensions. Through an orbifold projection, however, only the 4-dimensional $`𝒩=1`$ vector multiplet of the U(1)<sub>bulk</sub> can couple to the fields on two branes.<sup>7</sup><sup>7</sup>7 A detained analysis on the orbifold $`S^1/𝐙_2`$ is given in Ref. .
The 4-dimensional gauge coupling $`g_{\mathrm{bulk}}`$ is obtained from the 5-dimensional coupling $`g_{\mathrm{bulk}}^{(5)}`$ as
$`{\displaystyle \frac{1}{g_{\mathrm{bulk}}^2}}={\displaystyle \frac{2L}{(g_{\mathrm{bulk}}^{(5)})^2}}.`$ (16)
Assuming that $`g_{\mathrm{bulk}}^{(5)}`$ is of order one in the unit of the fundamental scale $`M_{}`$, we obtain the following relation:
$`(2L)^1=(4\pi \alpha _{\mathrm{bulk}})M_{},`$ (17)
which shows that the compactification length $`L`$ is $`(8\pi \alpha _{\mathrm{bulk}})^1100`$ times larger than the fundamental length scale $`M_{}^1`$ of the theory. This sufficiently suppresses the unwanted FCNC’s which would be induced by exchanges of bulk fields of masses around $`M_{}`$ .
On the other hand, the gravitational scale $`M_G`$ is given by
$`M_G^2=2M_{}^3L.`$ (18)
Thus, together with Eq. (17), we find that the fundamental scale $`M_{}`$ and the compactification scale $`L^1`$ are given by
$`M_{}=M_G(4\pi \alpha _{\mathrm{bulk}})^{1/2}2\times 10^{17}\mathrm{GeV},`$ (19)
$`L^1=2M_G(4\pi \alpha _{\mathrm{bulk}})^{3/2}2\times 10^{15}\mathrm{GeV}.`$ (20)
It is very interesting that these values are close to the ones discussed in Ref. .
Several comments are in order. First, the U(1)<sub>bulk</sub> gauge symmetry is unbroken in the present model so that one of scalar components of the $`E`$ and $`\overline{E}`$ fields is completely stable. This requires the reheating temperature $`T_R`$ of inflation to be lower than the mass of the lightest scalar field of order $`S10^5\mathrm{GeV}`$, in order for its present energy density not to exceed the critical density of the universe. Second, the dangerous $`D`$-term for the U(1)<sub>bulk</sub> does not appear in the model , since there is an unbroken charge conjugation symmetry defined as
$$\begin{array}{c}Q_1Q_2,Q_2Q_1,\\ \\ V_{\mathrm{bulk}}V_{\mathrm{bulk}},E\overline{E},\overline{E}E,\end{array}$$
where $`V_{\mathrm{bulk}}`$ is the U(1)<sub>bulk</sub> gauge superfield. The singlets $`Z`$ and $`Z^a`$ are assumed to transform properly so that the superpotential Eq. (2) is invariant under the charge conjugation symmetry Eq. (S0.Ex3). This is a consequence of the vector-like structure of the U(1)<sub>bulk</sub> gauge sector. Third, if there exists a non-Abelian gauge theory in the bulk other than the U(1)<sub>bulk</sub>, the radius of the extra dimension is stabilized as shown in Ref. .<sup>8</sup><sup>8</sup>8 The radius can also be stabilized by the mechanism of Ref. , if there are two non-Abelian gauge theories in the bulk which couple to suitable matters on a brane. Note that the stabilization is not disrupted by the Casimir energy induced by the SUSY breaking .
Finally, we should stress that the gravitino has a mass of order $`100\mathrm{GeV}1\mathrm{TeV}`$ in the present model although the mass spectrum of the other SUSY particles is the same as that of the gauge-mediated SUSY breaking models.<sup>9</sup><sup>9</sup>9 The $`B`$ term is of the order of the gravitino mass, where $`B`$ is defined as $`=\mu B\stackrel{~}{H}\stackrel{~}{\overline{H}}+\mathrm{h}.\mathrm{c}.`$ ($`\stackrel{~}{H}`$ and $`\stackrel{~}{\overline{H}}`$ are the scalar components of $`H`$ and $`\overline{H}`$). Thus, the small $`\mathrm{tan}\beta `$ region may be accommodated in the present model. This leads to an interesting phenomenological consequence. That is, the bino is most likely the dark matter in the present universe while the usual gauge mediation models predict the gravitino dark matter. It is a generic feature of the gauge-mediated SUSY breaking models in which the phenomenologically viable $`\mu `$ term arises from the supergravity effect . We should stress that the Randall-Sundrum brane-world scenario provides a natural solution to the $`\mu `$ problem in a large class of gauge-mediated SUSY breaking models , although we have adopted a specific model in Ref. to demonstrate our point.
Acknowledgments
We would like to thank Izawa K.-I. for a useful discussion. Y.N. thanks the Japan Society for the Promotion of Science for financial support. This work was partially supported by “Priority Area: Supersymmetry and Unified Theory of Elementary Particles (No. 707)” (T.Y.).
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# 1 Introduction
## 1 Introduction
An old and common idea in physics is that a particle makes its presence manifest via excitation of fields. If one puts a lot of particles together, one gets a macroscopic object, well described by classical physics, and correspondingly one expects the field excitations to be well described by a classical field theory. In particular, it seems obvious that wherever we trust this field theory as a good description of the low energy physics, a well-behaved solution to the field equations corresponding to that object should exist.
Type II string theory compactified on a Calabi-Yau manifold is described at low energies by a four dimensional $`𝒩=2`$ supergravity theory coupled to massles vector- and hypermultiplets. Quantum corrections to these theories are relatively well under control, and yet they are remarkably rich in content, with various intriguing connections to nontrivial physics and mathematics.
When string perturbation theory can be trusted, massive charged BPS particles in these theories can be described by D-branes wrapping nontrivial supersymmetric cycles in the Calabi-Yau manifold, or more generally by boundary states of the conformal field theory describing the relevant string perturbaton theory. When the low energy supergravity theory can be trusted, the same objects can be described by solutions to the field equations of motion. It turns out that not all charges support BPS states in the string theory, and that not all charges have BPS solutions in the supergravity theory. Thus, in suitable regimes, one naturally expects some sort of correspondence between supergravity solutions and D-branes.
Such a correspondence, while physically quite plausible, is in its consequences highly nontrivial. For instance, it would give rise to a number of powerful predictions about the existence of special Lagrangian submanifolds in Calabi-Yau manifolds, and the existence of boundary states in conformal field theories. However, as we will show in this paper, closer examination of this supposed correspondence reveals some intriguing puzzles. In particular, it turns out that the traditional assumption of the particle as a source localized in a single point of space leads to inconsistencies. Fortunately, once again, string theory finds its way out, and an interesting resolution to this paradox emerges.
The outline of this paper is as follows. In section 2, we briefly review the relevant geometry underlying low energy type IIB string theory compactified on a Calabi-Yau manifold. In section 3, the derivation of the attractor flow equations is revisited. We start from a duality invariant bosonic action, discuss an interpretation as a static string action, derive the spherically symmetric attractor flow equations in two different forms, and comment on a subtlety arising for vanishing cycles. In section 4, we analyze some properties of solutions, with special emphasis on conifold charges, leading to “empty holes”, and a short discussion of equal charge multicenter solutions. Then we tackle the existence issue: the attractor flow turns out to break down when the central charge has a regular zero, and this leads to a natural conjecture on the existence of BPS states . However, this natural conjecture leads to some puzzling paradoxes. This is illustrated in section 5, using the example of a certain BPS state at the Gepner point of the quintic, known to exist, but nevertheless having a regular zero of the central charge. A second puzzle is illustrated in the $`SU(2)`$ Seiberg-Witten model. In section 6, we propose a resolution to these puzzles; the key is to consider composite configurations. A thought experiment brings us rather naturally to the required configurations, in a spherical shell approximation. A stability check is made using test particle probes, and a representation as composite flows is given, making direct contact with 3-pronged strings in a suitble rigid limit. A smooth effective field theory picture for decay at marginal stability emerges, and Joyce’s stability conjecture for special Lagrangian submanifolds is recovered. We briefly comment on $`\mathrm{\Pi }`$-stability . Finally, in section 7, the analysis of general stationary multicenter solutions is initiated. Some properties of solutions can be inferred directly from the equations of motion. In particular, the intrinsic angular momentum of the multicenter composites is computed. Section 8 summarizes our conlusions, and indicates some open problems.
## 2 Geometry of IIB/CY compactifications
To establish notation and our setup, let us briefly review the low energy geometry of type IIB string theory compactified on a Calabi-Yau 3-fold. Some of the more technical elements of this section are only needed for the derivation of some more technical results further on.
We will follow the manifestly duality invariant formalism of . Consider type IIB string theory compactified on a Calabi-Yau manifold $`X`$. The four dimensional low energy theory is $`𝒩=2`$ supergravity coupled to $`n_v=h^{1,2}`$ massless abelian vectormultiplets and $`n_h=h^{1,1}+1`$ massless hypermultiplets, where the $`h^{i,j}`$ are the Hodge numbers of $`X`$. The hypermultiplet fields will play no role in the following and are set to zero.
The vectormultiplet scalars are given by the complex structure moduli of $`X`$, and the lattice of electric and magnetic charges is identified with $`H^3(X,)`$, the lattice of integral harmonic $`3`$-forms on $`X`$. The “total” electromagnetic field strength $``$ is (up to normalisation convention) equal to the type IIB self-dual five-form field strength, and is assumed to have values in $`\mathrm{\Omega }^2(M_4)H^3(X,)`$, where $`\mathrm{\Omega }^2(M_4)`$ denotes the space of 2-forms on the four dimensional spacetime $`M_4`$. The usual components of the field strength are retrieved by picking a symplectic basis $`\alpha ^I,\beta _I`$ of $`H^3(X,)`$:
$$=F^I\beta _IG_I\alpha ^I.$$
(2.1)
The total field strength satisfies the self-duality constraint:
$$=_{10},$$
(2.2)
where $`_{10}`$ is the Hodge star operator on the ten-dimensional space time, which factorises on the $`M_4\times X`$ compactification as $`_{10}=_4_X`$. To prevent overly heavy notation, we will also denote the Hodge dual in $`X`$ by a hat, so for any form $`\mathrm{\Gamma }`$ on $`X`$:
$$\widehat{\mathrm{\Gamma }}_X\mathrm{\Gamma }$$
(2.3)
Note that this operation is moduli-dependent. The constraint (2.2) relates the $`F`$ and $`G`$ components in (2.1). The (source free) equation of motion and the Bianchi identity of the electromagnetic field are combined in the equation $`d=0`$, implying locally the existence of a potential: $`=d𝒜`$.
The geometry of the vector multiplet moduli space, parametrized with $`n_v`$ coordinates $`z^a`$, is special Kähler . The (positive definite) metric
$$g_{a\overline{b}}=_a\overline{}_{\overline{b}}𝒦$$
(2.4)
is derived from the Kähler potential
$$𝒦=\mathrm{ln}(i_X\mathrm{\Omega }_0\overline{\mathrm{\Omega }}_0),$$
(2.5)
where $`\mathrm{\Omega }_0`$ is the holomorphic $`3`$-form on $`X`$, depending holomorphically on the complex structure moduli. It is convenient to introduce also the normalized 3-form
$$\mathrm{\Omega }=e^{𝒦/2}\mathrm{\Omega }_0.$$
(2.6)
The “central charge” of $`\mathrm{\Gamma }H^3(X,)`$ is given by
$$Z(\mathrm{\Gamma })_X\mathrm{\Gamma }\mathrm{\Omega }_\mathrm{\Gamma }\mathrm{\Omega },$$
(2.7)
where we denoted, by slight abuse of notation, the cycle Poincaré dual to $`\mathrm{\Gamma }`$ by the same symbol $`\mathrm{\Gamma }`$.
In the following we will frequently make use of the (antisymmetric, topological, moduli independent) intersection product:
$$\mathrm{\Gamma }_1,\mathrm{\Gamma }_2=_X\mathrm{\Gamma }_1\mathrm{\Gamma }_2=\mathrm{\#}(\mathrm{\Gamma }_1\mathrm{\Gamma }_2)$$
(2.8)
With this notation, we have for a symplectic basis $`\{\alpha ^I,\beta _I\}`$ by definition $`\alpha ^I,\beta _J=\delta _J^I`$. We will also often use the (symmetric, positive definite, moduli dependent) Hodge product:
$$\mathrm{\Gamma }_1,\widehat{\mathrm{\Gamma }_2}=\mathrm{\Gamma }_1,_X\mathrm{\Gamma }_2=_X\mathrm{\Gamma }_1_X\mathrm{\Gamma }_2.$$
(2.9)
When the $`\mathrm{\Gamma }_i`$ denote cohomology classes, their harmonic representative will always be assumed in (2.9).
Every harmonic $`3`$-form $`\mathrm{\Gamma }`$ on $`X`$ can be decomposed according to $`H^3(X,)=H^{3,0}(X)H^{2,1}(X)H^{1,2}(X)H^{0,3}(X)`$ as (for real $`\mathrm{\Gamma }`$):
$$\mathrm{\Gamma }=i\overline{Z}(\mathrm{\Gamma })\mathrm{\Omega }ig^{a\overline{b}}\overline{D}_{\overline{b}}Z(\mathrm{\Gamma })D_a\mathrm{\Omega }+c.c.,$$
(2.10)
where we introduced the Kähler covariant derivative on $`Z`$ and $`\mathrm{\Omega }`$:
$$D_a_a+\frac{1}{2}_a𝒦$$
(2.11)
This decomposition is orthogonal with respect to the intersection product (2.8), and diagonalizes the Hodge star operator:
$$_X\mathrm{\Omega }=i\mathrm{\Omega }\text{ and }_XD_a\mathrm{\Omega }=iD_a\mathrm{\Omega }$$
(2.12)
For further reference, we write down the following useful identities:
$`{\displaystyle _X}\mathrm{\Omega }\overline{\mathrm{\Omega }}`$ $`=`$ $`i`$ (2.13)
$`{\displaystyle _X}D_a\mathrm{\Omega }\overline{D}_{\overline{b}}\overline{\mathrm{\Omega }}`$ $`=`$ $`ig_{a\overline{b}}`$ (2.14)
$`(d+iQ+id\alpha )(e^{i\alpha }\mathrm{\Omega })`$ $`=`$ $`e^{i\alpha }D_a\mathrm{\Omega }dz^a,`$ (2.15)
where $`\alpha `$ is an arbitrary real function and $`Q`$ is the chiral connection:
$$Q=\mathrm{Im}(_a𝒦dz^a).$$
(2.16)
As an example of an application, one can easily check the following expressions for intersection and Hodge products:
$`\mathrm{\Gamma }_1,\mathrm{\Gamma }_2`$ $`=`$ $`2\mathrm{Im}[Z(\mathrm{\Gamma }_1)\overline{Z}(\mathrm{\Gamma }_2)+g^{a\overline{b}}D_aZ(\mathrm{\Gamma }_1)\overline{D}_{\overline{b}}\overline{Z}(\mathrm{\Gamma }_2)]`$ (2.17)
$`\mathrm{\Gamma }_1,\widehat{\mathrm{\Gamma }_2}`$ $`=`$ $`2\mathrm{Re}[Z(\mathrm{\Gamma }_1)\overline{Z}(\mathrm{\Gamma }_2)+g^{a\overline{b}}D_aZ(\mathrm{\Gamma }_1)\overline{D}_{\overline{b}}\overline{Z}(\mathrm{\Gamma }_2)],`$ (2.18)
## 3 The attractor flow equations revisited
We now turn to the investigation of 4d supergravity BPS solutions with charged sources corresponding to D3-branes wrapped around a nontrivial supersymmetric (i.e. special Lagrangian) 3-cycle $`\mathrm{\Gamma }`$ of $`X`$. In the mirror IIA picture this corresponds to BPS states with (mixed) 0-, 2-, 4- and 6-brane charge.
Such $`𝒩=2`$ supergravity solutions and the remarkable attractor mechanism emerging in this context were first studied, from supersymmetry considerations, in . An approach based on the bosonic action, which we will also follow here, was pioneered in . Further explorations were made in , and various solutions analyzed in . A rich connection with D-branes, geometry and arithmetic was discovered in . Some recent work on analogous phenomena in five dimensional theories includes .
Part of this section is a review of well known results, though the geometric, manifestly duality invariant setup we use may give a clarifying alternative point of view on some of these. Also, the strategy outlined here to obtain the BPS equations directly from the bosonic action will enable us in section 7 to do the same for the general stationary case (possibly non-static, with multiple centers having mutually nonlocal charges), adding further insight to the solutions of . Furthermore, some subtleties in the derivation of the flow equations will turn out to be relevant for a proper treatment of the solution for conifold charges later on, and finally, an interpretation of the reduced action as that of an effective stretched string will allow us to make contact with the $`317`$ brane description of BPS states in $`𝒩=2`$ QFT.
So we believe it’s worthwhile to revisit this derivation. However, the reader only interested in the resulting equations can safely skip the derivation and proceed with section 4.
### 3.1 Duality invariant formalism
The relevant bosonic part of the usual 4d low energy effective $`𝒩=2`$ supergravity action is, in 4d Planck units:
$$S_{4D}=\frac{1}{16\pi }_{M_4}d^4x\sqrt{G}R\mathrm{\hspace{0.17em}2}g_{a\overline{b}}dz^ad\overline{z}^{\overline{b}}\frac{1}{4\gamma ^2}_{M_4}F^IG_I$$
(3.1)
where $`\gamma `$ is a convention dependent number, $`F^I=dA^I`$ and the $`G_I`$ are obtained from the $`F^I`$ using the selfduality constraint (2.2). On the other hand, the bosonic 4d spacetime part of the low energy effective action of a probe D3-brane wrapped around a supersymmetric 3-cycle in the homology class $`\mathrm{\Gamma }`$, in a given background, is :
$$S_\mathrm{\Gamma }=|Z(\mathrm{\Gamma })|𝑑s+\frac{\sqrt{\pi }}{\gamma }\mathrm{\Gamma },𝒜,$$
(3.2)
with $`Z(\mathrm{\Gamma })`$ as in (2.7), $`d𝒜=`$, and $`,`$ denoting the intersection product (2.8). The integral is over the effective particle worldline.
Combining (3.1) and (3.2), assuming $`\mathrm{\Gamma }`$ to be electric with respect to the choice of symplectic basis (that is, $`\mathrm{\Gamma }`$ is a linear combination of the $`\alpha ^I`$), we see that an electromagnetic field produced by such a source with charge $`\mathrm{\Gamma }`$ satisfies, for any spatial surface $`S`$ surrounding the source:
$$_S=\sqrt{4\pi }\gamma \mathrm{\Gamma }$$
(3.3)
Now while the action (3.1) makes four dimensional general covariance manifest, it is not invariant under electromagnetic duality rotations (i.e. change of symplectic basis in (2.1)). A straightforward, manifestly covariant action exhibiting manifest duality invariance does not exist. In fact, since the 4D theory descends directly from type IIB supergravity, this problem is equivalent to the nonexistence of a straightforward generally covariant action for the self-dual four-form potential. However, a perfectly satisfactory, though not manifestly covariant action for self-dual forms has been known for quite a while , and this action (dimensionally reduced) will actually turn out to be very convenient for our purposes.
To write down this action for an arbitrary background metric, one introduces the usual shift and lapse vectors $`N_0`$ and $`N^i`$, putting the four dimensional metric in the form:
$$ds^2=N_0^2dt^2+G_{ij}(dx^i+N^idt)^2.$$
(3.4)
The shift vector determines a three dimensional 1-form $`𝐍=G_{ij}N^jdx^i`$. We will use boldface notation to refer to three dimensional quantities throughout. Thus $`𝐝=dx^i_i`$, the 3d Hodge dual (based on $`G_{ij}`$) is denoted by $`\mathbf{}`$, and the spatial part of the total electromagnetic field $``$ is
$$𝓕=_{ij}dx^idx^j=𝐝𝓐.$$
(3.5)
The $`H^3(X,)`$-valued 3-vector potential $`𝓐`$ is considered to be the fundamental variable (instead of the 4-vector potentials $`A^I`$ in the formulation based on (3.1)). The action obtained from with our compactification assumptions is then:
$$S_{e.m.}=\frac{1}{4\gamma ^2}dt_{M_3}_X𝓕_t𝓐(N_0𝓕\mathbf{}\widehat{𝓕}+𝐍\mathbf{}𝓕\mathbf{}𝓕)$$
(3.6)
The integral over $`X`$ yields simply the intersection product (2.8). Since the above expression does not refer to any choice of symplectic basis, it is indeed manifestly duality invariant. The equation of motion following from this action is the self-duality condition (2.2), with $`=d𝒜`$, where $`𝒜_0`$ arises as an integration constant.
Of course, since $`𝒜_0`$ does not exist off shell in this formulation, we can no longer use the coupling of the electromagnetic field to sources as in (3.2). Instead, its coupling to charges is implemented by imposing the constraint (3.3), which only involves the spatial fields. Again, no reference to a choice of basis is made. Note that the presence of charges will induce Dirac string singularities in $`𝓐`$, or require the introduction of a nontrivial bundle structure.
The coupling of the source to gravity and the scalars remains unchanged.
We will use (3.6) instead of the electromagnetic part of (3.1). In section 7 the full form of this action at nonzero $`𝐍`$ will be used to derive the BPS equations for the general stationary case, but it’s instructive (and sufficient for most of our purposes) to first consider some simpler cases.
### 3.2 Reduced action for static spherically symmetric configurations
In it was argued that time independent BPS configurations require a metric that can be expressed in the form
$$ds^2=e^{2U}(dt+\omega _idx^i)^2+e^{2U}dx^idx^i$$
(3.7)
This is the metric ansatz we will use throughout this paper. Usually we will also restrict to asymptotically flat spacetimes, i.e. $`U,\omega 0`$ at spatial infinity. Let us first consider static, spherically symmetric configurations. Then $`\omega =0`$ and $`U`$ is a function of the radial coordinate $`r=|𝐱|`$ only. Similarly we take the moduli $`z^a`$ to be function of $`r`$ only, and we can assume $`𝓕`$ to be of the form
$$𝓕=\frac{\gamma }{\sqrt{4\pi }}\mathrm{sin}\theta d\theta d\varphi \mathrm{\Gamma }$$
(3.8)
where $`\theta `$ and $`\varphi `$ are the usual angular coordinates and $`\mathrm{\Gamma }H^3(X,)`$ is the charge of the source. Then the total electromagnetic field is, with $`\tau 1/r`$:
$$=𝓕+_4\widehat{𝓕}=\frac{\gamma }{\sqrt{4\pi }}(\mathrm{sin}\theta d\theta d\varphi \mathrm{\Gamma }+e^{2U}d\tau dt\widehat{\mathrm{\Gamma }}),$$
(3.9)
which trivially satisfies the required equations of motion and duality constraints $`d=0`$ and $`=_4\widehat{}`$.
In terms of the inverse radial coordinate $`\tau =1/r`$, the total action per unit time reduces, under these assumptions, and dropping a total derivative proportional to $`\ddot{U}`$, simply to:
$$S_{eff}=S/\mathrm{\Delta }t=\frac{1}{2}_0^{\mathrm{}}𝑑\tau \{\dot{U}^2+g_{a\overline{b}}\dot{z}^a\dot{\overline{z}}^{\overline{b}}+e^{2U}V(z)\}(e^U|Z|)_{\tau =\mathrm{}}$$
(3.10)
where the dot denotes derivation with respect to $`\tau `$ and (cf. (2.18))
$`V(z)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{\Gamma },\widehat{\mathrm{\Gamma }}`$ (3.11)
$`=`$ $`|Z|^2+g^{a\overline{b}}D_aZ\overline{D}_{\overline{b}}\overline{Z}=|Z|^2+4g^{a\overline{b}}_a|Z|\overline{}_{\overline{b}}|Z|`$ (3.12)
with $`Z=Z(\mathrm{\Gamma })`$. The “potential” $`e^{2U}V(z)`$ is proportional to the electromagnetic energy density. The boundary term in (3.10) comes from (3.2). In principle, this reduced action has to be supplemented by the constraints coming from variations of the metric (consistent with spherical symmetry) other than the $`U`$ mode. In particular here this gives the constraint $`\dot{U}^2+\dot{z}^2e^{2U}V(z)=0`$. However, as we will see, this turns out to follow already from (3.10) (with the given source coupling and allowing arbitrary variations of the fields at $`\tau =\mathrm{}`$ in the variational principle). So we will simply proceed with the analysis of the action as it stands.
Note that (minus) this effective action per unit time can be interpreted as describing a nonrelativistic particle moving in $`(U,z)`$-space, subject to the potential $`e^{2U}V(z)`$, with time $`\tau `$ (fig. 1).
Only the initial ($`\tau =1/r=0`$) position of this effective particle is fixed. The $`\tau \mathrm{}`$ asymptotic behavior is given by requiring vanishing of the boundary terms at $`\tau =\mathrm{}`$ when varying the fields. This yields for $`\tau \mathrm{}`$
$`\dot{U}`$ $``$ $`e^U|Z|`$ (3.13)
$`\dot{z}^a`$ $``$ $`2e^Ug^{a\overline{b}}\overline{}_{\overline{b}}|Z|.`$ (3.14)
Incidentally, these are precisely the attractor flow equations, as we will see below. Note that this condition also implies the vanishing of the (conserved) effective particle’s energy, $`E_{eff}=\dot{U}^2+\dot{z}^2e^{2U}V(z)=0`$, which is precisely the “additional” constraint discussed earlier.
### 3.3 Interpretation as a static string action
In fact, the reduced action (3.11) can also be interpreted — and this is perhaps more natural — as a Polyakov action for a static string in $`(U,z)`$ space, with (variable) tension $`T=e^U\sqrt{V(z)}`$, in the background target space metric
$$ds^2=dt^2+dU^2+g_{a\overline{b}}dz^ad\overline{z}^{\overline{b}},$$
(3.15)
and worldsheet metric $`\text{diag}(T,T^1)`$ with respect to the worldsheet time resp. space coordinates $`t`$ and $`\tau `$. The vanishing of the effective particle’s energy $`E_{eff}`$ is equivalent to the Virasoro constraint, which can be used to transform the action to the Nambu-Goto form
$$S=𝑑t𝑑\tau e^U\sqrt{V}\sqrt{\dot{U}^2+\dot{z}^2}.$$
(3.16)
The asymptotic condition (3.13)-(3.14), being equivalent to the attractor flow equations, forces the endpoint of the string at $`\tau =\mathrm{}`$ to be fixed at an attractor point (see below). The other endpoint is fixed at the values of the moduli and $`U`$ at spatial infinity in the supergravity picture. The equations of motion determine the string to be a geodesic in $`(U,z)`$ space.
This interpretation, in a suitable rigid (i.e. gravity decoupling) limit of $`𝒩=2`$ supergravity (leading to Seiberg-Witten theory <sup>1</sup><sup>1</sup>1For an analysis of BPS solutions of pure low energy effective $`𝒩=2`$ SU(2) Yang-Mills theory, see . The solutions obtained there can be seen to be the rigid limit of supergravity attractor flows . and its generalizations), makes direct contact with the description of BPS states in $`𝒩=2`$ quantum field theories as stretched strings in a nontrivial background. For example, repeating the above analysis for a charge $`(n_e,n_m)`$ state in the pure $`SU(2)`$ Seiberg-Witten effective theory (or considering a suitable rigid limit of supergravity), one finds an effective reduced action with similar structure, with one modulus $`u`$, $`V(u)=|n_e+n_m\tau |^2/\mathrm{Im}\tau `$ (where $`\tau (u)`$ is the usual modular parameter of the SW Riemann surface), and evidently $`U0`$. For, say, a magnetic monopole, the attractor point turns out to be its vanishing mass point $`u=1`$ (see also section 5.2). Hence our effective string will be a geodesic stretched in the Seiberg-Witten plane, between an arbitrary modulus $`u_{\tau =0}`$ and $`u=1`$. Thus we arrive precisely at the picture of .
We will return to this point later on. In particular, the phenomenon of three-pronged strings appearing in this context will turn out to be related to the resolution of some intriguing paradoxes.
### 3.4 BPS equations of motion
The BPS equations of motion can be obtained from (3.10) by the usual Bogomol’nyi trick of completing squares. This can be done in two ways, yielding two different forms of the BPS equations. The first way is well known :
$`S_{eff}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _0^{\mathrm{}}}𝑑\tau \{(\dot{U}\pm e^U|Z|)^2+\dot{z}^a\pm 2e^Ug^{a\overline{b}}\overline{}_{\overline{b}}|Z|^2\}`$ (3.17)
$`\pm e^U|Z||_{\tau =0}^{\tau =\mathrm{}}(e^U|Z|)_{\tau =\mathrm{}}`$
leading to the BPS equations
$`\dot{U}`$ $`=`$ $`e^U|Z|`$ (3.18)
$`\dot{z}^a`$ $`=`$ $`2e^Ug^{a\overline{b}}\overline{}_{\overline{b}}|Z|.`$ (3.19)
This is the form of the equations found in . The other sign possibility is excluded by the asymptotic condition (3.13)-(3.14), or alternatively, by requiring physical acceptability: “wrong sign” solutions have runaway behavior, severe curvature singularities at finite distance, negative ADM mass, and are gravitationally repulsive.
The second way of squaring the action uses the Hodge product (2.9); if for a 3-form $`C`$ on $`X`$ we write
$$\mathbf{|}C\mathbf{|}^2=C,\widehat{C}$$
(3.20)
and we denote the position dependent phase of the central charge as
$$\alpha =\mathrm{arg}Z(\mathrm{\Gamma })$$
(3.21)
then we have, using (2.13)-(2.15):
$`S_{eff}`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle _0^{\mathrm{}}}𝑑\tau e^{2U}\mathbf{|}\mathrm{\hspace{0.17em}2}\mathrm{Im}[(_\tau +iQ_\tau +i\dot{\alpha })(e^Ue^{i\alpha }\mathrm{\Omega })]+\mathrm{\Gamma }\mathbf{|}^2`$ (3.22)
$`(e^U|Z|)_{\tau =0}`$
where $`Q_\tau =\mathrm{Im}(_a𝒦\dot{z}^a)`$, as in (2.16). (We have left $`U_{\tau =0}`$ arbitrary here, though in the asymptotically flat case this is zero of course.) Note that we take the holomorphic 3-form $`\mathrm{\Omega }`$ and the unnormalized holomorphic $`\mathrm{\Omega }_0`$ (cf. (2.6)) to be only dependent on the spacetime coordinates through the moduli $`z^a(\tau )`$. This is in contrast to refs. , where a (convenient) explicit position dependence of normalisation and phase of $`\mathrm{\Omega }_0`$ was chosen.<sup>2</sup><sup>2</sup>2One is free to make such a gauge choice, at least locally, since phase and normalisation of $`\mathrm{\Omega }_0`$ do not enter the action (3.1). However, from a Calabi-Yau geometrical point of view it is perhaps more natural to pick a dependence only through the moduli. In numerical computations, this has the further advantage that one can work with a fixed expression for $`\mathrm{\Omega }`$. Furthermore, in this way the phase $`\alpha `$ appears naturally in the equations, and this phase (which can be identified with the phase of the conserved supersymmetry ) will play a crucial role in the comparison with geometrical results on special Lagrangian manifolds .
The BPS equation following from (3.22) is again obtained by putting the square to zero:
$$2\mathrm{Im}[(_\tau +iQ_\tau +i\dot{\alpha })(e^Ue^{i\alpha }\mathrm{\Omega })]=\mathrm{\Gamma }.$$
(3.23)
However, by taking the intersection product with $`\mathrm{\Gamma }`$ on both sides of the equation, and using (3.21), it follows that $`Q_\tau +\dot{\alpha }=0`$, hence the BPS equation becomes simply
$$2_\tau \mathrm{Im}(e^Ue^{i\alpha }\mathrm{\Omega })=\mathrm{\Gamma }.$$
(3.24)
Conversely, by taking the intersection product of the latter equation with $`\overline{\mathrm{\Omega }}`$, using (2.15) and (2.13), and taking the imaginary part of the result, one obtains again $`Q_\tau +\dot{\alpha }=0`$, and hence (3.23). Now (3.24) readily integrates to
$$2\mathrm{Im}(e^Ue^{i\alpha }\mathrm{\Omega })=\mathrm{\Gamma }\tau +\mathrm{\hspace{0.17em}2}\mathrm{Im}(e^Ue^{i\alpha }\mathrm{\Omega })_{\tau =0}.$$
(3.25)
This is a powerful result, as it solves, in principle, the BPS equations of motion of the system. To bring it in a perhaps more familiar form, choose a symplectic basis $`\{\alpha ^I,\beta _I\}`$ of $`H^3(X,)`$, write $`\mathrm{\Gamma }=q_I\alpha ^I+p^I\beta ^I`$, define the holomorphic periods $`X^I=\alpha ^I,\mathrm{\Omega }_0`$, $`F_I=\beta _I,\mathrm{\Omega }_0`$, and take intersection products of this basis with the above equation. This gives:
$`2e^{U+𝒦/2}\mathrm{Im}(e^{i\alpha }X^I)={\displaystyle \frac{p^I}{r}}+c^I`$ (3.26)
$`2e^{U+𝒦/2}\mathrm{Im}(e^{i\alpha }F_I)={\displaystyle \frac{q_I}{r}}+d_I,`$ (3.27)
where we re-introduced $`r=1/\tau `$, and $`c^I,d_I`$ are constants. If, as in , we would pick an $`\mathrm{\Omega }_0`$-gauge where $`𝒦2U`$ and $`\alpha 0`$, one retrieves the expressions appearing in those references. Note also that the flow equations (3.18)-(3.19) are nothing but the projections of (3.23) on $`e^{i\alpha }\mathrm{\Omega }`$ resp. $`e^{i\alpha }D_a\mathrm{\Omega }`$.
Of course, finding the explicit flows in moduli space from (3.25) requires inversion of the periods to the moduli, which in general is not feasible analytically. However, in large complex structure approximations or numerically for e.g. the quintic, this turns out to be possible.
One final remark is in order. The solution (3.25) was derived from the action (3.22) under the implicit assumption that the quantity between brackets is not proportional to (the Poincaré dual of) a vanishing cycle $`\nu `$ (that is, a cycle for which the Hodge square $`\nu ,\widehat{\nu }=0`$, like the conifold cycle at a conifold point of moduli space). If that is the case, the Hodge square appearing in (3.22) is automatically zero, no matter what the expression inside the $`\mathbf{|}\mathbf{|}`$ is (as long as it is finitely proportional to a vanishing cycle). Actually, a we will see, such situations do occur naturally, and the previous remark should eliminate possible confusion there.
## 4 Attractors and existence of BPS states
### 4.1 Properties of some solutions
In what follows we will assume asymptotic flatness, i.e. $`U_{\tau =0}=0`$. Then solutions to (3.18)-(3.19) saturate the BPS bound
$$M=|Z_{\tau =0}|.$$
(4.1)
All mass is located in the fields: the “bare mass” contribution $`(e^U|Z|)_{\tau =\mathrm{}}`$ is zero. Indeed, (3.18) and (3.19) imply that both $`e^U`$ and $`|Z|`$ are monotonously decreasing functions satisfying the estimate $`e^U|Z||Z|/(1+|Z_{\mathrm{}}|\tau )`$, hence $`e^U|Z|0`$ when $`\tau \mathrm{}`$. More precisely, equation (3.19) implies
$$_\tau |Z|=4e^Ug^{a\overline{b}}_a|Z|\overline{}_{\overline{b}}|Z|0,$$
(4.2)
so the flows in moduli space described by (3.18) and (3.19) converge to minima of the central charge modulus $`|Z(\mathrm{\Gamma })|`$ (fig. 2).
Therefore the moduli values at the horizon are generically invariant under continuous deformations of the moduli at spatial infinity, and hence only dependent on the charge $`\mathrm{\Gamma }`$, a phenomenon referred to as the attractor mechanism. The attractor values of the moduli correspond to Calabi-Yau manifolds with very remarkable arithmetic properties, as explored in detail in .
#### 4.1.1 Black holes
The above estimate also implies that when $`Z_{\mathrm{}}0`$, the solution is a black hole with horizon at $`\tau =\mathrm{}`$. From the form of the metric (3.7) and direct analysis of equation (3.18) in the limit $`\tau \mathrm{}`$, the near horizon geometry can be seen to be $`AdS_2\times S^2`$:
$$ds_{NH}^2=\frac{r^2}{|Z_{\mathrm{}}|^2}dt^2+\frac{|Z_{\mathrm{}}|^2}{r^2}d𝐱^2.$$
(4.3)
The corresponding macroscopic entropy is $`A/4=\pi |Z_{\mathrm{}}|^2`$. These black holes have been studied extensively in the literature .
#### 4.1.2 Empty holes
When the D3-brane wraps a conifold cycle, i.e. a cycle vanishing at a conifold point, the minimal value of the central charge modulus is zero, and the above discussion of the generic case does no longer apply. However, since conifold cycles are known to exist (close to a conifold point) as special Lagrangian submanifolds, and therefore as physical BPS particles, it is natural to ask what the corresponding supergravity solution looks like.
Again, the flow in moduli space will converge to a point where $`|Z|`$ is minimal, which in this case is a point on the conifold locus, where $`Z`$ vanishes. At the conifold locus, the Calabi-Yau degenerates and we get an additional massless hypermultiplet in the low energy theory, so there we cannot necessarily trust our supergravity approximation. However, the results obtained are physically pleasing and interesting, so we will ignore this potential problem and proceed.
For simplicity, following , we will consider only one modulus, $`z`$, which we define to be the holomorphic ($`\mathrm{\Omega }_0`$) period of the vanishing cycle. Then the Kähler potential and metric close to the conifold point ($`z0`$) can be taken to be:
$`e^𝒦`$ $``$ $`k_1^2+{\displaystyle \frac{1}{2\pi }}|z|^2\mathrm{ln}|z|^2+k_2\mathrm{Re}z`$ (4.4)
$`g_{z\overline{z}}`$ $``$ $`{\displaystyle \frac{1}{2\pi k_1^2}}\mathrm{ln}|z|^2,`$ (4.5)
where $`k_1`$ and $`k_2`$ are positive constants. This geometry can be observed for instance close to the conifold point of the quintic, or (in a rigid limit) close to the massless monopole or dyon singularities in Seiberg-Witten theory.
The central charge of $`N`$ times the vanishing cycle close to $`z=0`$ is $`Z=\frac{N}{k_1}z`$ and the attractor flow equations in this limit are
$`\dot{U}`$ $`=`$ $`{\displaystyle \frac{N}{k_1}}e^U|z|`$ (4.6)
$`\dot{z}`$ $`=`$ $`2\pi k_1Ne^U{\displaystyle \frac{z}{|z|}}{\displaystyle \frac{1}{\mathrm{ln}|z|^2}},`$ (4.7)
with solution (approximately for $`z0`$) given by:
$`\mathrm{arg}z`$ $`=`$ const. (4.8)
$`|z|\mathrm{ln}|z|^1`$ $`=`$ $`\{\begin{array}{ccc}\pi k_1Ne^U_{}(\tau _{}\tau )\hfill & \text{for }& \tau <\tau _{}\hfill \\ 0\hfill & \text{for }& \tau \tau _{}\hfill \end{array}`$ (4.11)
$`U`$ $`=`$ $`{\displaystyle \frac{1}{4\pi k_1^2}}|z|^2\mathrm{ln}|z|^2+U_{},`$ (4.12)
where $`\tau _{}`$ and $`U_{}`$ are constants depending on initial conditions.
So here the attractor point $`z=0`$ is reached at finite nonzero coordinate distance $`r_{}=1/\tau _{}`$ from the origin. In the core region $`r<r_{}`$, the fields $`z`$ and $`U`$ are constant and the geometry is flat. There is no horizon and the core contains no energy ($`\dot{U}=\dot{z}=0`$ and $`V(z)=0`$), hence the name “empty hole” (fig. 3).
The solution $`z(\tau )`$ is once and $`U(\tau )`$ twice continuously differentiable at $`\tau =\tau _{}`$, so the curvature stays finite (and can be made arbitrary small by taking $`N`$ sufficiently large). However, higher derivatives diverge at $`\tau =\tau _{}`$, so strictly speaking the (two derivative) supergravity approximation breaks down here.<sup>3</sup><sup>3</sup>3This can perhaps be cured by including higher derivative terms or the new massless hypermultiplet in the effective action, smoothing out the solution, but presumably not changing it too much. Nevertheless, we believe this is a physically sensible solution. For instance, one could use it (in the straightforward multicenter extension given section 4.2) to compute the dynamics of a large number of slowly moving empty holes, in moduli space approximation, with sensible results. Note that due to the fact that these solutions are not black holes, one will obtain a moduli space geometry for nearly coincident centers which is completely different from the black hole one discussed in . In particular, there will presumably be no coalescence, in agreement with the physical expectation that no BPS bound states should exist for branes wrapping a conifold cycle .
We could also have used (3.25) to construct this solution (in particular cases this would in fact be a more powerful method to extract exact results, also away from the near conifold limit). However, naive application of this formula would lead to a field configuaration that is not constant inside the core: at $`\tau =\tau _{}`$, the central charge phase $`\alpha `$ jumps discontinuously from $`\alpha _{}`$ to $`\alpha _{}+\pi `$, and for $`\tau >\tau _{}`$, one gets the “solution” corresponding to the flow equations (3.18)-(3.19) with the opposite (=wrong) sign. As discussed in section 3.4, this is not an acceptable solution; it is not BPS, and physically ill-behaved.
The way out of this paradox is the remark given at the end of section 3.4: equation (3.25) needs only to be satisfied down to the radius where the conifold attractor point is reached. If we keep the fields constant below this radius, the BPS condition is automatically satisfied. This eliminates some confusion arising in the literature in this context.
Note that even though the solution (4.8)-(4.12) was derived in the near conifold limit, the conclusion that the attractor point is reached at finite $`\tau `$ is also true for moduli at infinity farther away from the conifold point, since the region where the approximation becomes valid will in any case be reached after finite $`\tau `$. Furthermore solutions at different $`N`$ are related by simple scaling; the core radius is proportional to $`N`$. So the solution will never be a black hole, no matter where we start in the moduli space, and no matter how many particles we put on top of each other. The attractor mechanism causes the mass to stay outside the Schwarzschild radius, protecting the configuration from gravitational collapse.
If the modulus $`z_0`$ at spatial infinity $`\tau =0`$ is sufficiently small, the core radius is given by $`r_{}=\frac{\pi k_1N}{|z_0|\mathrm{ln}|z_0|^1}`$. In the zero mass limit $`z_00`$, the core radius goes to infinity, leaving a completely flat space. If on the other hand one boosts up the particle while sending $`z_00`$, in such way that the total energy remains constant, one obtains in the limit $`z_00`$ the Aichelburg-Sexl shockwave metric for a massless particle moving at the speed of light. Again, this is physically sensible.
Finally note that we have derived the empty hole solution assuming all charge to be located at $`𝐱=0`$. However, exactly the same solution for $`U`$ and the moduli would have been obtained for any spherically symmetric charge distribution inside the core region. In particular the energy density and space curvature would have been the same. In that sense the charge is actually delocalized. It could for example be a spherical shell of radius $`r_{}`$ (this is perhaps the most natural location of the charges, as the “emptyness” of the core then becomes quite intuitive).
All this is of course very reminiscent of the enhançon mechanism of . One could say that the massless conifold particle is the “enhançon” curing the repulson singularity one would obtain for example by applying naively formula (3.25). The main difference is that there is no enhanced gauge symmetry in the core region, but rather an additional massless charged hypermultiplet.
It would be interesting to find out whether empty holes, like their black hole cousins , also have a Maldacena dual QFT description.
#### 4.1.3 No holes
As observed in , the flow equations (3.18)-(3.19) do not always have a solution: if the attractor point of the flow happens to be a simple zero of $`Z`$, at a regular point of moduli space, the flow will reach $`Z=0`$ at finite $`\tau =\tau _{}`$ and cannot be continued in a BPS way to the interior region $`\tau >\tau _{}`$ (see also fig. 4).
The basic difference with the previous case is the absence of the “damping” factor $`1/\mathrm{ln}|z|^2`$ in the inverse metric on the right hand side of (3.19) (or (4.7)), so that the constant field configuration at $`Z=0`$ is no longer a solution. On the other hand, by taking the charge sufficiently large, the curvature can be made again arbitrary small, so the absence of a supergravity solution should be quite meaningful.
Physically, one indeed doesn’t expect a BPS state with charge $`\mathrm{\Gamma }`$ to exist in a vacuum near a regular point where $`Z(\mathrm{\Gamma })=0`$: such a particle would be massless at $`Z=0`$, which (by integrating it out) should lead to a singularity in moduli space , in contradiction with the supposed regularity of the point under consideration.
### 4.2 Equal charge multicenter solutions
The single center configuration discussed above is readily extended to the multicenter case with equal (or parallel) charges in the centers. (Multicenter solutions with non-parallel charges are considerably more involved, and will be discussed in section 7.) One simply replaces $`\tau =1/|𝐱|`$ by
$$\tau \frac{1}{N}\underset{i=1}{\overset{N}{}}\frac{1}{|𝐱𝐱_i|},$$
(4.13)
where the $`𝐱_i`$ denote the positions of the particles, each with charge $`\mathrm{\Gamma }`$ (fig. 5).
Since the complete setup is formally the same as for the spherically symmetric case, so are the attractor flow equations. Therefore, everything said about the spherically symmetric case applies to the multicenter case as well.
For nearly coincident centers, the black hole near horizon geometry now becomes “fragmented” $`AdS_2\times S^2`$, as discussed in .
Though a proof for the general case is still lacking, it is expected that the moduli space geometry for the dynamics of slowly moving centers can be derived from the potential
$$L=d^3𝐱e^{4U}.$$
(4.14)
### 4.3 Existence of BPS states
The issue of existence of BPS states with given charge in theories with $`𝒩=2`$ supersymmetry in four dimensions is nontrivial and profound. The simplest example of such a theory is probably $`SU(2)`$ $`𝒩=2`$ Yang-Mills theory. The low energy dynamics of this theory was exactly solved in , where it was found that the BPS spectrum at weak coupling consists of the gauge boson and a tower of dyons of arbitrary integer electric charge and one unit of magnetic charge, while at strong coupling it consists solely of the magnetic monopole and the “elementary” dyon with one unit of electric and one unit of magnetic charge. Here “strong coupling” has the precise meaning of being inside a certain curve in the one dimensional moduli space, called the curve of marginal stability. At this curve, the various BPS particles have parallel central charges, so that they become only marginally stable against decay into constituents.
Similar phenomena are expected for type II string theory compactified on a Calabi-Yau 3-fold. Here the subject is intimately related to the existence of D-branes wrapped on supersymmetric cycles, since these are the objects that represent the BPS states. For instance in type IIB theory, at least in the large complex structure limit, existence of a BPS state of charge $`\mathrm{\Gamma }H^3(X,)`$ is equivalent with existence of a special Lagrangian submanifold in the homology class (Poincaré dual to) $`\mathrm{\Gamma }`$ . In type IIA theory, in the large volume limit, it is equivalent with existence of holomorphic submanifolds endowed with certain holomorphic bundles (or more precisely sheaves). At certain special points in moduli space, existence can be proven using the boundary state formalism. Recently, the problem has been studied intensively from various points of view: special Lagrangian submanifolds , holomorphic geometry and boundary states and low energy effective supergravity .
We will study this problem from the latter point of view, namely the low energy supergravity theory. The idea is as follows. If a certain charge supports a BPS state, one certainly would expect a corresponding 4d supergravity solution to exist, at least for sufficiently large charge, such that the supergravity approximation can be trusted. The converse statement is perhaps less evident, but with the knowledge that some charges indeed do not have BPS supergravity solutions (see section 4.1.3), it is quite tempting to conjecture an exact correspondence, at least for sufficiently large charges. Clearly, considering the degree of difficulty of the problem in other approaches, such a correspondence would be very powerful.
The above considerations were used in to arrive at the following proposal for an existence criterion for BPS states with given charge. Choose moduli $`z_{}^{a}{}_{0}{}^{}`$ at spatial infinity and a charge $`\mathrm{\Gamma }`$, and denote the minimal value of $`|Z(\mathrm{\Gamma })|`$ where the solution of (3.18)-(3.19) flows to by $`|Z|_{min}`$. There are three distinguished cases.
* Type a: $`|Z|_{min}0`$. The attractor flow exists and yields a regular BPS black hole solution. In this case one expects to have a BPS state in the theory with the given charge. Note that if the existence of a BPS state in a certain vacuum $`z_{}^{a}{}_{0}{}^{}`$ is thus established, it will also exist in any other vacuum that lies “upstream” the $`\mathrm{\Gamma }`$ attractor flow passing through the point $`z_{}^{a}{}_{0}{}^{}`$, where “upstream” means in the opposite direction of the flow given by (3.18)-(3.19). Since $`|Z|`$ has no maxima in moduli space , the upstream flows will tend to regions of moduli space at infinite distance, like the large complex structure limit. This also explains to a certain extent why BPS states are more likely to exist at large complex structure than in the bulk of modulispace.
* Type b: the flow tends to a singularity or a boundary of moduli space, where $`|Z|`$ might or might not vanish. More information is needed to decide whether the BPS state exists or not.
* Type c: $`|Z|_{min}=0`$, and this minimum is reached at a regular point in moduli space. As discussed in section 4.1.3, the flow breaks down and the charge is expected not to support a BPS state.
Though this proposed criterion works nicely for e.g. $`T^6`$ , it can fail in more general cases, as we will argue in the next section. More precisely, it turns out that some type c cases do correspond to BPS states present in the theory.
## 5 Puzzles and paradoxes
### 5.1 Puzzle 1: Solution suicide; states at the Gepner point of the quintic.
We start by considering the example of the quintic Calabi-Yau, first analyzed in great detail in . In particular, we will study BPS states in type IIB theory on the mirror quintic $`W`$ (or equivalently in type IIA on the quintic $`M`$ itself). This manifold can be obtained as a $`_5^3`$ quotient of the manifold in $`^4`$ given by the equation
$$W:x_1^5+x_2^5+x_3^5+x_4^5+x_5^55\psi x_1x_2x_3x_4x_5=0.$$
(5.1)
The transformation $`\psi \omega \psi `$ with $`\omega ^5=1`$ can be undone by a coordinate transformation $`x_1\omega ^1x_1`$, and thus the complex structure moduli space of $`W`$ can be parametrized by $`\psi ^5`$. The moduli space has three singularities: the Gepner point $`\psi ^5=0`$, which is a $`_5`$ orbifold singularity, the conifold point $`\psi ^5=1`$, where a 3-cycle vanishes, and the large complex structure limit $`\psi ^5=\mathrm{}`$, mirror to the large volume limit of the quintic.
In , building on , the D-brane spectrum of this theory was studied, mainly from the conformal field theory perspective. In particular the existence of a number of BPS states was established at the Gepner point $`\psi =0`$. These states were labeled as $`|00000_B,|10000_B,\mathrm{},|11111_B`$. The state $`|00000_B`$ corresponds to a D3-brane wrapped around the conifold cycle on the type $`IIB`$ side, and to a D6-brane on the type $`IIA`$ side.<sup>4</sup><sup>4</sup>4The identification of the type IIA D-brane charges depends on the chosen analytic continuation to large complex structure, so it has some intrinsic arbitrariness (see for some discussion of this point). It becomes massless at the conifold point $`\psi =1`$. The state $`|10000_B`$ has two units of D6-brane charge and five units of D2-brane charge in the type IIA theory, and the state $`|11000_B`$ has one unit of D6- and five units of D2-brane charge. The expected dimension of the deformation moduli space of these three states is respectively 0, 4 and 11.
According to the existence criterion of section 4.3, we should find “good” attractor flows with $`\psi =0`$ at spatial infinity for all these states; they should not be of type c. To address this question, one needs the exact moduli space geometry and central charges ($`\mathrm{\Omega }`$-periods) at arbitrary points in moduli space for the charges under consideration. From the results of , all this is indeed available, in terms of certain Meijer functions of the modulus $`\psi ^5`$. It is still hard then to tackle this problem analytically, but numerically using for instance Mathematica, it becomes quite tractable.
As an example, we show in fig. 6 the modulus of the central charge $`Z`$ as a function of $`\psi `$ for the state $`|11000_B`$. In this case we find indeed a nice regular BPS black hole solution, with $`|Z|_{min}1.61`$ at the attractor point $`\psi 0.85`$ (to make the supergravity approximation valid, we should actually put a large number $`N`$ of these charges on top of each other, but this simply rescales the solution). The same is true for $`|11100_B`$ ($`|Z|_{min}2.78`$ at $`\psi 0.51`$), for $`|11110_B`$ ($`|Z|_{min}4.58`$ at $`\psi 0.15`$), and for $`|11111_B`$ ($`|Z|_{min}7.43`$ at $`\psi 0.07`$). For $`|00000_B`$ we find an empty hole solution with attractor point $`\psi =1`$.<sup>5</sup><sup>5</sup>5Incidentally, in all cases we find $`|Z(\psi )|`$ to be symmetric under $`\psi \overline{\psi }`$, illustrating the rather special character of the boundary states constructed in .
However, as also noticed in , for the state $`|10000_B`$, we are in trouble. As shown in fig. 7, the attractor point $`\psi 1.46`$ is a regular zero of $`|Z|`$, so we have a type c situation: the supergravity solution does not exist!
Note that, though in conflict with the criterion of section 4.3, this result is not in contradiction with the physical expectation that a charge with $`Z=0`$ at a regular point cannot support BPS states in a neighborhood of that point: the zero $`\psi =\psi _{}1.46`$ and the Gepner point $`\psi =0`$, where the existence of the state is established, can be separated by a line of marginal stability where the state decays into lighter BPS constituents. So it is perfectly possible to have a BPS state with the given charge at $`\psi =0`$ and no such state close to $`\psi =\psi _{}`$.
Still, we clearly do have a physical problem here. This can be seen most clearly by considering the following thought experiment (fig. 8). Imagine a very large number of particles with the given charge on a huge sphere of radius $`R`$, in a vacuum with $`\psi =0`$. For $`R\mathrm{}`$ we expect to be allowed to neglect the backreaction of the particles on the bulk fields, and the description of these particles as CFT boundary states in the fixed background should be valid. Since we know from that the CFT boundary states corresponding to these particles exist and are BPS in the given vacuum, such a configuration should indeed exist and be BPS. Now give each of the particles the same very small inward velocity, and let us approximate the particle cloud as a uniformly charged spherical shell of adiabatically decreasing radius $`R`$. When the sphere becomes smaller, the collective backreaction becomes more important: outside the sphere, the fields will be given by the attractor flows (3.18)-(3.19); inside the sphere, the fields are constant. Note that this configuration is indeed BPS: the energy stored in the bulk fields outside the shell is $`E_{out}=|Z|_{r=\mathrm{}}(e^U|Z|)_{r=R}`$, the energy of the shell itself is $`E_{shell}=(e^U|Z|)_{r=R}`$, and the energy stored in the fields inside the shell is zero, adding up to a total energy $`E_{tot}=|Z|_{r=\mathrm{}}`$.
But if this motion goes on and nothing happens, we run into disaster: when $`R`$ becomes smaller than the nonzero radius $`r_{}`$ where the attractor flow breaks down, we no longer have a sensible solution! Moreover, by the physical argument given earlier, we actually expect that the particle cloud doesn’t even exist anymore at this point…
We propose a way out in section 6.
### 5.2 Puzzle 2: Monodromy murder; dyons in Seiberg-Witten theory
For our second (but closely related) puzzle, we consider the monopole in the Seiberg-Witten low energy effective theory for $`SU(2)`$ $`𝒩=\mathcal{2}`$ Yang-Mills . This theory can be obtained from the $`𝒩=2`$ supergravity theory describing the low energy physics of type II string theory compactified on a suitable Calabi-Yau manifold, in a certain rigid (=gravity decoupling) limit . The BPS solutions of this effective abelian theory (see for a discussion taking into account nonabelian corrections) can correspondingly be obtained as rigid limits of supergravity attractor flows . Because gravity decouples, $`U`$ is zero everywhere. The attractor flow equation for the modulus $`u(\tau )`$ is<sup>6</sup><sup>6</sup>6The factor $`\sqrt{2}`$ is due to the conventions used in . We take $`\mathrm{\Lambda }1`$.
$$\dot{u}=\sqrt{2}g^{u\overline{u}}_{\overline{u}}|Z|,$$
(5.2)
where $`g^{u\overline{u}}`$ is the inverse Seiberg-Witten metric and $`Z`$ is the central charge; for electric charge $`n_e`$ and magnetic charge $`n_m`$, this is $`Z=n_ea(u)+n_ma_D(u)`$, where $`a`$ and $`a_D`$ are given by certain hypergeometric functions .
Because $`Z(u)`$ is now analytic, the only possible minima of $`|Z|`$ are zeros. In fact, it is easy to see from (5.2) that the flows are lines of constant $`Z`$-phase, which of course necessarily end on a zero of $`Z`$. As before, to have a solution, the zero cannot be at a regular point, so it should be at the singularity $`u=1`$ where the monopole becomes massless, or at $`u=1`$ where the elementary dyon becomes massless. Therefore, the only solutions to (5.2) are of the empty hole type: the monopole, with attractor point $`u=1`$, and the elementary dyon, with attractor point $`u=1`$, plus of course their oppositely charged partners. In a neighborhood of their respective attractor points, with the choice of period cuts shown in fig. 9, the monopole has charge $`(n_e,n_m)=(0,1)`$, while the elementary dyon gets assigned the charge $`(1,1)`$ above the cut and $`(1,1)`$ below.
Again, we are facing a puzzle. It is well known that at weak coupling (that is, outside the line of marginal stability given by $`a_D/a`$), the BPS spectrum also contains a tower of dyons with $`n_m=\pm 1`$ and arbitrary (integer) $`n_e`$. These however correspond to “false flows” breaking down at a regular zero of $`Z=n_ea+n_ma_D`$ (on the line of marginal stability). The same problem arises for the purely electrically charged massive W-boson.
The paradox can be seen most sharply by starting with a $`(0,1)`$ monopole flow and performing a $`ue^{2i\pi }u`$ monodromy around $`u=\mathrm{}`$ (fig. 9). Doing this monodromy once should generate a higher dyon with charge $`(2,1)`$, doing it twice should yield a $`(4,1)`$ dyon, and so on. However, when circling around $`u=\mathrm{}`$, at a certain point, one arrives at a critical flow passing through the $`u=1`$ singularity. When one tries to “pull” the flow through this singularity, a catastrophe occurs: due to the nontrivial monodromy of the magnetic charge around $`u=1`$, the flow can no longer end on $`u=1`$; instead, past the singularity, it starts to diverge away from the flow just before criticality, and breaks down at the point (on the line of marginal stability) where $`Z_{(2,1)}`$ (analytically continued along the flow) becomes zero.
Physically, we don’t expect anything really drastic to happen when we vary the moduli at infinity just a little bit, yet we seem to find it can cause a complete breakdown of the solution.
So what is going on here?
## 6 Resolutions
### 6.1 Composite configurations
We now turn to the resolution of these puzzles. To get a first hint of what could do the job, consider again the situation described in section 5.1; the suicidal solution produced by an inmoving charged shell of charge $`\mathrm{\Gamma }`$. As explained there, we expect that the modulus at spatial infinity and the modulus where $`Z`$ becomes zero are separated by a line of marginal stability, so we expect the attractor flow to cross this line. Suppose that this is indeed what happens, say at $`r=r_{ms}`$, for the decay process $`\mathrm{\Gamma }\mathrm{\Gamma }_1+\mathrm{\Gamma }_2`$, and assume for simplicity that both $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ have well behaved attractor flows. Then what actually will happen when the shell shrinks is not the disaster scenario of section 5.1; instead, when the particle cloud has reached radius $`R=r_{ms}`$, the $`\mathrm{\Gamma }`$-particles will decay into $`\mathrm{\Gamma }_1`$\- and $`\mathrm{\Gamma }_2`$-particles (fig. 10). Now the $`\mathrm{\Gamma }_1`$\- and $`\mathrm{\Gamma }_2`$-particles cannot both continue to move inward, as this would be energetically impossible (the configuration would no longer be BPS because $`Z(\mathrm{\Gamma }_1)`$ and $`Z(\mathrm{\Gamma }_2)`$ acquire different phases on points of the $`\mathrm{\Gamma }`$-flow beyond the marginal stability line). Rather, the $`\mathrm{\Gamma }_1`$-particles will stay at the marginal stability locus $`r=r_{ms}`$, while the $`\mathrm{\Gamma }_2`$-particles move on, or vice versa. In the first case, when the $`\mathrm{\Gamma }_2`$-charges arrive at $`r=0`$, we have a BPS configuration (see below) consisting of a $`\mathrm{\Gamma }_2`$-charged center surrounded by a $`\mathrm{\Gamma }_1`$-shell at $`r=r_{ms}`$. Outside the shell the fields are given by the $`\mathrm{\Gamma }`$-attractor flow, and inside the shell by the $`\mathrm{\Gamma }_1`$-attractor flow. In the second case, we have a similar situation, with 1 and 2 interchanged.
To see that such configurations are indeed BPS, let us compute the total energy, say for the first case. The energy in the bulk fields outside the $`\mathrm{\Gamma }_1`$-shell is $`E_{out}=|Z(\mathrm{\Gamma })|_{\mathrm{}}(e^U|Z(\mathrm{\Gamma })|)_{r_{ms}}`$. The energy of the shell itself is $`E_{shell}=(e^U|Z(\mathrm{\Gamma }_1)|)_{r_{ms}}`$. The energy inside the shell is $`E_{in}=(e^U|Z(\mathrm{\Gamma }_2)|)_{r_{ms}}`$. So the total energy is
$$E_{tot}=|Z(\mathrm{\Gamma })|_{\mathrm{}}+\left(e^U(|Z(\mathrm{\Gamma }_1)|+|Z(\mathrm{\Gamma }_2)||Z(\mathrm{\Gamma }_1+\mathrm{\Gamma }_2)|)\right)_{r_{ms}}.$$
(6.1)
But since precisely at marginal stability, the quantity between brackets is zero, we find indeed $`E_{tot}=|Z(\mathrm{\Gamma })|_{r=\mathrm{}}`$, that is, the configuration is BPS.
Furthermore, when one would move the shell away from $`r=r_{ms}`$, the quantity between brackets becomes strictly positive, so this configuration is stable under such perturbations.
Another way of seeing this is by considering the force on a test particle of charge $`ϵ\mathrm{\Gamma }_1`$ at rest in the attractor flow field of a charge $`\mathrm{\Gamma }_2`$. This can be derived from (3.2). As shown in appendix A, the result is that this force can be derived from the potential
$$W=2ϵe^U|Z(\mathrm{\Gamma }_1)|\mathrm{sin}^2(\frac{\alpha _1\alpha _2}{2}),$$
(6.2)
where $`\alpha _i=\mathrm{arg}Z(\mathrm{\Gamma }_i)`$. This potential is everywhere positive, and becomes zero when $`\alpha _2=\alpha _1`$, that is, at marginal stability.
It is not difficult to extract the equilibrium radius $`r_{ms}`$ from the integrated flow equation (3.25). Taking the intersection product of $`\mathrm{\Gamma }_1`$ with this equation gives, denoting $`Z(\mathrm{\Gamma }_i)`$ in short as $`Z_i`$:
$$2\mathrm{Im}(e^Ue^{i\alpha }Z_1)=\mathrm{\Gamma }_1,\mathrm{\Gamma }\tau +\mathrm{\hspace{0.17em}2}\mathrm{Im}(e^{i\alpha }Z_1)_{\tau =0}.$$
(6.3)
At $`1/\tau =r=r_{ms}`$, the left hand side is zero, so
$$r_{ms}=\frac{\mathrm{\Gamma }_1,\mathrm{\Gamma }}{2\mathrm{Im}(e^{i\alpha }Z_1)_{r=\mathrm{}}}.$$
(6.4)
Using $`e^{i\alpha }=Z/|Z|`$ with $`Z=Z_1+Z_2`$ and $`\mathrm{\Gamma }_1,\mathrm{\Gamma }=\mathrm{\Gamma }_1,\mathrm{\Gamma }_2`$, this can be written more symmetrically as
$$r_{ms}=\frac{1}{2}\mathrm{\Gamma }_1,\mathrm{\Gamma }_2\frac{|Z_1+Z_2|}{\mathrm{Im}(\overline{Z_2}Z_1)}|_{r=\mathrm{}}.$$
(6.5)
Some interesting consequences of this identity will be discussed further on.
Having arrived at this picture of composite configurations in the approximation of spherical shells, a natural question to ask is whether supergravity also has nonspherical multicenter BPS solutions (with nonparallel charges). We will study this problem in section 7.
### 6.2 Forked flows
The composite configurations discussed above can be represented by composite flows (or “forked flows”): the flow starts as an ordinary $`\mathrm{\Gamma }`$-attractor flow, reaches a line of marginal stability, and then splits in a $`\mathrm{\Gamma }_1`$-flow and a $`\mathrm{\Gamma }_2`$-flow, corresponding to the two possible realization of the state as a charged center surrounded by a charged shell. The total energy of the configuration then equals the sum of the energies associated to each of the constituent flows, that is, for a $`\gamma `$-flow running from $`i`$ to $`f`$, $`E=(e^U|Z(\gamma )|)_f(e^U|Z(\gamma )|)_i`$.
Thus the generalization to composite spherically symmetric BPS states simply amounts to the generalization of simple attractor flows to composite attractor flows. Can we find such composite flows for the specific examples discussed in sections 5.1 and 5.2? Fortunately, it turns out we can. As shown in fig. 11, the $`\mathrm{\Gamma }=(2,1)`$ dyon in Seiberg-Witten theory can be realized as a flow splitting in a $`\mathrm{\Gamma }_1=(0,1)`$ monopole flow and a $`\mathrm{\Gamma }_2=2(1,1)`$ elementary dyon flow. This corresponds, in the supergravity regime with $`\mathrm{\Gamma }=N(2,1)`$, $`N`$ large, to a magnetic core with charge $`N(0,1)`$ surrounded by a dyonic shell with charge $`2N(1,1)`$, or vice versa. The intersection product of an elementary dyon and a monopole equals $`2`$.
For the quintic example outlined in section 5.1, we find a composite flow ending on two copies of the conifold point (fig. 11). In the conventions and notation of , the state $`|10000_B`$ under consideration has type IIA D-brane charge $`(Q_6,Q_4,Q_2,Q_0)=(2,0,5,0)`$, while the charges with vanishing mass at the two conifold point copies under consideration are $`(4,3,14,10)`$ and $`(6,3,19,10)`$,<sup>7</sup><sup>7</sup>7These slightly unnaturally looking values arise because the type IIA D-brane charges are naturally defined only at large volume (or large complex structure on the IIB mirror). Charges at arbitrary $`\psi `$ are defined by continuous transport coming from large $`\psi `$ in the wedge $`0<\mathrm{arg}\psi <2\pi /5`$. This procedure assigns charge $`(1,0,0,0)`$ to the state with vanishing mass at $`\psi =1`$. The states with vanishing mass at the other four copies of the conifold point get charges related to this one by the $`_5`$ monodromy around the Gepner point , which has no reason to have a particularly nice action when expressed in the type IIA D-brane basis. adding up to the required $`(2,0,5,0)`$. The intersection product of these two charges equals $`5`$.
The appearance of these composite flows is very reminiscent of the appearance of “3-pronged strings” in the “3-1-7 brane picture” of BPS states in $`𝒩=2`$ quantum field theories . This is no coincidence. As explained in section 3.3, in the Seiberg-Witten case for instance, there is an exact map between the attractor flows and the stretched strings of . Similarly, the composite flows arise precisely when the simple geodesic strings fail to exist and the 3-pronged strings take over, and here again there is an exact map between the flows and the strings.
Finally note one could imagine more complex configurations, involving more than one shell, corresponding to more than one flow split. For now, we will stick to the two charge case however.
### 6.3 Monodromy magic
This picture also offers a nice way to resolve the monodromy puzzle of section 5.2. Consider again $`N`$ monopole charges at $`r=0`$, in a vacuum $`u_{\mathrm{}}`$ such that the attractor flow is infinitesimally close to the critical one. When we further vary $`u_{\mathrm{}}`$ counterclockwise, the flow will pass through the $`u=1`$ point, say at $`r=r_c`$. Placed at this radius, an elementary dyon would be massless, so it could be created there at no cost in energy. And this is precisely what will happen when we continue to rotate $`u_{\mathrm{}}`$: $`2N`$ elementary dyons of charge $`(1,1)`$ are created! Due to the subtleties associated with monodromy, this is in full agreement with charge conservation. To get some physical feeling for this phenomenon, suppose we manipulate the $`u(r)`$-field in a certain region of space containing a piece of the surface $`r=r_c`$, in such way that here the $`u(r)`$-flow moves from passing just above to passing just below the $`u=1`$ singularity in the $`u`$-plane (fig. 12). Now imagine that just before the move a virtual monopole-antimonopole pair was created, to be destroyed again just after the move, and that the monopole happened to be at $`r>r_c`$ when the critical trajectory was crossed, while the antimonopole was at $`r<r_c`$. Then the spacetime trajectory of the monopole-antimonopole pair, mapped to moduli space via $`u(r,t)`$, encircles the point $`u=1`$. Consequently, there is a monodromy on the monopole charge, of which the net result is that we are left with two elementary dyons (with infinitesimal mass) when the monopole-antimonopole pair is destroyed again! If the resulting configuration is energetically favorable (so certainly if it is BPS), it will persist. This gives a physically reasonable mechanism to get the required dyons for the composite BPS state that is supposed to take over when we further rotate $`u_{\mathrm{}}`$. Once the dyon shell is present, we can continue the monodromy (making the dyon shell massive), till the (composite) flow passes through the $`u=1`$ singularity, where the above process is repeated with massless monopoles.<sup>8</sup><sup>8</sup>8The reader might be puzzled about how our configuration with a dyonic outer shell gets transformed into one with a magnetic outer shell. This can be understood from the discussion of empty holes in section 4.1.2: when approaching the flow passing through $`u=1`$, the distance between the dyonic shell and the “enhançon” radius $`r=r_{}`$, where $`u=1`$ is reached and below which the monopoles cannot be localized, shrinks to zero. Thus, at the critical flow, the roles of the monopole and dyon shells can be interchanged continuously.. In this way we can continue, creating all expected higher dyons.
Note that a local observer, placed in- or outside the sphere $`r=r_c`$, will not note anything peculiar when the transition takes place. Locally, everything changes perfectly smoothly.
### 6.4 Marginal stability, Joyce transitions and $`\mathrm{\Pi }`$-stability
From (6.5), it follows that when the moduli at infinity approach the line (or, if the dimension of moduli space is larger than one, the hypersurface) of marginal stability for the decay $`\mathrm{\Gamma }\mathrm{\Gamma }_1+\mathrm{\Gamma }_2`$, the shell radius $`r_{ms}`$ will diverge, eventually reaching infinity at marginal stability. This gives a nicely continuous 4d spacetime picture for the decay of the state when crossing marginal stability.
Furthermore, (6.5) tells us at which side of the marginal stability hypersurface the composite state can actually exist: since $`r_{ms}>0`$, it is the side satisfying
$$\mathrm{\Gamma }_1,\mathrm{\Gamma }_2\mathrm{sin}(\alpha _1\alpha _2)>0,$$
(6.6)
where $`\alpha _i=\mathrm{arg}Z(\mathrm{\Gamma }_i)_{r=\mathrm{}}`$. Sufficiently close to marginal stability, this reduces to
$$\mathrm{\Gamma }_1,\mathrm{\Gamma }_2(\alpha _1\alpha _2)>0,$$
(6.7)
which is precisely the stability condition for “bound states” of special Lagrangian 3-cycles found in a purely Calabi-Yau geometrical context by Joyce! (under more specific conditions, which we will not give here) .
Note also that, since the right hand side of (6.3) can only vanish for one value of $`\tau `$, the composite configurations we are considering here will actually satisfy
$$|\alpha _1\alpha _2|<\pi .$$
(6.8)
Another immediate consequence of (6.5) is the fact that these composite configurations can only<sup>9</sup><sup>9</sup>9at least for asymptotically flat space. For a space asymptotic to $`AdS_2\times S^2`$, the situation changes. occur for mutually nonlocal charges, that is, charges $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ with nonzero intersection product.
If the constituent $`\mathrm{\Gamma }_i`$ of the composite configuration for which $`\mathrm{\Gamma },\mathrm{\Gamma }_i>0`$ can be identified with a “subobject” of the state as defined in , the above conditions imply that the phases satisfy the $`\mathrm{\Pi }`$-stability criterion introduced in that reference. Though this similarity is interesting, it is far from clear how far it extends. $`\mathrm{\Pi }`$-stability is considerably more subtle than what emerges here. On the other hand, we have thus far only considered BPS configurations in a classical, spherical shell approximation, so also on this side the full stability story can be expected to be more complicated. We leave this issue for future work.
## 7 The general stationary multicenter case
In view of the emergence of composite BPS configuration in the spherical approximation, it is natural to look for more general multicenter solutions. This case is far more involved however. In particular, we have to give up the assumption that the configuration is static, and allow for more general, but still stationary, spacetimes.
### 7.1 BPS equations
Stationary (single center) BPS solutions of $`𝒩=2`$ supergravity were first studied in from supersymmetry considerations, in a specific space-dependent $`\mathrm{\Omega }_0`$-gauge (essentially the one described at the end of section 3.4). Here we will follow an approach based on the bosonic duality invariant action, similar to the one followed in section 3, and we let $`\mathrm{\Omega }_0`$ depend on position only through the the moduli.
Again, we will use the metric ansatz (3.7), but now with $`U`$ an arbitrary function of position $`𝐱`$, and $`\omega `$ not necessarily zero (but still time independent). We consider only the asymptotically flat case here, that is, $`U,\omega 0`$ when $`r\mathrm{}`$.
We will use boldface notation for 3d quantities as explained in section 3.1. The 3d Hodge dual with respect to the flat metric $`\delta _{ij}`$ will be denoted by $`\mathbf{}_0`$ , and for convenience we write $`\stackrel{~}{\omega }e^{2U}\omega `$. It will also turn out to be useful to define the following scalar product of spatial 2-forms $`𝓕`$ and $`𝓖`$:
$$(𝓕,𝓖)\frac{e^{2U}}{1\stackrel{~}{\omega }^2}_X𝓕[\mathbf{}_0\widehat{𝓖}\mathbf{}_0(\stackrel{~}{\omega }\widehat{𝓖})\stackrel{~}{\omega }+\mathbf{}_0(\stackrel{~}{\omega }\mathbf{}_0𝓖)].$$
(7.1)
Note that we have $`(𝓕,𝓖)=(𝓖,𝓕)`$ and for $`\stackrel{~}{\omega }`$ not too large $`(𝓕,𝓕)0`$.
With these assumptions and notations, the action (3.1), with the duality invariant electromagnetic action (3.6) substituted in place of the covariant one, becomes, putting $`\gamma \sqrt{4\pi }`$ and dropping a total derivative $`\mathrm{\Delta }U`$ from the gravitational action:
$`S_{4D}={\displaystyle \frac{1}{16\pi }}{\displaystyle 𝑑t_^3}`$ $`\{\mathrm{\hspace{0.17em}2}𝐝U\mathbf{}_0𝐝U{\displaystyle \frac{1}{2}}e^{4U}𝐝\omega \mathbf{}_0𝐝\omega `$ (7.2)
$`+2g_{a\overline{b}}𝐝z^a\mathbf{}_0𝐝\overline{z}^{\overline{b}}+(𝓕,𝓕)\}.`$
We will derive the BPS equation by “squaring” the action in a way inspired by (3.22). Let $`\alpha `$ be an arbitrary real function on $`^3`$, denote
$$𝐃𝐝+i(𝐐+𝐝\alpha +\frac{1}{2}e^{2U}\mathbf{}_0𝐝\omega ),$$
(7.3)
with $`𝐐`$ as in (2.16), and define the 2-form $`𝓖`$ as
$$𝓖𝓕2\mathrm{Im}\mathbf{}_0𝐃(e^Ue^{i\alpha }\mathrm{\Omega })+2\mathrm{Re}𝐃(e^Ue^{i\alpha }\mathrm{\Omega }\omega ),$$
(7.4)
Then we find for the integrand $``$ of (7.2), after some calculational effort involving repeated use of the identities (2.12) and (2.13)-(2.15),
$``$ $`=`$ $`(𝓖,𝓖)\mathrm{\hspace{0.17em}4}(𝐐+𝐝\alpha +{\displaystyle \frac{1}{2}}e^{2U}\mathbf{}_0𝐝\omega )\mathrm{Im}𝓖,e^Ue^{i\alpha }\mathrm{\Omega }`$ (7.5)
$`+𝐝[\mathrm{\hspace{0.17em}2}\stackrel{~}{\omega }(𝐐+𝐝\alpha )+4\mathrm{Re}𝓕,e^Ue^{i\alpha }\mathrm{\Omega }].`$
Thus if
$`𝓖`$ $`=`$ $`0`$ (7.6)
$`𝐐+𝐝\alpha +{\displaystyle \frac{1}{2}}e^{2U}\mathbf{}_0𝐝\omega `$ $`=`$ $`0,`$ (7.7)
we have a BPS solution to the equations of motion following from the reduced action (7.2) <sup>10</sup><sup>10</sup>10We will assume that these solutions also satisfy the equations of motion of the full action without restrictions on the metric, as in the spherically symmetric case, though we did not check this explicitly. (we will verify the saturation of the BPS bound below). Now from (7.7) and (7.3), we have $`𝐃=𝐝`$, and (7.6) becomes
$$𝓕+2𝐝\mathrm{Re}(e^Ue^{i\alpha }\mathrm{\Omega }\omega )=2\mathbf{}_0𝐝\mathrm{Im}(e^Ue^{i\alpha }\mathrm{\Omega }).$$
(7.8)
Since by construction $`𝐝𝓕=0`$ (away from sources), this implies $`𝐝\mathbf{}_0𝐝\mathrm{Im}(e^Ue^{i\alpha }\mathrm{\Omega })=0`$, so we can write
$$2\mathrm{Im}(e^Ue^{i\alpha }\mathrm{\Omega })=H,$$
(7.9)
with $`H`$ a $`H^3(X,)`$-valued harmonic function on $`^3`$ (possibly with source singularities). If we take the sources to be at positions $`𝐱_i`$ with charges $`\mathrm{\Gamma }_i`$, where $`i=1,\mathrm{},N`$, then from (7.8) and (3.3), we obtain
$$H=\underset{i=1}{\overset{N}{}}\mathrm{\Gamma }_i\tau _i+\mathrm{\hspace{0.17em}2}\mathrm{Im}(e^{i\alpha }\mathrm{\Omega })_{r=\mathrm{}},$$
(7.10)
with $`\tau _i=1/|𝐱𝐱_i|`$. Defining the 1-form
$$𝜻𝐝H,\mathrm{\Omega }=\underset{i=1}{\overset{N}{}}Z(\mathrm{\Gamma }_i)𝐝\tau _i,$$
(7.11)
we get from taking intersection products of $`𝐝H`$ given by (7.9) with $`\mathrm{\Omega }`$ and $`D_a\mathrm{\Omega }`$, and using (2.13)-(2.15):
$`𝐐+𝐝\alpha `$ $`=`$ $`e^U\mathrm{Im}(e^{i\alpha }𝜻)={\displaystyle \frac{1}{2}}e^{2U}𝐝H,H`$ (7.12)
$`𝐝U`$ $`=`$ $`e^U\mathrm{Re}(e^{i\alpha }𝜻)`$ (7.13)
$`𝐝z^a`$ $`=`$ $`e^Ug^{a\overline{b}}e^{i\alpha }\overline{D}_{\overline{b}}𝜻.`$ (7.14)
Using (7.12), equation (7.7) can be rewritten as:
$$\mathbf{}_0𝐝\omega =𝐝H,H.$$
(7.15)
Equations (7.9) and (7.10) generalize (3.25). Given the sources and the moduli at infinity, they yield the fields $`U(𝐱)`$, $`\alpha (𝐱)`$ and $`z^a(𝐱)`$. Equation (7.15) on the other hand gives $`\omega (𝐱)`$ (up to gauge transformations $`\omega \omega +𝐝f`$, which can be absorbed by a coordinate transformation $`ttf`$). Equations (7.13) and (7.14) generalize the flow equations (3.18)-(3.19).
Note that asymptotically for $`1/\tau =r\mathrm{}`$, the right hand side of (7.12) vanishes and $`\zeta _iZ(\mathrm{\Gamma }_i)𝐝\tau `$, implying
$$\alpha \mathrm{arg}Z(\mathrm{\Gamma })\text{ and }\zeta Z(\mathrm{\Gamma })𝐝\tau \text{when }r\mathrm{},$$
(7.16)
where $`\mathrm{\Gamma }=_i\mathrm{\Gamma }_i`$. Thus, far from all sources, we have again a simple attractor flow, corresponding to the total charge $`\mathrm{\Gamma }`$, as could be expected physically. In particular (7.13) gives $`𝐝U|Z(\mathrm{\Gamma })|d\tau `$, with $`\tau =1/r`$, establishing the saturation of the BPS bound on the mass:
$$M_{ADM}=|Z(\mathrm{\Gamma })|_{r=\mathrm{}}.$$
(7.17)
In the spherically symmetric case (and in the multicenter case with parallel charges), the above asymptotics become exact, and we retrieve the equations found earlier for those cases. Similarly, close to the center $`𝐱_i`$, we have
$$\alpha \mathrm{arg}Z(\mathrm{\Gamma }_i)\text{ and }\zeta Z(\mathrm{\Gamma }_i)𝐝\tau _i\text{when }𝐱𝐱_i,$$
(7.18)
and again we have asymptotically the flow equations for a single charge attractor, as could be expected physically. In particular the moduli at $`𝐱_i`$ will be fixed at the $`\mathrm{\Gamma }_i`$-attractor point.
The BPS equations of motion for the moduli and the metric obtained here can be seen to reduce to the equations found in in the $`\mathrm{\Omega }_0`$-gauge described at the end of section 3.4, except that we do not find the restricition $`\mathrm{𝐝𝐐}=0`$.
### 7.2 Some properties of solutions
Consider a multicenter solution, with distinct centers $`𝐱_i`$, $`i=1,\mathrm{},n`$, to the BPS equations
$`2e^U\mathrm{Im}(e^{i\alpha }\mathrm{\Omega })`$ $`=`$ $`H,`$ (7.19)
$`\mathbf{}_0𝐝\omega `$ $`=`$ $`𝐝H,H,`$ (7.20)
where
$$H=\underset{i=1}{\overset{n}{}}\mathrm{\Gamma }_i\tau _i+\mathrm{\hspace{0.17em}2}\mathrm{Im}(e^{i\alpha }\mathrm{\Omega })_{r=\mathrm{}},$$
(7.21)
as derived in the previous section. Acting with $`𝐝\mathbf{}_0`$ on equation (7.20) gives
$$0=\mathrm{\Delta }H,H,$$
(7.22)
so, using (7.21) and $`\mathrm{\Delta }\tau _i=4\pi \delta ^3(𝐱𝐱_i)`$, we find that for all $`i=1,\mathrm{},n`$:
$$\underset{j=1}{\overset{n}{}}\frac{\mathrm{\Gamma }_i,\mathrm{\Gamma }_j}{|𝐱_i𝐱_j|}=2\mathrm{Im}(e^{i\alpha }Z(\mathrm{\Gamma }_i))_{\mathrm{}}.$$
(7.23)
In the particular case of one source with charge $`\mathrm{\Gamma }_2`$ at $`𝐱=0`$ and $`m`$ sources with equal charge $`\mathrm{\Gamma }_1`$ at positions $`𝐱_i`$, this becomes
$$|𝐱_i|=\frac{\mathrm{\Gamma }_1,\mathrm{\Gamma }_2}{2\mathrm{Im}(e^{i\alpha }Z(\mathrm{\Gamma }_1))_{\mathrm{}}},$$
(7.24)
which is equal to the equilibrium distance $`r_{ms}`$ found in the spherical shell picture, equation (6.4).
In general, the moduli space of solutions to (7.23) will be quite nontrivial. Some general properties can be deduced relatively easy however. For instance, in a configuration made of only two different charge types $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ (distributed over an arbitrary number of centers), the charges of different type, if mutually nonlocal, will be driven to infinite distance from each other when $`(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2)`$-marginal stability is approached. This is similar to what we found for the spherical shell case. The stability condition (6.6) reappears as well. If on the other hand the two charges are mutually local (zero intersection), no BPS configuration exists with the two charges separated from each other, unless their phases are equal, that is, at marginal stability (then we can place the charges anywhere).<sup>11</sup><sup>11</sup>11If we consider a spacetime asymptotic to $`AdS_2\times S^2`$ instead of the asymptotically flat one we are assuming here, mutually local charges are no longer constrained by (7.23), because there will be an additional factor $`\mathrm{exp}[U(r=\mathrm{})]0`$ on the right hand side of (7.23).
For configurations made of more charge types, things get more complicated, but we will not go into this here.
Finally, note that because of the asymptotics discussed at the end of the previous section, we can expect the image of the moduli fields into moduli space for a multicenter configuration with only two different charge types to look like a fattened version of the composite flows we introduced earlier to represent the composite spherical shell configurations (fig. 13). Furthermore, we can expect that the more spherically symmetric the multicenter configuration becomes, the more this fattened version will approach the one dimensional composite flow. This is similar to what was found for spatial descriptions of dyons in $`𝒩=4`$ (effective) quantum field theories .
### 7.3 Angular momentum
It is well known from ordinary Maxwell electrodynamics that multicenter configurations with mutually nonlocal charges (e.g. the monopole-electron system) can have intrinsic angular momentum even when the particles are at rest. The same turns out to be true here.
We define the angular momentum vector $`𝐉`$ from the asymptotic form of the metric (more precisely of $`\omega `$) as
$$\omega _i=2ϵ_{ijk}J^j\frac{x^k}{r^3}+O(\frac{1}{r^3})\text{for }r\mathrm{}.$$
(7.25)
Plugging this expression in (7.20) and using (7.21) and (7.23), we find
$$𝐉=\frac{1}{2}\underset{i<j}{}\mathrm{\Gamma }_i,\mathrm{\Gamma }_j𝐞_{ij},$$
(7.26)
where $`𝐞_{ij}`$ is the unit vector pointing from $`𝐱_j`$ to $`𝐱_i`$:
$$𝐞_{ij}=\frac{𝐱_i𝐱_j}{|𝐱_i𝐱_j|}.$$
(7.27)
Just like in ordinary electrodynamics, this is a “topological” quantity: it is independent of the details of the solution and quantized in half-integer units (more precisely, when all charges are on the z-axis, $`2J_z`$).
The appearance of intrinsic configurational angular momentum implies that quantization of these composites will have some nontrivial features.
## 8 Conclusions
We have shown the emergence of some puzzles and paradoxes arising when one tries to construct four dimensional low energy effective supergravity solutions corresponding to certain BPS states in type II string theory compactified on a Calabi-Yau manifold, and demonstrated how these can be resolved by considering composite and extended configurations. We made connections to the enhançon mechanism, the 3-pronged string picture of QFT BPS states, $`\mathrm{\Pi }`$-stability and Joyce transitions of special Lagrangian manifolds. The problem was analyzed in a spherical shell approximation and by considering multicenter BPS solutions.
There are quite some problems however, new and old ones, we didn’t touch upon. The most prominent one is that we didn’t analyze to what extent these states really exist as BPS bound states in the full quantum theory. It seems quite likely that we now face the opposite problem we started with: instead of too little, we might now have too many possible solutions. In view of the nontriviality of quantum mechanics with mutually nonlocal charges, it is not unconceivable that a proper semiclassical treatment would eliminate some of these spurious solutions. But even at the classical level the existence issue is not completely settled. We did not show for instance that all solutions to (7.23) actually lead to well-behaved BPS solutions to the equations of motion; the same phenomenon causing the breakdown of some naively expected spherically symmetric solutions, namely hitting a zero of the central charge, could cause naively expected solutions to break down in this more complicated setting as well.
In this setup it seems also quite possible that a certain charge can have several different realizations as a BPS solution in a given vacuum, for example both as a single center and as a two center configuration. Crossing a line of marginal stability could then cause one realization to disappear, while leaving the other intact. The D-brane analog of this would presumably be a “jump” in its moduli space. This brings us to another interesting open question: is there a connection between D-brane moduli spaces and supergravity solution moduli spaces? And could those solution moduli spaces (for asymptotically flat or $`AdS_2\times S^2`$ spacetimes) teach us something about black hole entropy?
It could also be worthwhile to further explore the relation with $`\mathrm{\Pi }`$-stability, briefly mentioned in section 6.4.
Finally, this and other recent work illustrates an apparently recurrent theme in string theory: the resolution of singularities by creation of finitely extended D-brane configurations. It would be interesting to find out what the dielectric, non-commutative D-brane effects of can teach us about the states described in this paper.
Acknowledgements
I would like to thank Michael Douglas, Brian Greene, Calin Lazaroiu, Gregory Moore, Robert Myers, Christian Römelsberger, Walter Troost and Eric Zaslow for useful discussions and correspondence.
## Appendix A Potential for a test charge
From (3.2), it follows that the Lagrangian (with respect to the time coordinate $`t`$) for a test charge $`\mathrm{\Gamma }_t`$ at rest in the attractor flow field of a charge $`\mathrm{\Gamma }`$ is (denoting $`Z(\mathrm{\Gamma }_t)`$ in short as $`Z_t`$, and similarly for the other quantities involved)
$$L=e^U|Z_t|\frac{\sqrt{\pi }}{\gamma }\mathrm{\Gamma }_t,𝒜_0$$
(A.1)
where $`𝒜_0`$ is obtained from (3.9):
$$_i𝒜_0=\frac{\gamma }{\sqrt{4\pi }}e^{2U}_i\tau \widehat{\mathrm{\Gamma }}.$$
(A.2)
From (3.24), we get
$$\mathrm{\Gamma }=i_\tau (e^Ue^{i\alpha }\mathrm{\Omega })+\text{c.c.},$$
(A.3)
so, using (2.15) and (as shown in section 3.4) $`Q_\tau +\dot{\alpha }=0`$:
$$\mathrm{\Gamma }=ie^U\dot{U}e^{i\alpha }\mathrm{\Omega }+ie^Ue^{i\alpha }D_a\mathrm{\Omega }\dot{z}^a+\text{c.c.},$$
(A.4)
hence from (2.12) and again (2.15):
$`\widehat{\mathrm{\Gamma }}`$ $`=`$ $`e^U\dot{U}e^{i\alpha }\mathrm{\Omega }e^Ue^{i\alpha }D_a\mathrm{\Omega }\dot{z}^a+\text{c.c.}`$ (A.5)
$`=`$ $`e^{2U}_\tau (e^Ue^{i\alpha }\mathrm{\Omega })+\text{c.c.}.`$ (A.6)
Therefore
$`_i({\displaystyle \frac{\sqrt{\pi }}{\gamma }}\mathrm{\Gamma }_t,𝒜_0)`$ $`=`$ $`{\displaystyle \frac{1}{2}}e^{2U}_i\tau \mathrm{\Gamma }_t,\widehat{\mathrm{\Gamma }}`$ (A.7)
$`=`$ $`_i\left(e^U\mathrm{Re}(e^{i\alpha }Z_t)\right),`$ (A.8)
and thus (up to a constant)
$`L`$ $`=`$ $`e^U|Z_t|+e^U\mathrm{Re}(e^{i\alpha }Z_t)`$ (A.9)
$`=`$ $`e^U|Z_t|\left(1\mathrm{cos}(\alpha _t\alpha )\right)`$
$`=`$ $`2e^U|Z_t|\mathrm{sin}^2({\displaystyle \frac{\alpha _t\alpha }{2}}).`$
The force on the test particle is $`F_i=_iL`$, so we find for the force potential, as announced in section 6.1:
$$W=2e^U|Z_t|\mathrm{sin}^2(\frac{\alpha _t\alpha }{2}).$$
(A.11)
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# Far-infrared c-axis conductivity of flux-grown Y1-xPrxBa2Cu3O7 single crystals studied by spectral ellipsometry
## Abstract
The far-infrared c-axis conductivity of flux-grown Y<sub>1-x</sub>Pr<sub>x</sub>Ba<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> single crystals with 0.2$``$x$``$0.5 has been studied by sepectral ellipsometry. We find that the c-axis response exhibits spectral features similar to deoxygenated underdoped YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub>, i.e., a pseudogap develops in the normal state, the phonon mode at 320 cm<sup>-1</sup> exhibits an anomalous T-dependence and an additional absorption peak forms at T. This suggests that the T<sub>c</sub> suppression in flux grown Pr-substituted crystals is caused by a decrease of the hole content and/or by carrier localization rather than by pair breaking.
PACS: 74.25.Gz, 74.72.Bz, 78.20.-e
Superconductivity above the liquification temperature of nitrogen was first discovered in 1987 in the cuprate high T<sub>c</sub> compound YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> (Y-123) with a critical temperature of T<sub>c</sub>=92 K . It was subsequently shown that chemical substitution of Y by most rare-earth elements as well as by La leaves T<sub>c</sub> almost unaffected with T<sub>c</sub>$`>`$90 K . Ever since, the reason why Pr-123 does not become superconducting (SC) has been a puzzle . A number of models have been proposed which explain the absence of SC in Pr-123 either in terms of pair-breaking by the magnetic Pr-moments , a hole depletion of the CuO<sub>2</sub> planes caused by a mixed valency of Pr<sup>3+/4+</sup> or a partial substitution of Pr<sup>3+</sup> for Ba<sup>2+</sup> , a modification of the charge transfer between CuO chains and CuO<sub>2</sub> planes , or a hole redistribution and localization due to a strong hybridization of the Pr 4f and O 2p orbitals . Recently, it has been reported that Pr-123 can also be made SC with T<sub>c</sub>$`>`$90 K . This finding has led to renewed interest in the electronic properties of the Pr-123 compound. Yet the essential mechanism which decides whether Pr-123 is a superconductor with T<sub>c</sub>$`>`$90 K or a non SC insulator is still unknown.
In this paper we report ellipsometric measurements of the far-infrared (FIR) c-axis conductivity of partially Pr-substituted flux grown Y<sub>1-x</sub>Pr<sub>x</sub>Ba<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> crystals with 0.2$`x`$0.5. We show that the c-axis conductivity of the Pr-substituted crystals exhibits virtually the same spectral features as deoxygenated and thus underdoped YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> crystals, i.e. a spectral pseudogap forms in the normal state, the oxygen bond-bending phonon mode at 320 cm<sup>-1</sup> exhibits an anomalous T-dependence, and an additional broad absorption peak forms at low T. This finding implies that a similar effect causes the T<sub>c</sub> suppression in flux grown Pr substituted and deoxygenated Y-123 crystals, namely a depletion and/or a localization of the Zhang-Rice-type hole carriers of the CuO<sub>2</sub> planes. From our measurements we do not obtain any decisive information about the hole depletion mechanism in Pr-substituted samples.
The growth of the Pr-substituted Y<sub>1-x</sub>Pr<sub>x</sub>Ba<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> crystals has been described previously . Farily large crystals with a typical size of the ac-face of 3 by 0.5-1 mm have been used. They have been annealed in an O<sub>2</sub> gas stream for 1 day (d) at 600 C then cooled within 1d to 400 C and further annealed for 10 d at 400 C. A Zn-substituted YBa<sub>2</sub>Cu<sub>2.94</sub>Zn<sub>0.06</sub>O<sub>6.95</sub> crystal has been annealed under similar conditions but with the final step at 500 C for 5d. The midpoint SC transition temperature T<sub>c</sub> and the 10 to 90% halfwidth $`\mathrm{\Delta }`$T<sub>c</sub> have been determined by dc-SQUID magnetisation measurements at 5 Oe. The Pr-content and the Zn-content have been determined by EDX analysis. For the Pr-substituted crystals we obtained T<sub>c</sub>=81(3) K for x$``$0.2, T<sub>c</sub>=64(3) K for x$``$0.3, T<sub>c</sub>=48(4) K for x$``$0.4, and T<sub>c</sub>=24(5) K for x$``$0.5. For the Zn-substituted crystal we obtained T<sub>c</sub>=76(3) K for z$``$0.06.
The ellipsometric measurements have been performed at the National Synchrotron Light Source (NSLS) using a home built ellipsometer attached to a Nicolet Fast-Fourier spectrometer at the U4IR beamline . Some experiments have been done with the ellipsometer attached to a Bruker 113V using a conventional Hg arc light source. The optical measurements have been performed on the as grown clean ac surfaces (with the c-axis in the plane of incidence). The technique of ellipsometry provides significant advantages over conventional reflection methods in that (a) it is self-normalizing and does not require reference measurements and (b) the real and the imaginary parts of the dielectric function, $`\epsilon `$=$`\epsilon _1`$+i$`\epsilon _2`$, are obtained directly without a Kramers-Kronig transformation. Since only relative intensities of the reflected light are required, the ellipsometric measurements are more accurate and reproducible than conventional reflection measurements.
Figure 1 shows the spectra of the FIR c-axis conductivity, $`\sigma _{1c}`$, of the Y<sub>0.8</sub>Pr<sub>0.2</sub>Ba<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> crystal with T<sub>c</sub> =81(3) K in the normal and in the superconducting state. The room temperature spectrum is composed of an almost frequency independent electronic background on which five infrared active phonon modes are superimposed. The phonon modes correspond to those of fully oxygenated YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> . It has been previously shown for deoxygenated YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> that the strength of the phonon mode at 630 cm$`^1,`$ which corresponds to the vibration of apical oxygen neighboring an empty chain fragment, is a good indicator for the oxygen deficiency of a given sample. The absence of the 630 cm<sup>-1</sup> phonon mode thus confirms the fully oxidized state of our present crystal. With decreasing T the electronic background can be seen to develop a spectral gap in the normal state already well above T<sub>c</sub>. The characteristic features of this normal state spectral gap are identical to those of the pseudogap which has been observed in an underdoped deoxygenated YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6.75</sub> crystal with similar T<sub>c</sub> . In both crystals the normal state pseudogap has a similar size of $`\omega _{NG}`$700-800 cm<sup>-1</sup> (defined as the onset of the suppression of the conductivity with decreasing temperature) and its onset temperature T<sub>NG</sub> is around 200 K. Even the spectral shape of the pseudogaps compares very well. Close to T<sub>c</sub>, the in-plane oxygen bond-bending mode at 320 cm<sup>-1</sup> becomes strongly renormalized, its spectral weight decreases and its position is shifted by about 10 cm<sup>-1</sup> towards lower energies. Simultaneously, an additional broad peak appears around 500 cm<sup>-1</sup> as indicated by the arrow. Once more, the same spectral features have been observed in the FIR c-axis response of underdoped deoxygenated Y-123 . The broad low-T peak and the related strong anomaly of the 320 cm<sup>-1</sup> phonon mode have been successfully explained in terms of a model where the bilayer cuprate compounds like Y-123 are treated as a superlattice of intra- and interbilayer Josephson junctions . Within this model, the broad low-T peak corresponds to the transverse optical Josephson plasmon which arises from the out of phase oscillation of the intra-and the interbilayer longitudinal plasmon modes. The strong anomaly of the 320 cm<sup>-1</sup> phonon mode is explained as due to the drastic changes of the local electric fields acting on the in-plane oxygens as the Josephson currents set in the SC state . Overall, this close analogy implies that the T<sub>c</sub> suppression upon Pr-substitution is caused by a depletion and/or a localization of the mobile holes of the CuO<sub>2</sub> planes rather than by pair-breaking.
For comparison, we also investigated a YBa<sub>2</sub>Cu<sub>2.94</sub>Zn<sub>0.04</sub>O<sub>6.9</sub> crystal whose T<sub>c</sub> is suppressed by a similar amount to T<sub>c</sub>=76(3) K. Meanwhile, it is well established that the T<sub>c</sub>-suppression in Zn-substituted samples is caused by strong pair-breaking due to impurity scattering in the unitarity limit on the Zn impurities while the hole content of the CuO<sub>2</sub> plane is hardly affected . Figure 2 displays $`\sigma _{1c}`$ at different temperatures between room temperature and 10 K. It is evident that no sign of a pseudogap in the normal state c-axis conductivity occurs, despite that fact that the T<sub>c</sub> value is severely suppressed. The c-axis conductivity instead hardly changes in the normal state. If anything, it rather exhibits a very weak Drude-like behavior since $`\sigma _{1c}^{el}`$ increases slightly with decreasing T and towards low frequency. Note that a similar behavior has been observed in a Zn-free YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6.9</sub> crystal which was annealed under the same conditions being almost optimally doped with T<sub>c</sub>=92 K . Also the anomaly of the 320 cm<sup>-1</sup> phonon mode is much weaker for the Zn-substituted crystal than for the 20 % Pr-substituted crystal and there is no clear evidence for the additional low-T peak. For the Zn-substituted crystal a spectral gap forms only in the superconducting state below T<sub>c</sub>=76 K. Then the pair breaking effect due to the Zn-impurities is clearly evident. Firstly, the size of the spectral gap is significantly reduced to 2$`\mathrm{\Delta }_{SC}`$450 cm<sup>-1</sup>as compared to 2$`\mathrm{\Delta }_{SC}`$650 cm<sup>-1</sup> in pure optimally doped YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6.9</sub> . Secondly, the residual conductivity remains very large even at the lowest temperature of 10 K. This indicates that a large number of quasi-particles remain unpaired in the SC state due to the strong pair-breaking effect of the Zn-impurities. No such characterstic signature of the pair-breaking effect is observed for the Pr-substituted crystals.
Figure 3 shows the T-dependence of the FIR c-axis conductivity of the Pr<sub>0.3</sub>Y<sub>0.7</sub>Ba<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> crystal with T<sub>c</sub>=64(3) K which is more strongly Pr-substituted than the previous one. The normal state spectra are shown in Fig. 3(a). Once more a spectral pseudogap can be seen to develop in the normal state. The pseudogap now already starts to form around room temperature and its size clearly exceeds the measured spectral range, i.e. $`\omega _{NG}`$$`>`$700 cm<sup>-1</sup>. A corresponding increase of the onset temperature and of the size of the pseudogap with increasing underdoping of the CuO<sub>2</sub> planes has been previously observed for strongly deoxygenated Y-123 crystals . All the characteristic spectral features are very similar like in a deoxygenated Y-123 crystal with comparable T$`{}_{c}{}^{}`$60 K. On the other hand, the absolute value of the $`\sigma _{el}^{1c}`$ is significantly higher in the Pr-substituted crystal. This effect can be understood to be due to the presence of the fully oxygenated and thus metallic CuO chains. Evidence for metallic CuO chains has been obtained even in pure Pr-123 . A similar chain-related effect on $`\sigma _{el}^{1c}`$ has previously been observed for Ca-substituted Y<sub>1-x</sub>Ca<sub>x</sub>Ba<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> crystals for which the absolute value depends on the oxygen content and thus the metallicity of the CuO chains while the characteristic frequency- and T-dependence is determined mainly by the hole doping of the CuO<sub>2</sub> planes . Another difference as compared to the deoxygenated Y-123 crystals is the observation of two weak defect modes in the spectral range of 400 to 450 cm<sup>-1</sup>. These modes are present already at room temperature but become sharper and thus more pronounced at low T. These defect modes possibly originate from oxygen defects within the CuO chain layer other than the usual O(1) oxygen defects in Y-123. Alternatively, they might corresond to crystal field excitations of the magnetic Pr<sup>3+</sup> ions. Figure 3(b) shows the low temperature data in the SC state. It can be seen that the 320 cm<sup>-1</sup> phonon mode becomes strongly renormalized and that an additional peak develops around 420 cm<sup>-1</sup> again in close analogy to deoxygenated underdoped YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6.6</sub> with T$`{}_{c}{}^{}`$60 K . The additional peak is not quite as pronounced as for the Pr-free underdoped Y-123 crystals. In the context of the Josephson-plasmon superlattice model this difference can be explained as due to the larger quasiparticle conductivity of the Pr-substituted crystal which leads to a stronger damping of the transversal Josephson plasmon mode . As was noted above, this difference can be attributed to the metallic conductivity of the CuO chains and the subsequently higher normal state c-axis conductivity of the Pr-substituted crystals. Also due to the presence of the fully oxygenated metallic CuO chains, the interbilayer Josephson plasmon frequency, as deduced from the zero crossing of $`\epsilon _1`$ and/or from the low frequency slope of $`\epsilon _1`$, appears to be somewhat higher in the Pr-substituted crystals compared to deoxygenated Y-123 with similar T<sub>c</sub>. More details about the dependence of the inter- and intrabilayer Josephson plasmons on the presence of the metallic CuO chains, including fits with the model given in , will be presented in a forthcoming publication.
Figures 4 and 5 show the FIR c-axis conductivity of Pr<sub>0.4</sub>Y<sub>0.6</sub>Ba<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> with T<sub>c</sub>=48(4) K and Pr<sub>0.5</sub>Y<sub>0.5</sub>Ba<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> with T<sub>c</sub>=24(5) K, respectively. The absolute values $`\sigma _{1c}^{el}`$ once more are significantly higher for the fully oxygenated Pr-substituted crystals than for comparably underdoped deoxygenated Y-123. This circumstance allows us to see more clearly the characteristic features of the normal state pseudogap in such strongly underdoped 123-type samples. For the x=0.4 crystal it is evident that the size of the normal state pseudogap exceeds the measured spectral range by far, i.e., $`\omega _{NG}`$$`>`$$`>`$700 cm<sup>-1</sup>. This finding confirms our previous report that the size of the pseudogap of deoxygenated Y-123 crystals increases continuously on the underdoped side . The pseudogap forms around room temperature and it is not related to the formation of the additional low T peak nor to the anomaly of the 320 cm<sup>-1</sup> phonon mode. Notably, the characteristic shape of the anomaly of the 320 cm<sup>-1</sup> phonon mode and the additional peak, which both merge for the present sample, indicate that the size of the intrabilayer Josephson plasmon is somewhat smaller for this 40% Pr-substituted crystal than for a comparably underdoped pure YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6.5</sub> crystal with T<sub>c</sub>=52 K. Instead the spectral shape of the anomalous 320 phonon mode and the additional peak resemble that of a more strongly underdoped YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6.45</sub> crystal with T<sub>c</sub>=25 K . It worth noting that this effect is expected within the Josephson superlattice model since the separation of the CuO<sub>2</sub> planes is larger in the Pr-substituted crystal than in the pure Y-123 resulting in a smaller intrabilayer Josephson plasma frequency for the Pr-substituted sample. In contrast, as noted above, the interbilayer plasma frequency appears to be larger in the Pr-substituted crystals due to the presence of the metallic CuO chains. Figure 5 shows the FIR c-axis conductivity of our most heavily Pr-substituted Pr<sub>0.45</sub>Y<sub>0.55</sub>Ba<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> crystal with T<sub>c</sub>=24(5) K. It is remarkable that for this crystal, which is located close to the metal insulator transition around x=0.55, the normal state pseudogap suddenly is less pronounced and its size has decreased. This finding seems to indicate that the pseudogap disappears around the metal insulator transition. Note that it has not been possible to follow the evolution of the electronic c-axis conductivity in such detail in correspondingly underdoped deoxygenated Y-123 crystals for which the absolute values of $`\sigma _{1c}^{el}`$ are much lower .
Finally, we comment on the question whether our FIR c-axis conductivity data provide any information about the mechanism which causes the hole depletion of the CuO<sub>2</sub> planes upon Pr-substitution. Recent x-ray diffraction measurements revealed that flux-grown Pr-123 crystals are Cu deficient on the Cu(1) site and/or exhibit a partial substitution of Pr on the Ba-site . It was argued that these effects may be responsible for the hole depletion. The only anomalous features in the FIR phonon modes of our Y,Pr-123 crystals which may be indicative of structural disorder within the CuO chain layer are the weak defect mode at 420 cm$`^1,`$ whose origin is yet unknown, and a broadening of the phonon mode at 280 cm<sup>-1</sup> due to vibration of chain oxygen which is more pronounced for the Pr-substituted crystals than for deoxygenated crystals with similar T<sub>c</sub> . However, a similar defect mode around 400-450 cm<sup>-1</sup> and an even stronger broadening of the Cu(1) mode at 280 cm<sup>-1</sup> have recently been observed in Nd-123 which is superconducting with T<sub>c</sub>=93 K . Therefore, while these features may be related to a deficiency on the Cu(1) site, it is unlikely that they are responsible for the strong T<sub>c</sub>-suppression which occurs only for the Pr-substituted crystal but not for Nd substituted one. In our spectra we do not observe any anomalous broadening or shift for example of the phonon mode at 155 cm<sup>-1</sup> to which mainly Ba and Cu(1) contribute . While the masses of Ba and Pr are very similar one still expects this mode to be broadened if Ba<sup>2+</sup> is substituted by Pr<sup>3+</sup> since their electronic environment (oxygen configuration) should be rather different . To conclude, we do not observe any anomalous changes of the FIR active phonon modes which could be used as an indication that the decrease of the hole doping of the CuO<sub>2</sub> planes can be explained as a chemical doping effect. Our data rather favor the model of Fehrenbacher and Rice where the charge carriers of the CuO<sub>2</sub> planes are redistributed and localized in hybridized Pr 4f and O 2p orbitals . Note that this model is supported by recent x-ray absorption measurements .
In summary, by spectral ellipsometry we have studied the far-infrared c-axis conductivity of Pr-substituted Pr<sub>x</sub>Y<sub>1-x</sub>Ba<sub>2</sub>Cu<sub>2</sub>O<sub>7</sub> crystals with 0.2$``$x$``$0.5. We have shown that the c-axis conductivity of these crystals exhibits spectral features similar to deoxygenated and thus underdoped YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub>. A spectral pseudogap already forms in the normal state, the oxygen bond-bending phonon mode at 320 cm<sup>-1</sup> exhibits an anomalous T dependence and an additional broad absorption peak forms at low T. This finding suggests that the T<sub>c</sub> suppression in the Pr-substituted samples is caused by a decrease in the concentration or a localization of the mobile hole carriers of the CuO<sub>2</sub> planes.
We gratefully acknowledge G.P. Williams and L. Carr for technical help at the U4IR beamline at NSLS and E. Brücher and R.K. Kremer for performing the SQUID magnetisation measurements.
* Permanent address: Institute of Experimental Physics, Warsaw University, Hoża 69, 00-681 Warsaw, Poland
Figure Captions
Figure 1: Spectra of the real part of the FIR c-axis conductivity of Pr<sub>0.2</sub>Y<sub>0.8</sub>Ba<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> with T<sub>c</sub>=81(3) K in the normal- and the superconducting state. The position of the additional low temperature peak is indicated by the solid arrow.
Figure 2: Optical FIR c-axis conductivity of Zn-substituted YBa<sub>2</sub>Cu<sub>2.94</sub>Zn<sub>0.06</sub>O<sub>6.9</sub> with T<sub>c</sub>=76(3) K at different temperatures in the normal- and the superconducting state. The arrow indicates the onset of the spectral gap in the superconducting state.
Figure 3: Temperature dependence of the real part of the FIR c-axis conductivity of Pr<sub>0.3</sub>Y<sub>0.7</sub>Ba<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> with T<sub>c</sub>=64(3) K, (a) in the normal state and (b) in the superconducting state. The solid arrow marks the position of the additional low-temperature absorption peak.
Figure 4: Real part of the FIR c-axis conductivity of Pr<sub>0.4</sub>Y<sub>0.6</sub>Ba<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> with T<sub>c</sub>=48(4) K. The arrow indicates the position of the low temperature absorption peak which merges with the phonon mode at 320 cm<sup>-1</sup>.
Figure 5: Temperature dependence of the real part of the FIR c-axis conductivity of Pr<sub>0.5</sub>Y<sub>0.5</sub>Ba<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> with T<sub>c</sub>=24(5) K.
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# Quantum Zeno effect and the detection of gravitomagnetism.
## 1 Introduction.
In more than three–quarters of a century the theory of general relativity (GR) has achieved a great experimental triumph. Neverwithstanding, at this point it is also important to comment that all the current direct confirmations of GR are confirmations of weak field corrections to the Galilei–Newton mechanics . We must also add that one of the most important, and yet undetected, predictions of GR is the so called gravitomagnetic field , sometimes also called Lense–Thirring effect , which is generated by mass–energy currents. Its measurement would constitute a direct experimental evidence against an absolute inertial frame of reference, and would at the same time show the basic role that local inertial frames play in nature, i.e., it would be a direct proof that local inertial frames are influenced and dragged by mass–energy currents relative to other mass.
The first efforts in the detection of this gravitomagnetic field are quite old and have already included many interesting proposals .
An additional topic in connection with gravitomagnetism is related to its coupling with intrinsic spin, this issue is of fundamental interest since it comprises the inertial properties of intrinsic spin. It is noteworthy to comment that this point is under constant analysis .
In this work we introduce two experimental proposals that could lead to the detection of the coupling between intrinsic spin and the gravitomagnetic field. We analyze the role that the gravitomagnetic field of the Earth could have on a quantum system with spin $`1/2`$, i.e., our results could allow us to confront the effects of mass–energy density currents upon spin. In particular we deduce a Rabi formula, which depends on the coupling between the spin of the quantum system and the gravitomagnetic field of the Earth. Afterwards, the continuous measurement of the energy of the spin $`1/2`$ system is considered, and a Zeno effect is obtained.
## 2 Rabi transitions and the gravitomagnetic field.
Let us consider a spin $`1/2`$ system immersed in the gravitational field of a rotating uncharged, idealized spherical body with mass $`M`$ and angular momentum $`J`$. In the weak field and slow motion limit the metric, in the Boyer–Lindquist coordinates, reads
$`ds^2=c^2\left(1{\displaystyle \frac{2GM}{c^2r}}\right)dt^2+\left(1{\displaystyle \frac{2GM}{c^2r}}\right)^1dr^2`$
$`+r^2\left(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2\right){\displaystyle \frac{4GJ}{c^2r}}\mathrm{sin}^2\theta d\varphi dt.`$ (1)
The gravitomagnetic field in this case is approximately
$`\stackrel{}{B}=2{\displaystyle \frac{G}{c^2}}{\displaystyle \frac{\stackrel{}{J}3(\stackrel{}{J}\widehat{x})\widehat{x}}{|\stackrel{}{x}|^3}}.`$ (2)
We will assume that the expression that describes the precession of orbital angular momentum, immersed, for instance, in the gravitational field of the Earth, can be also used for the description of the dynamics in the case of intrinsic spin. This is a natural extension of general relativity .
Let us now denote the angular momentum of our spherical body by $`\stackrel{}{J}=J\widehat{z}`$, being $`\widehat{z}`$ the unit vector along the direction of the angular momentum. Our quantum particle is prepared such that $`\stackrel{}{S}=S_z\widehat{z}`$, it has vanishing small velocity and acceleration, and it is located on the $`z`$–axis, with coordinate $`Z`$.
There is a formal analogy between the weak field and slow motion of the gravitomagnetic field in general relativity and the magnetic field in electromagnetism . Following this analogy we may write down the interaction Hamiltonian (acting in the two–dimensional spin space of our spin $`1/2`$ system), which gives the coupling between $`\stackrel{}{B}`$ and the spin, $`\stackrel{}{S}`$, of our particle
$`H=\stackrel{}{S}\stackrel{}{B}.`$ (3)
Introducing expression (2) we may rewrite the interaction Hamiltonian as follows
$`H=2{\displaystyle \frac{GJ\mathrm{}}{c^2Z^3}}[|+><+||><|].`$ (4)
Here $`|+>`$ and $`|>`$ represent the eigenkets of $`S_z`$. Clearly, the introduction of the gravitomagnetic field renders two energy states
$`E_{(+)}=2{\displaystyle \frac{GJ\mathrm{}}{c^2Z^3}},`$ (5)
$`E_{()}=2{\displaystyle \frac{GJ\mathrm{}}{c^2Z^3}},`$ (6)
where $`E_{(+)}`$ ($`E_{()}`$) is the energy of the spin state $`+\mathrm{}/2`$ ($`\mathrm{}/2`$). Let us now define the frequency
$`\mathrm{\Omega }=\left(E_{(+)}E_{()}\right)/\mathrm{}=4{\displaystyle \frac{GJ}{c^2Z^3}}.`$ (7)
The present analogy allows us to consider the emergence of Rabi transitions . In order to do this let us now introduce a rotating magnetic field, which, at the point where the particle is located, has the following form
$`\stackrel{}{b}=b\left[\mathrm{cos}(wt)\widehat{x}+\mathrm{sin}(wt)\widehat{y}\right],`$ (8)
where $`\widehat{x}`$ and $`\widehat{y}`$ are two unit vectors perpendicular to the $`z`$–axis, and $`b`$ is a constant magnetic field.
Under these conditions the total Hamiltonian reads
$`H_T=2{\displaystyle \frac{GJ\mathrm{}}{c^2Z^3}}[|+><+||><|]`$
$`{\displaystyle \frac{eb\mathrm{}}{2mc}}[e^{iwt}|+><|+e^{iwt}|><+|].`$ (9)
Looking for a solution in the form $`|\alpha >=c_{(+)}(t)|+>+c_{()}(t)|>`$, we find the usual situation (our quantum system has been initially prepared such that $`c_{()}(0)=1`$ and $`c_{(+)}(0)=0`$.)
$`c_{()}(t)=\mathrm{exp}\left[i{\displaystyle \frac{E_{()}}{\mathrm{}}}t+{\displaystyle \frac{i}{2}}(w\mathrm{\Omega })t\right]\left[\mathrm{cos}(\mathrm{\Gamma }t)i{\displaystyle \frac{(w\mathrm{\Omega })}{2\mathrm{\Gamma }}}\mathrm{sin}(\mathrm{\Gamma }t)\right],`$ (10)
$`c_{(+)}(t)=i{\displaystyle \frac{eb}{2mc\mathrm{\Gamma }}}\mathrm{exp}\left[i{\displaystyle \frac{E_{(+)}}{\mathrm{}}}t{\displaystyle \frac{i}{2}}(w\mathrm{\Omega })t\right]\mathrm{sin}(\mathrm{\Gamma }t).`$ (11)
where $`\mathrm{\Gamma }=\sqrt{(\frac{eb}{2mc})^2+\frac{(w\mathrm{\Omega })^2}{4}}`$.
In this way we find
$`{\displaystyle \frac{|c_{()}(t)|^2}{|c_{()}(t)|^2+|c_{(+)}(t)|^2}}=\left[1+{\displaystyle \frac{(\frac{eb}{2mc\mathrm{\Gamma }})^2\mathrm{sin}^2(\mathrm{\Gamma }t)}{\mathrm{cos}^2(\mathrm{\Gamma }t)+\frac{(w\mathrm{\Omega })^2}{4\mathrm{\Gamma }^2}\mathrm{sin}^2(\mathrm{\Gamma }t)}}\right]^1.`$ (12)
Clearly, the Rabi transitions depend upon the coupling between spin and the gravitomagnetic field.
$`\left(4{\displaystyle \frac{GJ}{c^2Z^3}}w\right)^2=4\left[\mathrm{\Gamma }^2\left({\displaystyle \frac{eb}{2mc}}\right)^2\right].`$ (13)
## 3 Quantum Zeno effect and gravitomagnetism.
Let us now measure, continuously, the energy of our spin $`1/2`$ system, such that $`E`$ is the measurement output, and that this experiment lasts a time $`T`$. This kind of measuring process can be described by the so called effective Hamiltonian formalism , which is one of the models that exist in the topic of quantum measurement theory . In our case the corresponding effective Hamiltonian reads
$`H_{eff}=2{\displaystyle \frac{GJ\mathrm{}}{c^2Z^3}}[1+i{\displaystyle \frac{2\mathrm{}}{T\mathrm{\Delta }E^2}}(E{\displaystyle \frac{GJ\mathrm{}}{c^2Z^3}})]|+><+|`$
$`2{\displaystyle \frac{GJ\mathrm{}}{c^2Z^3}}[1+i{\displaystyle \frac{2\mathrm{}}{T\mathrm{\Delta }E^2}}(E+{\displaystyle \frac{GJ\mathrm{}}{c^2Z^3}})]|><|`$
$`{\displaystyle \frac{eb\mathrm{}}{2mc}}[e^{iwt}|+><|+e^{iwt}|><+|]i{\displaystyle \frac{E^2\mathrm{}}{T\mathrm{\Delta }E^2}}\mathrm{\Pi },`$ (14)
where $`\mathrm{\Pi }`$ is the unit operator in the spin space of our particle. Looking for solutions with the form $`|\alpha >=c_{(+)}(t)|+>+c_{()}(t)|>`$, we deduce
$`c_{()}(t)=\mathrm{exp}\left[i{\displaystyle \frac{E_{()}}{\mathrm{}}}t{\displaystyle \frac{(E_{()}E)^2}{T\mathrm{\Delta }E^2}}t+i\stackrel{~}{\mathrm{\Gamma }}t\right]`$
$`\times \left[c_{()}(0)\mathrm{cos}(\beta t)i{\displaystyle \frac{c_{()}(0)\stackrel{~}{\mathrm{\Gamma }}+(\gamma /\mathrm{})c_{(+)}(0)}{\beta }}\mathrm{sin}(\beta t)\right],`$ (15)
$`c_{(+)}(t)=\mathrm{exp}\left[i{\displaystyle \frac{E_{(+)}}{\mathrm{}}}t{\displaystyle \frac{(E_{(+)}E)^2}{T\mathrm{\Delta }E^2}}ti\stackrel{~}{\mathrm{\Gamma }}t\right]`$
$`\times \left[c_{(+)}(0)\mathrm{cos}(\beta t)+i{\displaystyle \frac{c_{(+)}(0)\stackrel{~}{\mathrm{\Gamma }}(\gamma /\mathrm{})c_{()}(0)}{\beta }}\mathrm{sin}(\beta t)\right],`$ (16)
where $`\stackrel{~}{\mathrm{\Gamma }}=\frac{(w\mathrm{\Omega })}{2}+\frac{i}{2T\mathrm{\Delta }E^2}\left[(E_{(+)}E)^2(E_{()}E)^2\right]`$, $`\beta ^2=(\gamma /\mathrm{})^2+\stackrel{~}{\mathrm{\Gamma }}^2`$, and finally $`\gamma =\frac{eb\mathrm{}}{2mc}`$.
Let us now suppose that the measurement output is the energy of the ground state, $`E_{()}`$, that we have a resonant perturbation, and that initially only the lowest energy state was populated, in other words, $`E=E_{()}`$, $`\mathrm{}w=E_{(+)}E_{()}`$, and $`c_{()}(0)=1`$, $`c_{(+)}(0)=0`$.
Hence (15) and (16) become
$`c_{()}(t)=\mathrm{exp}\left[i{\displaystyle \frac{E_{()}}{\mathrm{}}}t{\displaystyle \frac{(E_{(+)}E_{()})^2}{2T\mathrm{\Delta }E^2}}t\right]\left[\mathrm{cos}(\beta t)i{\displaystyle \frac{\stackrel{~}{\mathrm{\Gamma }}}{\beta }}\mathrm{sin}(\beta t)\right],`$ (17)
$`c_{(+)}(t)=i{\displaystyle \frac{\gamma }{\beta \mathrm{}}}\mathrm{exp}\left[i{\displaystyle \frac{E_{(+)}}{\mathrm{}}}t{\displaystyle \frac{(E_{(+)}E_{()})^2}{2T\mathrm{\Delta }E^2}}t\right]\mathrm{sin}(\beta t).`$ (18)
Let us now assume that $`\frac{(E_{(+)}E_{()})^4}{4T^2\mathrm{\Delta }E^4}>\gamma ^2/\mathrm{}^2`$, then
$`P_{()}(t)=\left[1+{\displaystyle \frac{\mathrm{sinh}^2(\frac{\gamma }{\mathrm{}}\stackrel{~}{\mathrm{\Omega }}t)}{\stackrel{~}{\mathrm{\Omega }}^2[\mathrm{cosh}(\frac{\gamma }{\mathrm{}}\stackrel{~}{\mathrm{\Omega }}t)+\frac{\mathrm{}(E_{(+)}E_{()})^2}{2T\gamma \stackrel{~}{\mathrm{\Omega }}\mathrm{\Delta }E^2}\mathrm{sinh}(\frac{\gamma }{\mathrm{}}\stackrel{~}{\mathrm{\Omega }}t)]^2}}\right]^1,`$ (19)
where $`\stackrel{~}{\mathrm{\Omega }}=\sqrt{\frac{\mathrm{}^2(E_{(+)}E_{()})^4}{4T^2\gamma ^2\mathrm{\Delta }E^4}1}`$, $`\gamma =\frac{eb\mathrm{}}{2mc}`$, and $`P_{()}(t)=\frac{|c_{()}(t)|^2}{|c_{()}(t)|^2+|c_{(+)}(t)|^2}`$.
In the case $`t\mathrm{}`$ this last expression reduces to
$`P_{()}^{(\mathrm{})}=\left[1+\left({\displaystyle \frac{c^2Z^3}{4GJ\mathrm{}}}\right)^2{\displaystyle \frac{ebT\mathrm{\Delta }E^2}{mc}}\left(\sqrt{1({\displaystyle \frac{c^2Z^3}{4GJ\mathrm{}}})^4({\displaystyle \frac{ebT\mathrm{\Delta }E^2}{mc}})^2}1\right)^2\right]^1.`$ (20)
Clearly, Rabi transitions are inhibited, and the asymptotic value that here appears depends explicitly upon the coupling between intrinsic spin and the gravitomagnetic field, i.e., $`J`$ emerges in expression (20).
At this point it must be commented that the behavior of spin leads, in some cases, to the emergence of a non–geometric element in gravity .
In this work Ahluwalia has considered two different classes of flavor–oscillation clocks. The first one comprises the superposition of different mass eigenstates, associated to a quantum test particle, such that all the terms of the corresponding superposition have the same spin component. The second class of flavor–oscillation clocks, contains, at least, two distinct spin projections.
If the gravitomagnetic field is absent, then both clocks redshift identically in the corresponding gravitational field. Nevertheless, if the source of the gravitational field has a nonvanishing angular momentum, then these redshifts do not coincide any more . This fact depends not only upon the gravitomagnetic component of the gravitational field, but also on the quantum mechanical features of the employed quantum test particle. In other words, here a non–geometric element appears when gravitational and quantum mechanical phenomena are considered simultaneously.
Clearly, in the present essay we have a quantum system with spin immersed in a nonvanishing gravitomagnetic field. Nevertheless, our case is an eigenstate of the spin operator $`S_z`$, something that in Ahluwalia’s second class of flavor–oscillation clocks does not happen. This last remark means that our quantum system is closer to Ahluwalia’s first class of flavor–oscillation clocks than to his second one.
Finally, we must add that it is now possible to test, experimentally, the quantum Zeno effect , particularly using Penning traps to analyze Rabi transitions .
Acknowledgments.
The author would like to thank A. A. Cuevas–Sosa his help, and D.-E. Liebscher for the fruitful discussions on the subject. It is also a pleasure to thank R. Onofrio for bringing Refs. 6 and 11 to my attention. The hospitality of the Astrophysikalisches Institut Potsdam is also kindly acknowledged. This work was supported by CONACYT Posdoctoral Grant No. 983023.
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# Abstract
## Abstract
This letter is a critique of Barbero’s constrained Hamiltonian formulation of General Relativity on which current work in Loop Quantum Gravity is based. While we do not dispute the correctness of Barbero’s formulation of general relativity, we offer some criticisms of an aesthetic nature. We point out that unlike Ashtekar’s complex $`SU(2)`$ connection, Barbero’s real $`SO(3)`$ connection does not admit an interpretation as a space–time gauge field. We show that if one tries to interpret Barbero’s real $`SO(3)`$ connection as a space–time gauge field, the theory is not diffeomorphism invariant. We conclude that Barbero’s formulation is not a gauge theory of gravity in the sense that Ashtekar’s Hamiltonian formulation is. The advantages of Barbero’s real connection formulation have been bought at the price of giving up the description of gravity as a gauge field.
## Introduction
In the eighties, Ashtekar introduced complex “new variables” on the phase space of General Relativity, which greatly simplified the form of the constraints. These variables are a (densitised) soldering form and a complex $`SU(2)`$ connection. Ashtekar’s motivation was to formulate General Relativity in a manner similar to Yang-Mills fields, so that ideas and techniques used in quantising gauge theories could be imported into Quantum Gravity. This motivation was an attractive one, since the “gauge” description of Nature is a unifying idea that seems to permeate diverse branches of physics. Ashtekar suggested the use of a gauge field (“the connection representation”) as the basic configuration variable in canonical gravity instead of the more traditional “metric representation”. This area has been an active one . Witten’s paper on the solution of $`2+1`$ gravity was partly inspired by Ashtekar’s suggestion. In 2+1 gravity the connection representation has proved extremely useful and yields a considerably simpler formulation of quantum gravity than the metrical description. In both $`2+1`$ and in $`3+1`$ dimensions, the connection variable is the pull-back of a space–time connection to a spatial slice $`𝒮`$. However, in 3+1 dimensions, progress has been hampered by the ‘reality conditions’, which are necessary because of the use of complex variables on the real phase space of General Relativity.
In 1994, Barbero wrote a very influential paper pointing out that a small modification of Ashtekar’s original canonical transformation leads to a new Hamiltonian formulation of General Relativity in which the basic variables are <sup>1</sup><sup>1</sup>1 These real variables were earlier considered by Ashtekar and discarded in favour of the complex (“new”) variables of because the latter simplified the constraints. a real $`SU(2)`$ connection <sup>2</sup><sup>2</sup>2 Or equivalently $`SO(3)`$; we will not make this relatively fine distinction in this paper. and a real densitised triad. The form of the Hamiltonian constraint in Barbero’s formulation is not as simple as in Ashtekar’s original formulation. But this is a small price to pay. There are significant advantages in using real variables on the phase space of General Relativity, because the ‘reality conditions’ which had to be imposed in Ashtekar’s original formulation, are no longer necessary. As a result, Barbero’s Hamiltonian formulation (BHF) has gained wide acceptance and is currently the basis of Loop Quantum Gravity. A lot of work has been done on the space of real $`SU(2)`$ connections on manifolds. Since $`SU(2)`$ is a compact group it has an invariant (Haar) measure, and one is able to achieve a high degree of mathematical control over the space of connections.
It was pointed out by Immirzi that Barbero’s canonical transformation could be slightly generalised: a one parameter family of canonical transformations were possible and all of them led to a Hamiltonian formulation based on a real $`SU(2)`$ connection. This free parameter $`\beta `$ is known as the “Immirzi parameter”, and does not appear to be fixed by any theoretical considerations.
There are several puzzling features about Barbero’s Hamiltonian formulation. Viewed as a gauge theory (and this was Ashtekar’s original motivation), the gauge group of General Relativity is certainly non-compact. Depending on the approach, one might either believe the gauge group to be the Lorentz group or the Poincare group . It is generally agreed that the gauge group of General Relativity must be non–compact. How then is it possible to formulate GR as gauge theory of a compact group, as Barbero seems to? Of course, it is possible to formulate GR as a gauge theory of a complex SU(2) group as in the original Ashtekar formulation. Over the complex numbers, there is no distinction between compact and non–compact gauge groups. The puzzle is that Barbero’s Hamiltonian formulation uses a gauge group which is both real and compact. Since BHF was derived by a canonical transformation from a diffeomorphism invariant theory, it must of course be diffeomorphism invariant. However, there does not presently exist any manifestly covariant Lagrangian formulation of General Relativity as a gauge theory in which the gauge group is real and compact. How does one understand this?
Holst has given a Lagrangian formulation equivalent to General Relativity and shown that on Legendre transformation, it results in BHF. A curious feature of Holst’s derivation is that he starts with a gauge theoretic formulation with a non–compact $`SO(3,1)`$ gauge group. Yet he is able to make contact with Barbero’s formulation, which is based on an $`SO(3)`$ gauge group. It is not clear a priori how the reduction in the gauge group takes place.
Another puzzling feature of Barbero’s Hamiltonian formulation is that there are many of them. For every nonzero, real value of the Immirzi parameter $`\beta `$ there is a Hamiltonian formulation of General Relativity with as much claim to validity as Barbero’s original formulation. It is not easy to understand the origin of the ‘Immirzi ambiguity’. It first appears as a parameter in a canonical transformation that one performs on the phase space of GR. However, instead of dissappearing from the final results of the theory (as one might have hoped), it appears in the spectrum of operators and also in the final expression for Black Hole Entropy as calculated using Loop Quantum Gravity . If the Immirzi parameter is indeed present in Quantum Gravity, it would seem that a new fundamental constant (which can only be fixed by experiment) has entered into physics. Such quantization ambiguities do occur in other areas of physics (e.g $`\theta `$ vacua), but they are well understood to arise from a multiply connected configuration space. The Immirzi ambiguity does not appear to be of topological origin and does require understanding.
We emphasize that we do not question the correctness of BHF as a Hamiltonian formulation or its equivalence of BHF to General Relativity. This equivalence is assured since BHF was derived from the ADM formulation of GR by a canonical transformation. The aim of this paper is to understand the basis of BHF as a reformulation of General Relativity. Is BHF a gauge theory of gravity? How does it happen that the gauge group is both real and compact? How does the diffeomorphism invariance of BHF jive with the absence of any manifestly covariant Lagrangian for GR as a real $`SO(3)`$ gauge theory. In this paper we will answer these questions and thereby clarify some aspects of Barbero’s formulation.
Barbero’s Hamiltonian formulation can be derived by making a canonical transformation starting from a constrained Hamiltonian Formulation (CHF) due to Ashtekar- the extended phase space construction (EPS). We note in passing that the EPS was an intermediate step in Ashtekar’s original derivation of the new variables. The EPS consists of the following ingredients: the basic variables are $`(\stackrel{~}{E}_i^a,K_a^i)`$, which form a canonically conjugate pair. The space–time meaning of these variables is that $`\stackrel{~}{E}_i^a`$ is a densitised triad on a spatial slice $`𝒮`$ and $`K_a^i`$ is the extrinsic curvature tensor of $`𝒮`$ with one index converted into a triad index. The three metric tensor $`q_{ab}`$ of $`𝒮`$ is a derived object, which can be constructed from $`\stackrel{~}{E}_i^a`$ as explained in . Similarly, the extrinsic curvature of $`𝒮`$ can also be expressed in terms of the basic fields. The constraints of the theory are
$`ϵ_{ijk}K_a^j\stackrel{~}{E}^{ak}`$ $``$ $`0`$
$`D_a[\stackrel{~}{E}_k^aK_b^k\delta _b^a\stackrel{~}{E}_k^cK_c^k]`$ $``$ $`0`$
$`\sqrt{q}R+{\displaystyle \frac{2}{\sqrt{q}}}\stackrel{~}{E}_i^{[a}\stackrel{~}{E}_j^{b]}K_a^iK_b^j`$ $``$ $`0,`$
where $`D_a`$ is the covariant derivative associated with $`q_{ab}`$ and $`R`$, its scalar curvature. The Hamiltonian is a combination of constraints. (We assume throughout this paper that space is closed so that we can drop spatial boundary terms.) The EPS formulation is strongly diffeomorphism invariant (mod $`SO(3)`$ gauge) with the above space–time interpretation for the basic variables $`(\stackrel{~}{E}_i^a,K_a^i)`$. We recall that a Hamiltonian theory is strongly diffeomorphism invariant (SDI) if a) there are constraints whose brackets reflect the Lie algebra of the diffeomorphism group. b) the basic variables transform as is expected of them from their space–time interpretation. For this criterion to be applied, one must first declare the space–time interpretation of the basic variables in the theory.
BHF differs from the EPS only by a canonical transformation. The basic variables in BHF are $`(\stackrel{~}{E}_i^a,A_a^i)`$, where $`\stackrel{~}{E}_i^a`$ is the same as before and
$$A_a^i:=\mathrm{\Gamma }_a^i+\beta K_a^i,$$
(1)
where $`\beta `$ is the Immirzi parameter and $`\mathrm{\Gamma }_a^i`$ are the triad spin coefficients. The constraints of the theory are now
$`𝒟_a\stackrel{~}{E}_i^a`$ $``$ $`0`$ (2)
$`\stackrel{~}{E}_i^bF_{ab}^i`$ $``$ $`0`$ (3)
$`ϵ^{ijk}\stackrel{~}{E}_i^a\stackrel{~}{E}_j^bF_{abk}2{\displaystyle \frac{(1+\beta ^2)}{\beta ^2}}\stackrel{~}{E}_{[i}^a\stackrel{~}{E}_{j]}^b(A_a^i\mathrm{\Gamma }_a^i)(A_b^j\mathrm{\Gamma }_b^j)`$ $``$ $`0,`$ (4)
where $`𝒟`$ is the covariant derivative associated with the Barbero connection (1). Note that the Scalar constraint (4) explicitly contains the Immirzi parameter. If one now endows $`A_a^i`$ with a space–time interpretation coming from (1) as a linear combination of the triad spin coefficients of $`𝒮`$ and the extrinsic curvature of $`𝒮`$, it is clear that the new formulation is still SDI (mod $`SO(3)`$ gauge), since the starting point was.
The original Ashtekar variables were derived by making the canonical transformation (1) with $`\beta =i`$ (or equivalently $`i`$). This choice is distinguished in several ways.
Why $`\beta =i`$ is special:
1. Simplification of the Hamiltonian constraint: The expression (4) for the Hamiltonian constraint in BHF contains a messy term multiplied by $`(1+\beta ^2)`$. For $`\beta =\pm i`$ this term disappears and we get a simplified form for the Hamiltonian constraint.
2. Full local Lorentz invariance: As emphasized by Immirzi the time gauge is not needed to derive the original Ashtekar variables. They can be derived from a manifestly covariant action maintaining full local Lorentz invariance. It follows that the theory is manifestly SDI, (i.e., SDI and not just modulo gauge).
3. Space–time interpretation for the connection: The Ashtekar connection (as defined by (1) with $`\beta =i`$) can also be interpreted as a space–time connection. More precisely, the Ashtekar connection is the pull back to $`𝒮`$ of a space–time connection 1-form. This is evident because Ashtekar’s formulation can be derived from a manifestly covariant Lagrangian in which one of the basic fields is an $`SL(2,IC)`$ space–time connection. The Ashtekar connection then appears as the pull-back of the space–time $`SL(2,IC)`$ connection to $`𝒮`$. Thus we would regard Ashtekar’s Hamiltonian Formulation (AHF) as a gauge theoretic reformulation of General Relativity.
For values of $`\beta `$ other than $`\pm i`$ these three properties do not obtain.
1. As Barbero points out, the form of the Hamiltonian constraint is more complicated for $`\beta `$ real, but one can choose to accept it.
2. It appears that one must choose the “time gauge” in order to arrive at Barbero’s formulation. In this formulation the full local Lorentz invariance of the Ashtekar formulation is lost. One can also choose to live with this, since this is only a choice of gauge.
3. The Barbero connection does not have a space–time interpretation as the pull back of a space–time connection 1-form to $`𝒮`$. In the case $`\beta =i`$, one had the option of interpreting the $`A_a^i`$ defined by (1) as a space–time connection. If one attempts to do this for real $`\beta `$, one finds that the theory violates Strong Diffeomorphism Invariance. Under diffeomorphisms tangential to $`𝒮`$, the $`A_a^i`$ defined by (1) does transform like a connection 1-form. However, for diffeomorphisms that are normal to $`𝒮`$, the canonical transformation properties of the $`A_a^i`$ do not reflect its proposed space–time interpretation as the pullback of a space–time connection 1–form. This point is explained and proved below.
Claim: Barbero’s connection cannot be interpreted as a space–time connection
Proof: Consider a solution $`(,g)`$ of Einstein’s equations and a loop $`\gamma `$ in $``$. If a connection transforms correctly under space–time diffeomorphisms (i.e, as a 1-form), the trace of the Holonomy of the connection along $`\gamma `$ should depend only on the loop $`\gamma `$ and should not depend on the slicing of $``$. It is easy to construct an example of an empty space solution to Einstein’s equations in which the Holonomy of the Barbero Connection changes from being trivial for one slice containing $`\gamma `$ to non –trivial for another slice containing $`\gamma `$. Consider flat space–time $`(,\eta )`$ with standard Minkowski coordinates $`(t,x,y,z)`$ and a loop $`\gamma `$ which is a circle of radius $`R`$ described by $`t=\sqrt{1+R^2},z=0,x=R\mathrm{cos}\theta ,y=R\mathrm{sin}\theta ,0\theta 2\pi `$. The flat slice $`𝒮_1`$ defined by $`t=\sqrt{1+R^2}`$ contains $`\gamma `$. Since both the intrinsic and extrinsic curvature of $`𝒮_1`$ vanish, the holonomy of the Barbero connection is the identity and its trace is $`3`$. However, the same loop is also contained in the hyperbolic slice, $`𝒮_2`$ defined by $`t^2x^2y^2=1`$. We now compute the trace of the Holonomy $`trH(A)`$ of the Barbero connection along the same loop $`\gamma `$.
At the point $`y^\mu =(t=1,\stackrel{}{x}=0)𝒮_2`$ the standard orthonormal frame $`\widehat{e}^I\mu =\delta ^I_\mu `$ ($`I=0,1,2,3`$ is a frame index and $`\mu =0,1,2,3`$ is a space-time index) has the property that $`\widehat{e}^0_\mu `$ is normal to $`𝒮_2`$. Let us move this frame to other points $`x^\mu 𝒮_2`$ by the Lorentz transformation
$$e^I{}_{\mu }{}^{}(x):=\mathrm{\Lambda }_\mu {}_{}{}^{\nu }(x)\widehat{e}^I{}_{\nu }{}^{},$$
(5)
where
$$\mathrm{\Lambda }_\mu {}_{}{}^{\nu }:=\delta _\mu {}_{}{}^{\nu }+(1x.y)^1(x_\mu +y_\mu )(x^\nu +y^\nu )2x_\mu y^\nu .$$
(6)
$`\mathrm{\Lambda }_\mu ^\nu `$ has the property that $`x_\mu =\mathrm{\Lambda }_\mu {}_{}{}^{\nu }y_{\nu }^{}`$, which ensures that $`e^0_\mu `$ is normal to $`𝒮_2`$ everywhere. (Our frame satisfies the ‘time gauge’.) The space-time connection 1-form is given by $`A_\alpha ^{IJ}`$
$$A_\alpha {}_{}{}^{IJ}:=e_\mu ^I_\alpha e^{\mu J}$$
(7)
($`A^{IJ}`$ with two internal indices is the space-time connection and should not be confused with the Barbero connection $`A^i`$, which has one internal index.) The contraction $`A^{IJ}:=t^\alpha A_\alpha `$ of the connection one-form with the tangent vector $`t^\alpha `$ to the curve $`\gamma `$ is easily worked out. The non-vanishing components are $`A^{01}=A^{10}=y,A^{02}=A^{20}=x,A^{12}=A^{21}=1t`$ Constructing Barbero’s connection by the formula
$$A^i=1/2ϵ^{ijk}A^{jk}+\beta A^{0i},$$
(8)
we find that $`A^1=\beta y,A^2=\beta x,A^3=1t`$. An elementary calculation then yields the result
$$\mathrm{Tr}H(A)=1+2\mathrm{cos}(2\pi \sqrt{1+R^2(1+\beta ^2)})$$
(9)
Notice that the trace of the Holonomy of the Barbero connection along the same loop $`\gamma `$, depends on $`\beta `$ and (except for the special values $`\beta =\pm i`$) is not the equal to $`3`$. Briefly, the trace of the holonomy of the Barbero Connection along a loop $`\gamma `$ is not just a property of the loop but also depends on the slicing.
What this means is that it is impossible to attach a gauge theoretic space–time interpretation to Barbero’s connection. As a spatial connection, Barbero’s connection is certainly well defined and, in fact, transforms correctly under diffeomorphisms that preserve the spatial slice. But unlike Ashtekar’s connection, Barbero’s connection does not admit an interpretation as a space–time gauge field. Such an interpretation is not consistent with the Poisson bracket relations between $`A_a^i`$ and the Scalar constraint. It appears that the “extra” term proportional to $`1+\beta ^2`$ ruins this bracket. Thus BHF is not a gauge theory of gravitation.
This observation removes most of the puzzles we had raised at the beginning of this letter. BHF is not a gauge theory of gravitation and so one cannot identify Barbero’s $`SO(3)`$ with the gauge group of gravity. There is no conflict between Barbero’s Hamiltonian formulation and our expectation that the gauge group of gravity must be noncompact.
This observation also removes the puzzle of how Holst is able to derive a compact gauge group from a non compact one. The variable defined as a connection variable by Holst is not the pullback of a space–time connection. Rather, certain components of the space–time connection are defined to be components of a new $`SO(3)`$ connection. This is Barbero’s connection and it has no space–time significance. One sometimes sees it implied that the reduction in the gauge group from $`SO(3,1)`$ to $`SO(3)`$ takes place because of our choice of the “time gauge”. It is indeed true that once we make this gauge choice<sup>3</sup><sup>3</sup>3 Holst’s orginal derivation of BHF from his covariant Lagrangian contained a small logical gap: he used the time gauge fixed action to perform the Legendre transformation. This gap has since been filled in papers by S. Alexandrov and Nuno Barros e Sa , our freedom to make additional gauge transformations is curtailed from $`SO(3,1)`$ to $`SO(3)`$. However this does not mean that the gauge group has been reduced. The pullback of the connection to a spatial slice is still an $`SO(3,1)`$ connection, in spite of our gauge choice.
While this point is elementary, it is worth making in detail and we digress briefly to do so. The central point here is a geometrical notion: the holonomy group of a connection . Given a $`G`$ connection $`A`$ on a manifold $`𝒮`$, we define its holonomy group to be the collection $`H_p(A)=\{H_{\gamma _p}(A)\}`$ of elements of $`G`$ which arise as holonomies of $`A`$ along closed loops $`\gamma _p`$ based at $`p`$. For example, if the connection is pure gauge, the holonomy group is trivial. It is easily seen that the holonomy groups based at $`p`$ and $`p^{}𝒮`$ are related by conjugation and therefore isomorphic. Further, under a local gauge transformation, the holonomy group of $`A`$ transforms by conjugation
$$H_p(A)g(p)H_p(A)g^1(p).$$
(10)
where $`g(p)G`$. We say that a connection $`A`$ is reducible if its holonomy group is a proper subgroup of $`G`$. Clearly, the reducibility of a connection is a gauge invariant notion and is independent of gauge choice. A general $`SO(3,1)`$ connection is not reducible to $`SO(3)`$.
One may of course, simply abandon the “gauge interpretation” of the theory and view the Barbero connection as a purely spatial connection which one uses for technical reasons to produce quantum states which are functionals of the extrinsic curvature. But then one must be aware that one has given up the gauge interpretation. We argue here that there are strong aesthetic reasons for retaining the gauge interpretation. One of the principal motivations of the Ashtekar program was the gauge description of gravity. This appears to be a unifying thread across the different forces of nature: they are all described by a space-time gauge field. The view we would like to offer here is that this was an important motivation of the original Ashtekar program and should be retained. From this point of view, the Immirzi parameter is not a free parameter but must be fixed to the special value $`i`$. This is one possible resolution of the “Immirzi Ambiguity”. Similar views have been expressed (see footnote 3 of ) by Alexandrov .
If one gives up the gauge interpretation of gravity, the Immirzi parameter appears not to be fixed by theory. This would not be a problem if the parameter dissappeared <sup>4</sup><sup>4</sup>4as in Alexandrov’s path integral quantisation . from all physical predictions of the theory. However, this is not the case: the Immirzi parameter does appear in the calculated value of Black Hole entropy in Loop Quantum Gravity. This phenomenon does not appear to be well understood (unlike the $`\theta `$ vacua of QCD). Rovelli and Thiemann offer a finite dimensional example of the Immirzi ambiguity. We do not find this example convincing: their original system does not suffer from any ambiguity. The ambiguity is introduced ‘by hand’ by changing the original configuration space in a $`\beta `$ dependent manner. Thus their “quantisation prodedure” does not quantise the original system at all, but quantises a one parameter family of distinct systems .
If one accepts the idea that canonical transformations made by the theorist can introduce parameters into the physical predictions of the quantum theory, one does lose predictive power. Admittedly, this loss is very small in the case of the Immirzi ambiguity. The ambiguity is to the extent of a single parameter which can be fixed by comparing a single prediction of the theory with an independent calculation (say the Black Hole entropy) or (thinking wishfully) an experiment. However, we would argue that in a field theory like General Relativity, there are infinitely many degrees of freedom (two at each spatial point) and one could, in principle, contemplate making separate canonical transformations in each of them, thereby introducing an enormous ambiguity into the theory. In the absence of any internal criterion to curb this ambiguity, the physical predictions of the theory may depend on an infinite number of parameters. Such a theory would have no predictive power.
If one wishes to describe Euclidean Gravity rather than Lorentzian Gravity, the gauge group is automatically compact and the Ashtekar variables are real. Barbero’s connection (with $`\beta =1`$) agrees with Ashtekar’s and one does have a real formulation of Euclidean gravity as a gauge theory. However, for $`\beta `$ not equal to unity, the same problem arises: Barbero’s connection is not a space–time connection and if one attempts to interpret it as such, one violates diffeomorphism invariance. It appears then that the values $`\beta =1`$ (for Euclidean) and $`\beta =i`$ (for Lorentzian Gravity) are very special. One could hope to relate these two by a Wick rotation as discussed in
In this paper we argue strongly for maintaining the gauge aspect of gravity in the approach to quantum gravity, even though the gauge group is non–compact and therefore not as tractable as say, $`SU(2)`$. It does not appear to us a strong argument to say that we study compact gauge groups because we do not know how to deal with non–compact gauge groups with mathematical rigour. The non–compactness of the gauge group appears to us an essentially physical feature of General Relativity, which is closely related to the Minkowskian signature of the space–time metric and light cones. Phenomena like infinite red shift, seen in Black Hole physics need non–compact groups for a gauge theoretical description. We would suggest that one must learn to deal with non–compact gauge groups. One could return either to the original Ashtekar variables with a complex connection or to the real Palatini tetrad formulation, with non–compact gauge group $`SO(3,1)`$. In our opinion these two problems – non–compact gauge groups and complex gauge groups- are the same problem in different guises. Evading both problems simultaneously is impossible, if one is interested in dealing with the physical effects of Lorentzian General Relativity as a gauge theory.
Acknowledgement: It is a pleasure to thank Richard Epp, B.R. Iyer, Sukanya Sinha and Madhavan Varadarajan for extended discussions and Ramesh Anishetty, Ghanashyam Date, N.D. Haridass, T.R. Govindrajan, Romesh Kaul and H.S Sharatchandra for their critical comments on this work.
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# 1 Introduction
## 1 Introduction
In the first order formalism of four-dimensional gravitation theory , the independent dynamical variables are the vierbein 1-form $`E`$ (giving the metric $`G`$) and the connection 1-form $`A`$. The vierbein dependence of the connection is given by the field equation with respect to $`A`$, whereas Einstein equation results from the field equation with respect to $`E`$. The action will be written, following , in a “topological” form, i.e. in such a way that it can be interpreted as an action of the 1-form fields $`E`$ and $`A`$ on a differentiable manifold $`M`$, without reference to any a-priori background metric. The latter point is known to be an essential characteristic of topological theories, and trying to exploit this feature belongs to the spirit of the modern attempts towards a construction of quantum gravity (see for reviews and further references).
Since the theory possesses two local symmetries – the diffeomorphism and local Lorentz invariances – one has to perform a gauge fixing for both. We will choose a gauge fixing of the Landau type, within the BRST formalism . Much in the same way as in topological theories, this requires the introduction of a nondynamical, background metric $`g`$. This construction closely parallels the one performed for the Chern-Simons theory in . It should be clear that the background metric, being introduced only in the gauge fixing part of the theory, should not affect in any way the physical outcome, as it has been shown for instance in for the – perturbative – quantum version of the Chern-Simons theory.
An interesting features of topological theories such as Chern-Simons or $`BF`$ theory, is the presence of a “vector supersymmetry” – a supersymmetry whose generator is a vector valued operator . In case the manifold admits isometries generated by Killing vectors – e.g. the space-time translations if the background metric is flat – the vector supersymmetry is a symmetry of the gauge-fixed action. It happens that its generator together with the BRST symmetry generator form an algebra which closes on the generators of the isometries of the translation type . The vector supersymmetry has been shown to play a key role in the ultraviolet finiteness of the topological theories .
Another interesting features – actually shared by any gauge theory, provided its gauge fixing be of the Landau type – is the so-called ghost equation , which restricts the coupling of the ghosts and implies the nonrenormalization of their field amplitude<sup>3</sup><sup>3</sup>3A review of the properties of topological theories mentionned above may be found in Chapters 6 and 7 of the book ..
The purpose of the present note is to show the existence of such a vector supersymmetry for Einstein gravity in the Palatini formalism. We shall in fact see that the vector supersymmetry is a direct consequence of the field equation of the Feynman-Dewitt-Faddeev-Popov ghost associated to diffeomorphism invariance. On the other hand, the ghost equation related to local Lorentz invariance will be seen to be algebraically associated with rigid Lorentz invariance. We shall also see that this supersymmetry, like in topological theories, yields the Sorella operator $`\delta `$ used in order to solve the BRST cohomology and to construct the invariants of the theory. The operator $`\delta `$ has been given for gravity in .
To the contrary of the topological theories of a Yang-Mills connection (Chern-Simons or BF), where the supersymmetry generators are the components of a vector and where the superalgebra closes on the translations, Einstein gravity in the Palatini formalism studied in the present paper will be seen to admit a supersymmetry possessing generators which are components of one vector and one antisymmetric tensor, the full algebra containing now all the ten Poincaré generators – in the case of a flat background metric at least<sup>4</sup><sup>4</sup>4Such a tensor supersymmetry has been pointed out in , in the case of a Chern Simons model in a gravitational background in the vielbein formalism..
Although the present work will only be concerned with the classical aspects of the theory, the results are of interest since, as we already said, they reveal the link between the construction of the observables via the $`\delta `$ operator of Sorella , on one hand, and the gauge fixing, through the ghost equation, on the other hand.
## 2 Symmetries, Gauge Fixing and BRST Invariance
The Einstein gravity Lagrangian in the first order formalism of Palatini may be written as :
$$S_{\mathrm{inv}}=\frac{1}{4}_M\epsilon _{IJKL}E^IE^JF^{KL}(A)+S_{\mathrm{matter}}(E,A,\mathrm{\Phi }).$$
(2.1)
The integral is taken over some differentiable 4-manifold $`M`$, $`E^I`$ is the vierbein 1-form, with $`I=0,\mathrm{},3`$ a tangent plane Lorentz index. $`F^{KL}`$ is the curvature 2-form
$$F^{IJ}(A)=dA^{IJ}+A^{IK}A_K^J$$
(2.2)
of the Lorentz connection<sup>5</sup><sup>5</sup>5If the connection is self-dual, (2.1) is the Ashtekar action . $`A^{IJ}`$, the latter being taken as an independant variable<sup>6</sup><sup>6</sup>6In a particular coordinate frame with $`x=`$ $`(x^\mu ,\mu =0,\mathrm{},3)`$, $`E^I=E_\mu ^Idx^\mu `$, $`F^{IJ}={\displaystyle \frac{1}{2}}F_{\mu \nu }^{IJ}dx^\mu dx^\nu `$, etc.. $`\epsilon _{IJKL}`$ is the rank four totally antisymmetric tensor, normalized by $`\epsilon _{0123}=1`$. In the following, the exterior multiplication symbol $``$ will be omitted. $`S_{\mathrm{matter}}`$ is some action for minimally coupled matter fields $`\mathrm{\Phi }`$, which we don’t need to specify. We shall in fact omit this part in the following, for the sake of simplicity.
The field equations given by the variations of this action read
$$\begin{array}{c}\frac{\delta S_{\mathrm{inv}}}{\delta E^I}=\frac{1}{2}\epsilon _{IJKL}E^JF^{KL},\hfill \\ \frac{\delta S_{\mathrm{inv}}}{\delta A^{IJ}}=\epsilon _{IJKL}E^KDE^L,\hfill \end{array}$$
(2.3)
where $`D`$ is the covariant exterior derivative: $`DE^I=dE^IA^I{}_{J}{}^{}E_{}^{J}`$. It is known that they lead to the usual specification of a torsion-free connection function of the vierbein<sup>7</sup><sup>7</sup>7This is true for pure gravity. In case of coupling with matter, the second field equation does not automatically lead to a vanishing torsion . One may then choose to stay with a non-vanishing action, or to impose a supplementary condition. and to the Einstein equation, in a Riemanian space-time with metric
$$G_{\mu \nu }=E_\mu ^IE_\nu ^J\eta _{IJ},$$
(2.4)
where $`\eta _{IJ}`$ is the Minkowsky metric<sup>8</sup><sup>8</sup>8We consider a Lorentzian signature. But everything applies as well to the Euclidean case. $`\mathrm{diag}(1,1,1,1,)`$, used to lower and rise the tangent space indices $`I,J,\mathrm{}`$.
The action (2.1) is invariant under the diffeomorphisms, written in infinitesimal form, the infinitesimal parameter being a vector field $`\xi `$:
$$\delta _{(\xi )}\phi =_\xi \phi ,\phi =E^I,A^{IJ},$$
(2.5)
where $`_\xi `$ is the Lie derivative along the vector $`\xi `$. It is also invariant under the local Lorentz transformations – written in infinitesimal form, with local parameters $`\omega ^{IJ}=\omega ^{JI}`$:
$$\begin{array}{c}\delta _{(\omega )}E^I=\omega ^I{}_{J}{}^{}E_{}^{J},\hfill \\ \delta _{(\omega )}A^{IJ}=d\omega ^{IJ}+\omega ^I{}_{K}{}^{}A_{}^{KJ}+\omega ^J{}_{K}{}^{}A_{}^{IK}.\hfill \end{array}$$
(2.6)
In view of the gauge fixing procedure it is convenient to express these local invariances in the form of a nilpotent BRST operation $`s`$ defined by :
$$\begin{array}{c}sE^I=_\xi E^I+\omega ^I{}_{J}{}^{}E_{}^{J},\hfill \\ sA^{IJ}=_\xi A^{IJ}+d\omega ^{IJ}+\omega ^I{}_{K}{}^{}A_{}^{KJ}+\omega ^J{}_{K}{}^{}A_{}^{IK},\hfill \\ s\xi =\frac{1}{2}\{\xi ,\xi \},(\text{or: }s\xi ^\mu =\xi ^\lambda _\lambda \xi ^\mu ),\hfill \\ s\omega ^I{}_{J}{}^{}=_\xi \omega ^I{}_{J}{}^{}+\omega ^I{}_{K}{}^{}\omega _{}^{K}{}_{J}{}^{},\hfill \end{array}$$
(2.7)
with $`s^2=0`$. The infinitesimal parameters $`\xi ^\mu (x)`$ \- the components of the vector $`\xi `$ – and $`\omega ^I{}_{J}{}^{}(x)`$ are now Grassmann (i.e. anticommuting) number fields – the Faddeev-Popov ghosts. The bracket $`\{,\}`$ is the Lie bracket<sup>9</sup><sup>9</sup>9In a particular coordinate frame, the Lie bracket of 2 vectors $`u`$, $`v`$ takes the form
$$\{u,v\}^\mu =u^\lambda _\lambda v^\mu \pm v^\lambda _\lambda u^\mu ,$$
with the sign $`+`$ is both $`u`$ e $`v`$ are odd, and the sign $``$ otherwise. Even (odd) refers to the commuting (anticommuting) character of the object..
In order to gauge fix the theory with respect to its local symmetries – diffeomorphism and local Lorentz invariances – we introduce antighosts<sup>10</sup><sup>10</sup>10Despite of an unfortunate but usual terminology, the antighosts are independent of the ghosts . $`\overline{\xi }_I`$, $`\overline{\omega }_{IJ}`$ and Lagrange multipliers $`\lambda _I`$, $`b_{IJ}`$, with the following nilpotent BRST transformations:
$$s\overline{\xi }_I=\lambda _I,s\lambda _I=0,s\overline{\omega }_{IJ}=b_{IJ},sb_{IJ}=0.$$
(2.8)
The gauge fixing part of the action is then defined as:
$$\begin{array}{c}S_{gf}=s_Md^4x\sqrt{g}g^{\mu \nu }\left(_\mu \overline{\xi }_IE_\nu ^I+\frac{1}{2}_\mu \overline{\omega }_{IJ}A_\nu ^{IJ}\right)\hfill \\ =_Md^4x\sqrt{g}g^{\mu \nu }\left(_\mu \lambda _IE_\nu ^I+\frac{1}{2}_\mu b_{IJ}A_\nu ^{IJ}\right)\hfill \\ +_Md^4x\sqrt{g}g^{\mu \nu }\left(_\mu \overline{\xi }_IsE_\nu ^I+\frac{1}{2}_\mu \overline{\omega }_{IJ}sA_\nu ^{IJ}\right),\hfill \end{array}$$
(2.9)
which is automatically BRST invariant. Note that in order to contract the world indices $`\mu `$, $`\nu `$ we have introduced a (BRST-invariant) background metric $`g_{\mu \nu }`$ – not to be confounded with the physical, dynamical metric $`G_{\mu \nu }`$ defined in (2.4).
This particular gauge fixing, which is of the Landau type, is completely determined by BRST invariance and by the “gauge conditions” – i.e. by the field equations for the Lagrange multipliers:
$$\frac{\delta S}{\delta \lambda _I}=_\mu \left(\sqrt{g}g^{\mu \nu }E_\nu ^I\right),\frac{\delta S}{\delta b_{IJ}}=_\mu \left(\sqrt{g}g^{\mu \nu }A_\nu ^{IJ}\right),$$
(2.10)
where $`S`$ is the total action
$$S=S_{\mathrm{inv}}+S_{\mathrm{gf}},$$
(2.11)
which is BRST invariant by construction: $`sS=0`$.
On the other hand, the field equations for the ghosts $`\xi `$ and $`\omega `$ are
$$\begin{array}{c}\frac{\delta S}{\delta \xi ^\mu }=\left(\sqrt{g}g^{\nu \lambda }_\lambda \overline{\xi }_I_\mu E_\nu ^I+_\nu \left(\sqrt{g}g^{\nu \lambda }_\lambda \overline{\xi }_IE_\mu ^I\right)\right)\hfill \\ \frac{1}{2}\sqrt{g}g^{\nu \lambda }_\lambda \overline{\omega }_{IJ}_\mu A_\nu ^{IJ}+\frac{1}{2}_\nu \left(\sqrt{g}g^{\nu \lambda }_\lambda \overline{\omega }_{IJ}A_\mu ^{IJ}\right),\hfill \\ \frac{\delta S}{\delta \omega ^{IJ}}=\sqrt{g}g^{\mu \nu }(_\mu \overline{\xi }_IE_{\nu J}+_\mu \overline{\omega }_{IK}A_{\nu J}{}_{}{}^{K}(IJ))+_\mu \left(\sqrt{g}g^{\mu \nu }_\nu \overline{\omega }_{IJ}\right).\hfill \end{array}$$
(2.12)
## 3 Ghost Equation and Vector Supersymmetry
Let us introduce the condensed notation
$$\begin{array}{cc}\{𝒜_\mu ^i,i=1,2\}=\{E_\mu ^I,A_\mu ^{IJ}\},\hfill & \\ \{\overline{C}_i,i=1,2\}=\{\overline{\xi }_I,\overline{\omega }_{IJ}\},\hfill & \{B_i,i=1,2\}=\{\lambda _I,b_{IJ}\},\hfill \end{array}$$
(3.1)
under which the gauge conditions (2.10) and the equation for the ghost $`\xi `$ (2.12) now read<sup>11</sup><sup>11</sup>11Summations $`_I`$ and $`\frac{1}{2}_{IJ}`$ over repeated indices are implicit.
$$\frac{\delta S}{\delta B_i}=_\mu \left(\sqrt{g}g^{\mu \nu }𝒜_\nu ^i\right),$$
(3.2)
$$\frac{\delta S}{\delta \xi ^\mu }=\underset{i=1,2}{}\left(\sqrt{g}g^{\nu \lambda }_\lambda \overline{C}_i_\mu 𝒜_\nu ^i+_\nu \left(\sqrt{g}g^{\nu \lambda }_\lambda \overline{C}_i𝒜_\mu ^i\right)\right).$$
(3.3)
What we claim here is that, under circonstances to be specified later on, the theory is invariant under the vector supersymmetry transformations (we use the notation (3.1))
$$\begin{array}{c}\delta _{(\epsilon )}^\mathrm{S}\xi ^\mu =\epsilon ^\mu ,\hfill \\ \delta _{(\epsilon )}^\mathrm{S}B_i=\epsilon ^\mu _\mu \overline{C}_i,\hfill \\ \delta _{(\epsilon )}^\mathrm{S}\phi =0,\phi \xi ^\mu ,B^i,\hfill \end{array}$$
(3.4)
where the infinitesimal parameter $`\epsilon ^\mu `$ is a vector field – taken as commuting, to the contrary of $`\xi ^\mu `$, so that the supersymmetry operator $`\delta _{(\epsilon )}^\mathrm{S}`$ is an antiderivation. The latter together with the BRST operator $`s`$ obey the superalgebra anticommutation relations
$$s^2\phi =0,(\delta _{(\epsilon )}^\mathrm{S})^2\phi =0,\{s,\delta _{(\epsilon )}^\mathrm{S}\}\phi =_\epsilon \phi ,$$
(3.5)
for all fields $`\phi `$, where $`_\epsilon `$ is the Lie derivative along the vector $`\epsilon `$.
In order to check the possible invariance of the action under the vector supersymmetry, we first note that this is trivially the case for the gauge invariant part (2.1) of the total action (2.11). Next, since the gauge fixing part (2.9) is a BRST variation:
$$S_{\mathrm{gf}}=_Md^4x\sqrt{g}g^{\mu \nu }\underset{i=1,2}{}s\left(_\mu \overline{C}_i𝒜_\nu ^i\right),$$
(3.6)
we can use the anticommutation relation (3.5) and thus write
$$\delta _{(\epsilon )}^\mathrm{S}S=_Md^4x\sqrt{g}g^{\mu \nu }\underset{i=1,2}{}_\epsilon \left(_\mu \overline{C}_i𝒜_\nu ^i\right).$$
(3.7)
Partial integrations<sup>12</sup><sup>12</sup>12We assume allthrough the absence of boundary terms contributions. then yield
$$\delta _{(\epsilon )}^\mathrm{S}S=_Md^4x\sqrt{g}\left(_\epsilon g^{\mu \nu }+_\lambda ^{(\mathrm{g})}\epsilon ^\lambda g^{\mu \nu }\right)\underset{i=1,2}{}\left(_\mu \overline{C}_i𝒜_\nu ^i\right),$$
(3.8)
where $`_\mu ^{(\mathrm{g})}`$ is the covariant derivative with respect to the background metric $`g_{\mu \nu }`$. The first conclusion is that, generically, the vector supersymmetry transformations (3.4) are not an invariance of the theory. However they will indeed represent an invariance if and only if the parenthesis in the integrant of the right-hand side of (3.8) vanishes:
$$_\epsilon g^{\mu \nu }+_\lambda ^{(\mathrm{g})}\epsilon ^\lambda g^{\mu \nu }=0,$$
(3.9)
which is easily shown to be equivalent to the condition that the vector $`\epsilon `$ be a Killing vector field of the background metric: $`g^{\mu \nu }`$:
$$_\epsilon g^{\mu \nu }=0,\text{or:}_\mu ^{(\mathrm{g})}\epsilon _\nu +_\nu ^{(\mathrm{g})}\epsilon _\mu =0.$$
(3.10)
In such a case, the vector supersymmetry invariance may be expressed by the functional identity (still using the notation (3.1))
$$\delta _{(\epsilon )}^\mathrm{S}S_Md^4x\epsilon ^\mu \left(\frac{\delta }{\delta \xi ^\mu }+\underset{i=1,2}{}_\mu \overline{C}_i\frac{\delta }{\delta B_i}\right)S=0.$$
(3.11)
It is illustrative to consider a flat background metric, e.g. the Minkowski one: $`g_{\mu \nu }`$ $`=`$ $`\eta _{\mu \nu }`$. In this case the general solution of the condition (3.10) reads
$$\epsilon ^\mu =a^\mu +b^{\mu \nu }x_\nu ,\text{with}a^\mu ,b^{\mu \nu }=b^{\nu \mu }\text{constants}.$$
(3.12)
The right hand side of the anticommutator in (3.5) is then an infinitesimal rigid Poincaré transformation of parameters $`a^\mu `$ and $`b^{\mu \nu }`$.
Note that one could have derived the identity (3.11) directly from the equation (3.3) for the diffeomorphism ghost $`\xi `$, integrated with the vector field $`\epsilon `$, and performing some partial integations. Thus we clearly see how the vector supersymmetry is linked to the diffeomorphism ghost equation. In this respect the situation differs from the one encountered in Yang-Mills topological theories (Chern-Simons, BF), where there is no such narrow relation between a ghost equation and the vector supersymmetry – although they both hold in a Landau type gauge, too.
Another difference with the topological Yang-Mills case is that the weaker condition (3.10) for supersymmetry invariance holds in the present case, whereas it reads there $`_\mu ^{(\mathrm{g})}\epsilon _\nu =0`$: the superalgebra (3.5) closes there on isometries of the translation type only - true translations instead of general Poincaré transformations for a flat background metric.
It is known that in the Yang-Mills case the vector symmetry operator may be expressed in the form of the so-called operator $`\delta `$ of Sorella used to construct the invariants of the theory, and characterized by the algebraic relation
$$[\delta ,s]=d,$$
(3.13)
where $`d`$ is the exterior derivative. In the present case, too, there exists such an operator $`\delta `$. And, remarkably, it is linked to our vector supersymmetry, hence to the diffeomorphism ghost equation, in the following way. Considering the supersymmetry transformation rules (3.4) for a constant<sup>13</sup><sup>13</sup>13Vector supersymmetry invariance will then hold for a flat Minkovskian background metric. vector field $`\epsilon `$, we define the action of the operator $`\delta `$ as given by these transformations, with $`\epsilon ^\mu `$ replaced by the differential $`dx^\mu `$:
$$\begin{array}{c}\delta \xi ^\mu =dx^\mu ,\hfill \\ \delta B_i=d\overline{C}_i,\hfill \\ \delta \phi =0,\phi \xi ^\mu ,B^i,\hfill \end{array}$$
(3.14)
which is the result of – up to the action on the Lagrange multiplier fields $`B_i`$, not considered there. We can easily check the commutation rule (3.13), and also that $`\delta `$ commutes with $`d`$. Note that, as in , we can write the first of eqs. (3.14) in a coordinate independent way as
$$\delta \eta ^I=E^I,$$
(3.15)
where $`\eta ^IE_\mu ^I\xi ^\mu `$ is the “tangent space translation ghost” .
Before concluding, we could ask for the role of the Lorentz ghost equation, the second of eqs. (2.12) The answer is much the same as in ordinary gauge theories for the Yang-Mills ghost equation . Integrating the Lorentz ghost equation in space-time, integrating by part and using the gauge conditions (2.10), we obtain
$$\delta _{IJ}^{(\mathrm{SL})}S_Md^4x\left(\frac{\delta }{\delta \omega ^{IJ}}\overline{\xi }_I\frac{\delta }{\delta \lambda ^J}+\overline{\xi }_J\frac{\delta }{\delta \lambda ^I}\overline{\omega }_{IK}\frac{\delta }{\delta b^J_K}+\overline{\omega }_{JK}\frac{\delta }{\delta b^I_K}\right)S=0,$$
(3.16)
which is very similar to the result of , the Yang-Mills gauge invariance being now replaced by the local Lorentz invariance. One also may check the anticommutation rule
$$\{s,\delta _{IJ}^{(\mathrm{SL})}\}=\delta _{IJ}^{(\mathrm{L})},$$
(3.17)
where the right-hand side is an infinitesimal generator of rigid Lorentz transformation.
## 4 Conclusions
We have found a direct relation between the diffeomorphism ghost equation (the first of eqs. (2.12)) and the existence of a vector supersymmetry (3.4) – or of the Sorella operator $`\delta `$ (3.14), (3.15). This appears to be a characteristic features of theories invariant under “active diffeomorphisms”, i.e. diffeomorphisms which act on the dynamical fields only<sup>14</sup><sup>14</sup>14In a quantum context, these are diffeomorphisms acting quantum mechanically on the field operators – vierbein, connection and matter fields .. In such theories, the diffeomorphism ghost $`\xi `$ is a dynamical field, which actually means that it obeys an equation of motion.
Our results do not depend on the specificity of the invariant action taken to define the theory. They obviously also hold in the case of a self-dual connection $`A^{IJ}`$, in which case the action (2.1) is that of Ashtekar . As we have mentionned, the presence of minimally coupled matter is allowed, as well as other type of actions, provided they share the same “topological-like” character, i.e. provided they are constructed with the vierbein and connection as independent variables, and with invariance under active diffeomorphisms,
We have also emphasized the differences of the present diffeomorphism invariant theory with respect to the topological theories for Yang- Mills fields, such as the role of the ghost equation, but also the existence of more supersymmetry thanks to a weaker condition of invariance.
During completion of this work, the author became aware of a recent preprint – published by now – which gives an alternative derivation of vector supersymmetry in topological theories. This method may be applied to the gravitational case, too – and has been indeed applied in the published version of .
### Acknowledgments
The author is very grateful to François Gieres for critical remarks and for having communicated to him results prior to publication. Thanks are due to one of the referees for suggestions of references and useful comments.
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# Abstract
## Abstract
Fermi Transport is useful for describing the behaviour of spins or gyroscopes following non-geodesic, timelike world lines. However, Fermi Transport breaks down for null world lines. We introduce a transport law for polarisation vectors along non-geodesic null curves. We show how this law emerges naturally from the geometry of null directions by comparing polarisation vectors associated with two distinct null directions. We then give a spinorial treatment of this topic and make contact with the geometric phase of quantum mechanics. There are two significant differences between the null and timelike cases. In the null case (i) The transport law does not approach a unique smooth limit as the null curve approaches a null geodesic. (ii) The transport law for vectors is integrable, i.e the result depends only on the local properties of the curve and not on the entire path taken. However, the transport of spinors is not integrable: there is a global sign of topological origin.
## 1 Introduction
The spin four-vector $`s^a`$ of a gyroscope (not acted on by external torques) moving along a timelike geodesic is parallel transported along the geodesic. Similarly, the polarisation vector $`v^a`$ of a light ray following a null geodesic is parallel transported along the null geodesic. If a gyroscope follows a timelike curve which is not a geodesic, the parallel transport rule no longer applies: parallel transport does not preserve the orthogonality of the spin four-vector to the tangent vector for non-geodesic curves. The correct transport law for spin four-vectors along timelike curves is Fermi transport. However, Fermi transport does not apply to null curves: the transport equation breaks down. What is the appropriate transport law for polarisation vectors along null, non-geodesic curves? The purpose of this paper is to answer this question.
We will start in section II with a review of Fermi transport and introduce in section III a new transport law for null, non-geodesic curves and describe its properties. We then show in section IV that this law derives naturally from the geometry of null vectors just as Fermi transport derives naturally from the geometry of timelike vectors. Section V is a more sophisticated spinorial discussion of the transport law. Section VI connects the spinorial discussion of Section V to the geometric phase of a two state quantum system. Section VII is a concluding discussion.
## 2 Review of Fermi Transport
Let $`(,g)`$ be a space-time manifold with a Lorentzian metric $`g`$ of signature $`(+,,,)`$. Let $`𝒞`$ be a smooth time-like curve in $``$ and $`p`$ a point on the curve. In local coordinates, the curve is described as $`x^a(\tau )`$, where $`\tau `$ is an arbitrary parameter which increases into the future. We define the tangent vector $`𝐭`$ by $`t^a=dx^a/d\tau `$ and the acceleration $`\dot{𝐭}`$ by $`\dot{t}^a=t^b_bt^a`$, where $`_b`$ is the covariant derivative. If $`s^a`$ is a vector at $`p`$ which is orthogonal to $`𝐭`$ ($`𝐬T_p,(𝐬𝐭)_p=0`$), Fermi transport gives us a vector at every point of $`𝒞`$ defined by the transport law
$$Ds^a/d\tau =F^{ab}s_b,$$
(1)
where
$$F^{ab}=\frac{(\dot{t}^at^bt^a\dot{t}^b)}{t^ct_c},$$
(2)
or more abstractly, $`𝐅=\dot{𝐭}𝐭/(𝐭𝐭)`$. The transport law (1) has the following properties:
1. Vectors orthogonal to the tangent vector at $`p`$ are transported to vectors orthogonal to the tangent vector at other points of $`𝒞`$.
2. The transport law is covariant under reparametrization of the curve. If $`\tau `$ is changed to $`\tau ^{}`$, where $`f:=d\tau /d\tau ^{}>0`$, $`𝐭`$ and $`\dot{𝐭}`$ transform as follows
$$𝐭f𝐭\dot{𝐭}f^2\dot{𝐭}+\alpha _1𝐭,$$
(3)
where $`\alpha _1`$ is some function on $`𝒞`$.
As a result $`𝐅f𝐅`$ and so the vector field defined on $`𝒞`$ by Fermi transport does not depend on the parametrization of $`𝒞`$.
3. Inner products between vectors $`𝐬_1`$ and $`𝐬_2`$ are maintained under Fermi transport
$$\frac{d(𝐬_1𝐬_2)}{d\tau }=\frac{D(𝐬_1𝐬_2)}{d\tau }=F^{ab}s_{1a}s_{2b}+F^{ab}s_{1b}s_{2a}=0$$
(4)
from the antisymmetry of $`F^{ab}`$. In particular, an orthonormal triad $`(𝐞_1,𝐞_2,𝐞_3)_p`$ of vectors orthogonal to $`𝐭`$ at $`p`$ can be Fermi transported along $`𝒞`$ to give an orthonormal triad orthogonal to $`𝐭`$ everywhere on $`𝒞`$.
4. For geodesics Fermi transport reduces to parallel transport, since $`𝐅`$ vanishes on geodesics.
Although the treatment given above is in the context of an arbitrary curved space-time, the notion of Fermi transport has nothing to do with curvature.This is evident because the entire discussion takes place in the neighborhood of a single open curve. By an appropriate choice of co-ordinates, all the Christoffel symbols can be made to vanish on this curve . The curvature of space-time thus plays no essential role in the discussion of Fermi transport along an open curve, which should be viewed as a special relativistic kinematic effect. It is thus possible to conduct the whole discussion in flat space–time. By parallel transport we can identify the tangent space at any point $`p^{}`$ along the curve $`𝒞`$ with $`T_p`$. The subsequent discussion is hence entirely within $`T_p`$.
Fermi transport arises from the geometry of timelike vectors in $`T_p`$. For a non–geodesic curve the vector field obtained by parallel transport of the tangent vector at one point does not agree with the local tangent vector at other points. Fermi transport allows us to set up a correspondence between vectors orthogonal to distinct timelike vectors, viz the tangent vectors at different points of a non-geodesic curve. This correspondence is geometrically natural but is not integrable. By integrability we mean that the result of transporting a given initial vector to a final point along the curve depens only on the local properties of the curve at the final point, and not on the entire path upto that point. In this sense, the Fermi transport rule is not integrable. This reflects the curvature of the space of timelike directions, which can be identified with the unit hyperboloid in Minkowski space.
## 3 Transport along Null Non-Geodesics
As is evident from (2), the Fermi transport rule breaks down for null curves, since the denominator of (2) vanishes. For geodesic null curves, polarisation vectors are parallel transported along the curve. But what about non-geodesic null curves? This is the question that we address and answer in this paper. In this section we will simply write down a transport law for polarisation vectors along null, non–geodesic curves. We will then show that this law has geometrically natural properties as does Fermi transport.
Let $`𝒩`$ be a smooth null curve described in local coordinates by $`x^a(\tau )`$, where $`\tau `$ is an arbitrary parameter which increases into the future. We use $`l^a=dx^a/d\tau `$ to denote the tangent vector to the null curve. We will also need $`\dot{l}^a=l^b_bl^a`$ and $`\ddot{l}^a=l^b_b\dot{l}^a`$. We will assume that $`𝒩`$ is nowhere geodesic i.e, $`\dot{𝐥}`$ is nowhere a scalar multiple of $`𝐥`$. By differentiating $`𝐥.𝐥=0`$ with respect to $`\tau `$, we deduce in succession $`𝐥.\dot{𝐥}=0`$ and $`\dot{𝐥}.\dot{𝐥}+𝐥.\ddot{𝐥}=0`$. Let us write $`H_pT_p`$ for the space of all vectors in $`T_p`$ which are orthogonal to $`𝐥`$ at $`p`$. $`H_p`$ also includes $`𝐥`$, since $`𝐥`$ is null. The Lorentz metric at $`p`$, pulled back to $`H_p`$, has signature $`(0,,)`$. We need to “mod out” by the null direction $`l^a`$ to get a non-degenerate metric. We define a polarization vector to be an equivalence class of vectors in $`H_p`$ which differ by a multiple of $`𝐥`$ ($`𝐯_1𝐯_2𝐯_2=𝐯_1+\lambda 𝐥`$ for some $`\lambda `$). Our transport law will tell us how to transport polarization vectors along $`𝒩`$. (This “modding out” is standard for null curves ).
Given a polarisation vector (an equivalence class) at $`p`$, let us take a representative element $`v^a`$ from $`H_p`$. We transport this vector along $`𝒩`$ using the rule
$$\frac{Dv^a}{D\tau }=K^{ab}v_b,$$
(5)
where $`K^{ab}`$ is an antisymmetric tensor defined by
$$K^{ab}=\frac{(\ddot{l}^a\dot{l}^b\dot{l}^a\ddot{l}^b)}{(\dot{l^c}\dot{l_c})}.$$
(6)
Finally we evaluate the equivalence class of $`v^a`$ (mod out by $`l^a`$) to produce a polarisation vector field along $`𝒩`$.
This law, designated $`K`$ transport from now on, has the following properties:
1. It is straightforward to check that polarization vectors (equivalence classes of vectors perpendicular to $`𝐥`$) at $`p`$ are transported to polarization vectors at other points of $`𝒩`$.
2. The transport law is reparametrization covariant: Under reparametrization of $`𝒩`$, $`𝐥f𝐥`$ and
$$\dot{𝐥}f^2\dot{𝐥}+\alpha _2𝐥,\ddot{𝐥}f^3\ddot{𝐥}+\alpha _3\dot{𝐥}+\alpha _4𝐥,$$
(7)
where the $`\alpha `$’s are some functions on $`𝒩`$. Using (7) in (5,6), we can drop the $`\alpha _2`$ and $`\alpha _4`$ terms, which are proportional to $`𝐥`$: they either vanish when they are contracted with $`v_b`$ in (5) or are modded out when we pass to polarisation vectors. As a result,
$$\dot{𝐥}\dot{𝐥}f^4\dot{𝐥}\dot{𝐥},\ddot{𝐥}\dot{𝐥}(f^2\ddot{𝐥}+\alpha _3\dot{𝐥})(f^3\dot{𝐥})=f^5\dot{𝐥}\ddot{𝐥}.$$
(8)
We find that
$$K^{ab}fK^{ab}$$
(9)
and so (5) provides a reparametrisation invariant transport law.
3. Inner products between polarization vectors are preserved. This follows from antisymmetry of $`K`$.
A point worth noting is that $`K`$ transport involves the second derivative $`\ddot{𝐥}`$ of the tangent vector $`𝐥`$. This is quite unlike Fermi transport which only involves the first derivative $`\dot{𝐭}`$ of the tangent vector. We will see below that this is an unavoidable consequence of the geometry of null vectors.
Unlike Fermi transport, the $`K`$ transport law does not have a smooth limit as the null curve becomes a geodesic. This will be discussed in the concluding section.
Note that we are not claiming to transport all vectors along $`𝒩`$ in a geometrically natural manner. Our rule is only meant for polarisation vectors, i.e. equivalence classes of vectors perpendicular to the tangent vector.
## 4 Comparing Polarisation vectors on distinct Null vectors
We will now show that the transport rule (5) originates naturally from the geometry of null directions, just as the Fermi transport rule (1) derives from the geometry of timelike directions . If we parallel transport the tangent vector $`𝐥_p^{}`$ along $`𝒩`$ from $`p^{}`$ to $`p`$, we find that for non-geodesic curves, the parallel transported tangent vector does not agree with the local tangent vector $`𝐥_p`$. The parallel tansport of a vector orthogonal to the tangent vector at $`p^{}`$ is not in general orthogonal to the local tangent vector $`𝐥_p`$, but to the parallel transported tangent vector. We need to find a way to compare vectors orthogonal to two distinct null vectors. We show below that this can be done in a geometrically natural manner for polarisation vectors.
Let $`M`$ be a four dimensional vector space with Lorentzian metric ($`\eta `$) of signature $`(+,,,)`$. ($`M`$ is a model for $`T_p`$, the tangent space at the point $`p`$ of $``$.) The set of future pointing ($`l^0>0`$), null ($`l^al^b\eta _{ab}=0`$) vectors in $`M`$ forms the future light cone and the set of null directions (defined as future pointing null vectors modulo extent) is a sphere. We will sometimes refer to this sphere as the celestial sphere or the sky (although, strictly speaking, this terminology should be reserved for past pointing null directions).
If $`L`$ is a null direction and $`𝐥`$ a null vector belonging to $`L`$, we define $`H_L`$ to be the space of vectors $`v^a`$ in $`M`$ orthogonal to $`L`$, $`v^al_a=0`$. $`H_L`$ also includes $`L`$. We define a polarisation vector $`𝐩_L`$ to be an equivalence class of vectors in $`H_L`$ differing by an arbitrary multiple of $`𝐥`$ ($`v^av^a+\lambda l^a`$ $`\lambda `$ arbitrary). The vector space of polarisation vectors defined by $`L`$ is written $`P_L`$.
Given two distinct null directions $`L_1`$ and $`L_2`$ and a polarisation vector $`𝐩_1P_{L_1}`$, there is a geometrically natural choice of a polarisation vector from $`P_{L_2}`$. We pick from $`𝐩_1`$ the unique element $`𝐰H_{L_1}`$ which is orthogonal to $`L_2`$ and define $`𝐩_2`$ as the class $`𝐰+\lambda _2𝐥_2`$ containing the vector $`𝐰`$. More explicitly, pick any $`𝒗_\mathrm{𝟏}P_{L_1}`$. Requiring that
$$(𝐯_1+\lambda _1𝐥_1)𝐥_2=0$$
(10)
uniquely fixes $`\lambda _1`$:
$$\lambda _1=\frac{𝐯_1𝐥_2}{𝐥_1𝐥_2},$$
(11)
which is well defined, since $`𝐥_1𝐥_2>0`$, for $`𝐥_1`$ and $`𝐥_2`$ distinct. We write
$$𝐰=𝐯_1\frac{𝐯_1𝐥_2}{𝐥_1𝐥_2}𝐥_1$$
(12)
and define $`𝐩_2`$ to be the equivalence class $`𝐰+\lambda _2𝐥_2`$ containing $`𝐰`$:
$$𝐯_2=𝐯_1\frac{𝐯_1𝐥_2}{𝐥_1𝐥_2}𝐥_1+\lambda _2𝐥_2.$$
(13)
Let us write $`\delta 𝐥_{12}=𝐥_2𝐥_1`$ and note that $`𝐥_2𝐥_2=0=(𝐥_1+\delta 𝐥_{12})(𝐥_1+\delta 𝐥_{12})=2\delta 𝐥_{12}𝐥_1+\delta 𝐥_{12}\delta 𝐥_{12}`$ implies
$$𝐥_1\delta 𝐥_{12}=\frac{1}{2}\delta 𝐥_{12}\delta 𝐥_{12}.$$
(14)
We can now write
$$𝐯_2=𝐯_1\frac{2𝐯_1\delta 𝐥_{12}}{\delta 𝐥_{12}\delta 𝐥_{12}}𝐥_1+\lambda _2𝐥_2$$
(15)
and by suitable choice of $`\lambda _2`$
$$𝐯_2=𝐯_1\frac{2(𝐯_1\delta 𝐥_{12})\delta 𝐥_{12}}{(\delta 𝐥_{12}\delta 𝐥_{12})}.$$
(16)
This reverses the $`\delta 𝐥_{12}`$ component of $`𝐯_1`$ and leaves the component of $`𝐯_1`$ which is orthogonal to $`\delta 𝐥_{12}`$ unchanged. The rule which associates polarisation vectors in $`P_1`$ to polarization vectors in $`P_2`$ reverses orientation and therefore cannot be continuously deformed to the identity. This can be seen quite clearly going to a frame in which $`𝐥_1`$ represents a light ray going in the positive $`z`$ direction and $`𝐥_2`$ a light ray going in the negative $`z`$ direction. Polarization vectors of $`𝐥_1`$ and $`𝐥_2`$ can then be identified with the $`xy`$ plane but with opposite orientations.
The rule for identifying $`P_1`$ and $`P_2`$ can also be stated succinctly as follows. Consider the two dimensional subspace of $`M`$ orthogonal to both $`𝐥_1`$ and $`𝐥_2`$. The projector on to this subspace is $`h_{ab}=\eta _{ab}(𝐥_1𝐥_2)^1(l_{1a}l_{2b}+l_{2a}l_{1b})`$. This subspace can be identified with both $`P_1`$ and $`P_2`$ (by taking its elements to represent classes) and this gives an identification of $`P_1`$ with $`P_2`$. As the example of the last paragraph shows, the identification reverses orientation. The natural volume form $`ϵ_{ab}:=(𝐥_1𝐥_2)^1ϵ_{abcd}l_1^cl_2^d`$ on the two dimensional subspace reverses sign when $`𝐥_1`$ and $`𝐥_2`$ are interchanged.
Since the map from $`P_1`$ to $`P_2`$ cannot be continuously deformed to the identity, it does not have a smooth limit as $`L_2`$ tends to $`L_1`$. Such a smooth limit would be necessary to define transport along a smooth null curve. However, if we repeat the process by considering three null directions $`L_1`$, $`L_2`$, $`L_3`$ and go from $`P_1`$ to $`P_2`$ and $`P_2`$ to $`P_3`$ using the rule (16) twice, we get an orientation preserving map from $`P_1`$ to $`P_3`$. This map does have a smooth limit as $`P_1`$, $`P_2`$ and $`P_3`$ approach each other. It is intuitively clear that since we need three null directions (rather than two) to take a smooth limit, the transport law we derive for null vectors will depend on $`𝐥,\dot{𝐥}`$ and $`\ddot{𝐥}`$, in contrast to the Fermi transport law, which only depends on $`𝐭`$ and $`\dot{𝐭}`$. Writing $`𝐥_1`$, $`𝐥_2`$, $`𝐥_3`$ for elements of $`L_1`$, $`L_2`$, $`L_3`$ and $`\delta 𝐥_{23}=(𝐥_2𝐥_3)`$ we find as in (16),
$$𝐯_3=𝐯_2\frac{2𝐯_2\delta 𝐥_{23}}{\delta 𝐥_{23}\delta 𝐥_{23}}\delta 𝐥_{23}$$
(17)
using (16) we arrive at the following expression for $`𝐯_3𝐯_1`$
$$\frac{4(𝐯_1\delta 𝐥_{12})(\delta 𝐥_{12}\delta 𝐥_{23})\delta 𝐥_{23}2(𝐯_1\delta 𝐥_{12})(\delta 𝐥_{23}\delta 𝐥_{23})\delta 𝐥_{12}2(𝐯_1\delta 𝐥_{23})(\delta 𝐥_{12}\delta 𝐥_{12})\delta 𝐥_{23}}{(\delta 𝐥_{12}\delta 𝐥_{12})(\delta 𝐥_{23}\delta 𝐥_{23})}.$$
(18)
We can now take the limit as the three null directions approach each other. Let $`𝐥(\tau )`$ be the tangent vector of a smooth null curve. We expand $`𝐥(\tau )`$ in a Taylor series.
$`𝐥_1`$ $`=`$ $`𝐥(\tau \mathrm{\Delta }\tau )=𝐥(\tau )\dot{𝐥}(\tau )\mathrm{\Delta }\tau +1/2\ddot{𝐥}(\tau )(\mathrm{\Delta }\tau )^2+\mathrm{}`$ (19)
$`𝐥_2`$ $`=`$ $`𝐥(\tau )`$ (20)
$`𝐥_3`$ $`=`$ $`𝐥(\tau +\mathrm{\Delta }\tau )=𝐥(\tau )+\tau \dot{𝐥}(\tau )\mathrm{\Delta }+1/2\ddot{𝐥}(\mathrm{\Delta }\tau )^2+\mathrm{}`$ (21)
where the ellipsis stand for higher order terms than we need. As a result,
$`\delta 𝐥_{12}`$ $`=`$ $`\dot{𝐥}(\tau )\mathrm{\Delta }\tau 1/2\ddot{𝐥}(\mathrm{\Delta }\tau )^2+\mathrm{}`$ (22)
$`\delta 𝐥_{23}`$ $`=`$ $`\dot{𝐥}(\tau )\mathrm{\Delta }\tau +1/2\ddot{𝐥}(\mathrm{\Delta }\tau )^2+\mathrm{}`$ (23)
The leading term in the denominator of (18) is of order $`(\mathrm{\Delta }\tau )^4`$:
$$(\dot{𝐥}(\tau )\dot{𝐥}(\tau ))(\mathrm{\Delta }\tau )^4.$$
(24)
In the numerator the term of order $`(\mathrm{\Delta }\tau )^4`$ is
$$[4(𝐯_1\dot{𝐥})(\dot{𝐥}\dot{𝐥})\dot{𝐥}2(𝐯_1\dot{𝐥})(\dot{𝐥}\dot{𝐥})\dot{𝐥}2(𝐯_1\dot{𝐥})(\dot{𝐥}\dot{𝐥})\dot{𝐥}](\mathrm{\Delta }\tau )^4$$
(25)
which vanishes. The first non vanishing term is of order $`(\mathrm{\Delta }\tau )^5`$. After some straight forward algebra, we evaluate $`(𝐯_3𝐯_1)/(2\mathrm{\Delta }\tau )`$ and find that the limit $`\mathrm{\Delta }\tau 0`$ exists and yields the $`K`$ transport law (5) of section II.
As the reader may have noticed, the entire discussion of this section depends only on the conformal metric and not the metric itself. The definition of $`H_p`$ and “modding out” by $`𝐥`$ are unchanged under conformal transformations and so is the rule (16) for comparing polarisation vectors between fibres. One may therefore expect that the $`K`$ transport rule is conformally invariant. It is easily checked that it is. The $`K`$ transport of a vector using a conformally rescaled metric only results in trivial rescalings of the polarisation vector (to be expected because parallel transport preserves the norm). Under conformal transformations there is no change in the direction of polarisation of the $`K`$ transported vector.
## 5 Spinorial formulation
As one might expect, the discussion of the last section can be formulated quite naturally in terms of spinors . Let $`(V,ϵ_{AB})`$ be a complex 2-dimensional vector space endowed with an antisymmetric non-degenerate tensor $`ϵ_{AB}`$. Elements of $`V`$ are written $`\xi ^A`$. The complex conjugate of $`\xi ^A`$ is written $`\overline{\xi }^A^{}`$, with $`A^{}`$ a “primed” or “dotted” spinor index. A pair $`AA^{}`$ of spinor indices can be converted into a vector index $`a`$ by using the standard correspondence between vectors and spinors:
$$v^a=\sigma _{AA^{}}^av^{AA^{}},$$
(26)
The components of $`\sigma ^a`$ are $`\sigma ^0=I`$, $`\sigma ^1=\sigma ^x`$,$`\sigma ^2=\sigma ^y`$,$`\sigma ^3=\sigma ^z`$, and $`I`$ is the $`2\times 2`$ identity matrix and $`(\sigma ^x,\sigma ^y,\sigma ^z)`$ are the standard Pauli matrices. We will write such relations (26) as $`v^av^{AA^{}}`$.The spinor $`\xi ^A`$ defines a future pointing null vector $`𝐥^a\overline{\xi }^A^{}\xi ^A`$. Altering $`\xi ^A`$ by a phase does not alter the vector all and multiplying $`\xi ^A`$ by a real number alters the extent of the null vector, but not its direction.
Define the following equivalence relation on (non-zero elements of) $`V`$:
$$\xi ^A\alpha \xi ^A,$$
(27)
where $`\alpha `$ is any non zero complex number. The set of equivalence classes form a sphere $`S^2`$, which is the set of future pointing null directions - the sky of the previous section. Non zero elements of $`V`$ form a fibre bundle with the base equal to $`S^2`$ and the fibre isomorphic to the set of non zero complex numbers. The phase of this non-zero complex number determines a “flag plane” or polarisation direction. This is easily seen as follows. Let $`L_1`$ be a point on $`S^2`$ and $`\xi _1^A`$ a point on the fibre over $`L_1`$. The null vector $`l^a\overline{\xi }_1^A^{}\xi _1^A`$ belongs to the null direction $`L_1`$. Let us pick an arbitrary spinor $`\xi _2`$, distinct from $`\xi _1`$, so that $`\xi _{2A}\xi _1^A`$ is nonzero. One can always multiply it by a suitable complex number so that $`\xi _1^A\xi _{2A}=1`$. The space-like unit vector
$$v^a\frac{1}{\sqrt{2}}(\overline{\xi }_1^A^{}\xi _2^A+\overline{\xi }_2^A^{}\xi _1^A)$$
(28)
is clearly orthogonal to $`l_{1a}\overline{\xi }_1^A^{}\xi _1^A`$
$$l_{1a}v^a=\frac{1}{\sqrt{2}}\overline{\xi }_{1A^{}}\xi _{1A}(\overline{\xi }_1^A^{}\xi _2^A+\overline{\xi }_2^A^{}\xi _1^A)=0$$
(29)
A different choiceof $`\xi _2`$, obtained by adding a multiple of $`\xi _1`$ to it, only changes the vector $`v^a`$ in (28) by a multiple of $`l_1^a`$. Thus $`\xi _1^A`$ determines an equivalence class of unit vectors $`v^a`$ orthogonal to $`l_a`$, i.e., a unit polarization vector. As can be easily verified, altering $`\xi _1^A`$ by a phase $`e^{i\theta }`$ leads to a rotation of the polarization vector by an angle $`2\theta `$. Thus, $`\xi _1^A`$ and $`\xi _1^A`$ define the same polarization vector. The correspondence is two to one. We thus have a map from the fibre $`(L)`$ over a null direction $`L`$ to the unit circle of polarization vectors defined by $`L`$. Our discussion now will be entirely on the spinor bundle.
Given two distinct points $`L_1`$ and $`L_2`$ on the base and a point $`\xi _1^A`$ on the fibre $`(L_1)`$ over $`L_1`$, there is a natural way to pick a point $`\xi _2^A`$ the $`(L_2)`$ fibre over $`L_2`$. We pick the unique point $`\xi _2^A`$ which satisfies
$$\xi _1^A\xi _{2A}=1.$$
(30)
(This choice when translated into vectors agrees with the discussion of section IV). The rule (30) is well defined only if $`L_1`$ and $`L_2`$ are distinct points. If $`\xi _1^A`$ is altered by a phase $`\xi _1^Ae^{i\theta }\xi _1^A`$, $`\xi _2^A`$ picks up the opposite phase: $`\xi _2^Ae^{i\theta }\xi _2^A`$. Thus, rule (30) maps a circle winding in the anticlockwise sense to a circle winding in the clockwise sense. The map defined by (30) from $`(L_1)`$ to $`(L_2)`$ cannot be continuously deformed to the identity and the rule (30) does not admit a smooth limit as $`L_2`$ approaches $`L_1`$.
As in Section 4, we can solve this problem by considering three points, $`L_1`$, $`L_2`$, $`L_3`$ on $`S^2`$. Given $`\xi _1^Aϵ(L_1)`$ we pick $`\xi _2^A`$ from $`(L_2)`$ accordingly to the rule (30) and repeat the process to pick $`\xi _3^A`$ from $`L_3`$ using $`\xi _2^A\xi _{3A}=1`$. The map from $`L_1`$ to $`L_3`$ does admit a smooth limit as $`L_1,L_2`$ and $`L_3`$ approach each other. We will use this below to derive the spinorial form of the transport law (5). If $`L_1,L_2`$ and $`L_3`$ are three distinct null directions, the map from $`L_1`$ to $`L_3`$ (via $`L_2`$) does depend on $`L_2`$. If a different choice $`L_2^{}`$ is made, one can check that the point $`\xi _3`$ on $`(L_3)`$ determined by $`\xi _1`$ is multiplied by a complex number $`\chi `$, where
$$\chi =\frac{(\xi _1^A\xi _{2^{}A})(\xi _2^B\xi _{3B})}{(\xi _1^D\xi _{2D})(\xi _2^{}^C\xi _{3C})}.$$
(31)
$`\chi `$ depends only on the four null directions $`L_1,L_2,L_3`$ and $`L_2^{}`$ and not on the representatives chosen from each fibre. $`\chi `$ is called the cross ratio of these four null directions. The fact that the map from $`L_1`$ to $`L_3`$ does depend on $`L_2`$, ($`\chi `$ is not the identity) shows that the discrete rule for comparing fibres over distinct null directions is not integrable.
We will now take the continuous limit of the discrete rule and recover the transport law (5) in spinorial language. We are given a curve $`L(\tau )`$ of null directions and a point $`\xi (0)`$ on the fibre over $`L(0)`$. What we seek is a geometrically natural “lift” of this curve i.e. we need to find $`\xi (\tau )`$ so that $`\xi (\tau )(L(\tau ))`$.
Fix a spin frame $`(\dot{i}^A,o^A)(i^Ao_A=1)`$ and write $`\xi (\tau )=\gamma (\tau )(i^A+z(\tau )o^A)`$. $`z`$ is a stereographic coordinate on the set of null directions and $`\gamma `$ is a coordinate on the fibre. The problem now is: give a smooth curve $`z(\tau )`$, determine $`\gamma (\tau )`$ using the rule (16). We expand $`z(\tau )`$in a Taylor series and write
$`z_1`$ $`=`$ $`z(\tau \mathrm{\Delta }\tau )=z(\tau )\dot{z}\mathrm{\Delta }\tau +{\displaystyle \frac{1}{2}}\ddot{z}(\mathrm{\Delta }\tau )^2+\mathrm{}`$ (32)
$`z_2`$ $`=`$ $`z(\tau )`$ (33)
$`z_3`$ $`=`$ $`z(\tau +\mathrm{\Delta }\tau )=z(\tau )+\dot{z}\mathrm{\Delta }\tau +{\displaystyle \frac{1}{2}}\ddot{z}(\mathrm{\Delta }\tau )^2.`$ (34)
Evidently,
$`\xi _1`$ $`=`$ $`\xi (\tau +\mathrm{\Delta }\tau )=\gamma (\tau \mathrm{\Delta }\tau )(i^A+z(\tau \mathrm{\Delta }\tau )o^A)`$ (35)
$`\xi _2`$ $`=`$ $`\xi (\tau )=\gamma (\tau )(i^A+z(\tau )o^A)`$ (36)
$`\xi _3`$ $`=`$ $`\xi (\tau +\mathrm{\Delta }\tau )=\gamma (\tau \mathrm{\Delta }+\tau )(i^A+z(\tau +\mathrm{\Delta }\tau )o^A).`$ (37)
Using the rule (16) for determining $`\xi _2`$ and $`\xi _3`$ we find from $`\xi _1^A\xi _{2A}=\xi _2^A\xi _{3A}=1`$,
$`\gamma (\tau )\gamma (\tau \mathrm{\Delta }\tau )(z(\tau )z(\tau \mathrm{\Delta }\tau )`$ $`=`$ $`1`$ (38)
$`\gamma (\tau )\gamma (\tau +\mathrm{\Delta }\tau )(z(\tau +\mathrm{\Delta }\tau )z(\tau )`$ $`=`$ $`1.`$ (39)
We eliminate $`\gamma (\tau )`$ from these equations and find using the Taylor expansion for $`z(\tau )`$
$$\frac{\gamma (\tau +\mathrm{\Delta }\tau )\gamma (\tau \mathrm{\Delta }\tau )}{\gamma (\tau \mathrm{\Delta }\tau )}=\frac{(\dot{z}\mathrm{\Delta }\tau 1/2\ddot{z}(\mathrm{\Delta }\tau )^2+\mathrm{})(\dot{z}\mathrm{\Delta }\tau +1/2\ddot{z}(\mathrm{\Delta }\tau )^2+\mathrm{})}{(\dot{z}\mathrm{\Delta }\tau +\mathrm{})}.$$
(40)
Evaluating
$$\frac{\gamma (\tau +\mathrm{\Delta }\tau )\gamma (\tau \mathrm{\Delta }\tau )}{\gamma (\tau \mathrm{\Delta }\tau )2\mathrm{\Delta }\tau }$$
we find in the limit $`\mathrm{\Delta }\tau 0`$
$$\gamma ^1\dot{\gamma }=(1/2)\ddot{z}/\dot{z}$$
(41)
which is the transport law (5) in spinorial form. Its this form, it is apparent that (41) can be integrated to yield
$$\gamma (\tau )=\gamma (0)\sqrt{(\dot{z}(0)/\dot{z}(\tau ))}.$$
(42)
Although the discrete rule (30) for comparing points on distinct fibres is not integrable, its continuous limit (the K-transport rule) is. This is a feature of K–transport which is not shared by Fermi transport. It is interesting that the discrete rule is not integrable, while the continuum limit is. This indicates that as the four null directions $`L_1,L_2,L_2^{},L_3`$ tend to each other, the phase discrepancy between alternative paths vanishes sufficiently fast that there is none left in the continuum limit.
## 6 Relation to the Geometric Phase in two state quantum mechanics
The discussion so far has been entirely Lorentz invariant. In particular the last section has been in the language of $`SL(2,IC)`$ spinors. We will now attempt to make contact with the notion of the Geometric Phase in quantum mechanics . The motivation is as follows. One can think of the spinors of the previous section as describing the state vectors of a two state quantum mechanical system. The equivalence relation (27) defining the fibres is precisely the one which takes one from Hilbert space to the space of physical states (ray space). The relation (30) associates a member $`\xi _2`$ of the fibre over the point $`L_2`$ is associated with a particular member $`\xi _1`$ of the fibre over $`L_1`$. In quantum mechanics, there is a notion of two state vectors, corresponding to different rays, being ”in phase”, which leads to the geometric phase. It is natural, therefore, to ask whether the relation (30) has an analogue in quantum mechanics. At first sight, there is an obstacle. The relation (30) breaks down when the two null vectors $`L_1`$ and $`L_2`$ are coincident. The geometric phase convention in quantum mechanics breaks down when the two rays are orthogonal. Nevertheless, there is a close correspondence which is developed in this section. In order to do this we need to reduce the structure group from $`SL(2,IC)`$ to $`SU(2)`$. The $`SL(2,IC)`$ invariant structures that we described in the last section will now be described in terms of $`SU(2)`$ spinors.
In order to break the structure group down from $`SL(2,IC)`$ to $`SU(2)`$, we introduce on the two complex dimensional vector space $`(V,ϵ_{AB})`$ an additional structure: a positive definite Hermitian inner product $`G_{A^{}A}`$. One can think of $`G_{AA^{}}`$ as the spinor corresponding to a timelike four vector. Thus making a choice of $`G`$ is like making a choice of the four-velocity of a frame of reference, which still leaves freedom to make spatial rotations. The group of transformations that preserves both $`ϵ_{AB}`$ and $`G_{AA^{}}`$ is $`SU(2)`$. By choice of spin frame $`(\iota ^A,o^A)`$ we can arrange that
$$G_{A^{}A}=\iota _A\iota _A^{}+o_Ao_A^{}$$
(43)
and use $`G_A^{}^A`$ to define a $``$ operation taking a spinor $`\xi ^A`$ to a new spinor $`\xi ^A`$ transforming in the same way.
$$\xi ^A:=\overline{\xi }^A^{}G_A^{}^A$$
(44)
We will sometimes use Dirac notation $`|\xi >`$ for the element $`\xi ^A`$ of $`V`$ and $`<\xi |`$ for the element $`\xi _A^{}`$ of $`V^{}`$ (the dual of $`V`$). Note that $`<\xi |`$ is not $`\xi _A`$, for $`<\xi |\xi >`$ is positive definite whereas $`\xi _A\xi ^A`$ vanishes. It is easily checked that $`\xi _A^{}=\xi _A`$ and that $`\xi ^{}`$ is orthogonal to $`\xi ,i.e<\xi ^{}|\xi >=0`$. The action of $``$ on the sphere of null directions (the sky) is easy to visualise. By explicit computation are sees that $`\iota ^{}=o`$, $`o^{}=\iota `$ and so, if $`\xi ^A=i^A+zo^A`$, $`\xi ^{}{}_{}{}^{A}=\overline{z}i^A+o^A`$, so $``$ sends each point on the sky to its antipode. The subgroup of $`SL(2,IC)`$ which preserves the relation of antipodality is $`SU(2)`$, which acts on the sky by rotations.
We now identify the $`SU(2)`$ spinors with state vectors of a two state system and the sky with the corresponding ray space, which is a sphere. (Historically, this arose in the context of polarised light, which can be represented by a pair of complex numbers, and the sphere was discovered by Poincaré after whom it is named. The definition that two states are in phase when their inner product is real and positive was proposed by Pancharatnam ). The rule $`\xi _1^A\xi _{2A}=1`$ for comparing points on distinct fibres can be rewritten as
$$\xi _{1A}^{}\xi _2^A=<\xi _1^{}|\xi _2>=1$$
(45)
i.e, we require that $`|\xi _2>`$ be in phase with $`|\xi _1^{}>`$. (We are not concerned here with the modulus of the complex number $`<\xi _1^{}|\xi _2>`$, but only its phase. ) The rule (45) is well defined if $`\xi _1`$ and $`\xi _2`$ are on distinct fibres, or, equivalently, if $`\xi _2`$ and $`\xi _1^{}`$ are not antipodal.
As in the case of vectors, the rule (30) for passing from the fibre $`(L_1)`$ over $`L_1`$ to the fibre $`(L_3)`$ over $`L_3`$ (via $`L_2`$ ) will depend on $`L_2`$. This dependence is captured by the cross ratio (31). The phase of the complex number (31) is a measure of the non-integrable nature of the rule (30) for comparing points on distinct fibres. We can rewrite this quantity as the phase of the complex number
$$<\xi _1|\xi _2^{}><\xi _2^{}|\xi _3><\xi _3|\xi _2^{}^{}><\xi _2^{}^{}|\xi _1>,$$
(46)
which has a simple geometric interpretation. Consider the four points $`L_1,\stackrel{~}{L_2},\stackrel{~}{L_2^{}}`$ and $`L_3`$ on the celestial sphere, where $`\stackrel{~}{L_2}`$ and $`\stackrel{~}{L_2^{}}`$ are points antipodal to $`L_2`$ and $`L_2^{}`$ respectively. The phase of $`\chi `$ has an interpretation which is well known in the Geometric phase literature : it is equal to half the solid angle subtended at the center of the sphere by the geodesic rectangle $`L_1,\stackrel{~}{L_2},L_3,\stackrel{~}{L_2^{}},L_1`$. It follows that the change in the plane of polarisation in following the route $`L_1,\stackrel{~}{L_2},L_3,\stackrel{~}{L_2^{}},L_1`$ is equal to the solid angle subtended by this rectangle. Although this change in the plane of polarisation has been computed in language pertaining to a given frame of reference, it is of course Lorentz invariant, from the earlier discussion.
## 7 Concluding Discussion
We have presented a transport law (5) which is the replacement for Fermi transport in the case of null curves. We have also shown how this transport law arises naturally from the geometry of null vectors. The $`K`$ transport law has a natural description in terms of $`SL(2,IC)`$ spinors. This description also brings out close analogies with the geometric phase, once it is specialised to $`SU(2)`$ spinors by choosing a timelike observer. In the rest of this section we compare the $`K`$ transport law and Fermi transport.
The main difference between Fermi transport and $`K`$ transport is due to the difference between the geometry of timelike directions and the geometry of null directions. The set of timelike directions can be identified with a time-like 3-hyperboloid, whose isometry group is the entire Lorentz group. In contrast, the null directions are identified with a 2-sphere and the Lorentz group acts on the sphere by conformal transformations. There is consequently no Lorentz invariant meaning to the statment that two null directions are “near” each other. By a suitable Lorentz transformation, any two distinct null directions can be made antipodal. As a result, there is no Lorentz and reparametrization invariant measure of the “acceleration” of a null curve. If the direction of the tangent vector of a null curve changes “slightly” in one Lorentz frame, this deviation can be made as large as one pleases in some other Lorentz frame.
Fermi transport reduces smoothly to parallel transport when the timelike curve becomes a timelike geodesic. In contrast, the transition from null curves to null geodesics is a singular one. This is reflected in the absence of a smooth limit for $`K`$ transport. As an example, let $``$ be Minkowski space with standard $`(t,x,y,z)`$ Cartesian co-ordinates and consider the null curve: $`x=R\mathrm{cos}(\mathrm{\Omega }\tau ),y=R\mathrm{sin}(\mathrm{\Omega }\tau ),z=\tau ,t=(\sqrt{1+R^2\mathrm{\Omega }^2})\tau `$, where $`R`$ and $`\mathrm{\Omega }`$ are constants. This curve describes a particle moving at the speed of light along a helical path.
The $`K`$ tensor for this null curve is easily worked out to be $`K=\mathrm{\Omega }dxdy`$ and is independent of $`R`$. The transport rule simply says that the polarisation vector rotates about the $`z`$-axis with angular velocity $`\mathrm{\Omega }`$. This is of course the angular velocity of the frame made up by the spatial tangent and the normal, i.e the Serret-Frenet frame.
In the limit that $`R`$ tends to $`0`$ with $`\mathrm{\Omega }`$ finite, the null curve does become a geodesic curve. However, $`K`$-transport does not reduce to parallel transport Thus, the limit of $`K`$ transport to null geodesics is not smooth. In fact, one can easily do the above Minkowski-space calculation for motion at the speed of light along a general smooth space curve whose spatial curvature nowhere vanishes. We can choose the length of the curve as a parameter, and choose polarisation vectors to be purely spatial. As in the case of the helix, one finds that $`K`$ transport reduces to Serret-Frenet transport, which is an integrable rule. One other way to check this is to compute the rate of change of the cosine of the angle made by the transported vector $`𝐯`$ with the acceleration vector $`\dot{𝐥}`$. A straightforward calculation using $`K`$–transport shows that
$$\frac{d}{d\tau }(\frac{𝐯\dot{𝐥}}{\sqrt{\dot{𝐥}\dot{𝐥}}})=0$$
Since the acceleration vector vanishes for geodesics, the rule must and does become ill-defined in the geodesic limit, as does Serret-Frenet transport.
Although the discrete rule given above for comparing polarisation vectors on distinct fibres is not integrable, the limit of this rule for smooth curves is integrable. This is a feature of the $`K`$ transport rule, which is different from Fermi transport. We note that if the rule had not been integrable in the continuous limit, we would have been able to define a two-form on the base space (the sky) whose integral around a closed curve would give the total rotation on traversing that curve. But, as is well known, there is no Lorentz invariant notion of area on the space of null vectors, except the one which assigns zero to each such area. With hindsight, the integrability is consistent with, and even forced by, Lorentz invariance.
Finally, we close with the remark that on spinors, $`K`$ transport is only locally integrable. (42) fixes $`\gamma (\tau )`$ only up to a sign. This sign is unimportant in discussing the transport of polarisation vectors. However, in transporting spinors along null, non–geodesic curves one can see a non integrable phase of topological origin with values $`\pm 1`$. The result of spinor transport is not affected by continuous deformations of the null curve $`𝒩`$ within the class of null, everywhere non–geodesic curves, connecting the end points of $`𝒩`$, but is affected by changes in $`𝒩`$ which cannot be continuously deformed away. Let $`z(\tau )`$ be a smooth closed simple (i.e non self-intersecting) curve in the space of null directions, with $`\dot{z}(\tau )`$ being nowhere zero. The complex number $`\dot{z}(\tau )`$ encircles the origin once in the complex plane for such a curve. Equation (42) then shows that transporting a spinor once along this curve results in a phase difference of $`\pi `$, which is in principle observable by interference. Acknowledgments: J.S. thanks V.P. Kattabomman for a discussion on transport rules.
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# 1 Introduction
## 1 Introduction
Recently a powerful apparatus to work with BPS excitations in superstring theory has been developed (see for some reviews). Despite this progress, we still have a poor understanding of the dynamics when SUSY is violated. In fact, even in the simplest situation of $`D`$-anti-$`D`$-brane ($`D\overline{D}`$) systems we know only the details of their topological content rather than the dynamics of the annihilation process .
First, it is argued that the annihilation process of a $`D\overline{D}`$-system is related to the tachyon rolling down to the bottom of its energy functional . The tachyon field here describes the lowest energy excitations of strings stretched between the $`D`$\- and $`\overline{D}`$-branes . This excitation has an imaginary mass which signals an instability in the $`D\overline{D}`$-system and leads to the annihilation. Second, knowing which charges are excited, one could trace what kinds of branes should be remnants of the annihilation . It is not clear, though, how to determine their position in their moduli space after the annihilation has happened. Moreover, it is not known how the tachyon rolls down to the bottom of its energy functional during the annihilation process. The main obstacle is the lack of knowledge of the tachyon potential .
Rather than looking for the tachyon potential in this note we would like to implement another approach. We know that in many occasions $`D`$-branes supply a good microscopic description of various low energy physics phenomena. In particular, the closest example to our approach is presented in ref. . In this paper the $`D`$-brane description of instantons in SUSY Yang-Mills (SYM) theories is given. Concretely, the gauge connection corresponding to the YM instanton is recovered from a microscopic theory which describes a $`D1`$-brane in the $`D5`$-brane background in type I string theory. What is most important for us is that the $`D1`$-brane theory in question is absolutely restricted by its (4,0) SUSY invariance .
Inspired by these considerations, we would like to probe a $`D9`$-$`\overline{D9}`$-brane annihilation by a $`D1`$-brane in type I string theory. Rather than dealing with the light-cone action for the $`D1`$-brane we consider its covariant Green-Schwarz (GS) formulation . For the theory to be non-anomalous it is necessary to consider the number of $`D9`$-branes to be the number of $`\overline{D9}`$-branes plus 32 .
The annihilation process is viewed on the $`D1`$-brane as a renormalization-group (RG) flow . In fact, after the annihilation some of the strings stretched between the $`D1`$-brane and the $`D9\overline{D9}`$-system should become massive and decouple from the IR limit of the probe theory. In the limit the theory describes the $`D1`$-brane in the background of only 32 $`D9`$-branes with some gauge bundle on the latter.
Here we study the theory on the probe brane which is already a low energy intermediate step in the RG evolution. It is an approximation to an as yet unknown microscopic theory which contains both $`D`$\- and $`\overline{D}`$-branes. Thus, we do not expect to recover from our probe theory explicitly the way the tachyon rolls down to the bottom of its energy functional. However, we could hope to extract from it some information about tachyon classical solutions which respect SUSY, i.e. tachyon values after the annihilation. In particular, one would hope that there is some symmetry which restricts possible background values of the tachyon after the annihilation . In fact, before the annihilation there is some non-SUSY theory which describes the $`D1`$-brane in the presence of both $`D9`$\- and $`\overline{D9}`$-branes. Via RG flow the theory evolves to a superconformal limit with a proper background value of the tachyon. We believe that there should be some hidden (non-linearly realized) SUSY of the theory, which forces it to flow in such a rigid way .
It is at this point that $`\kappa `$-invariance comes into the game. In fact, as is well established , this symmetry is related to a linearly realized SUSY on the world-sheet: one could formulate the $`D1`$-brane theory with an explicit SUSY both on the world-sheet and in the target space. Then $`\kappa `$-symmetry appears from the world-sheet SUSY after the integration over auxiliary fields . Hence, the presence of the $`\kappa `$-invariance is a sign that there is a linearly realized SUSY on the $`D1`$-brane world-sheet. In our case, we expect the invariance to be present only after the annihilation is completed.
It is for this reason that we are looking for backgrounds for the $`D1`$-brane, which respect $`\kappa `$-invariance. Conditions that invariance imposes on the theory constrain the possible values of the tachyon field. We find some equations which establish that this field is covariantly constant on light-like surfaces in the target space. The latter should be supplemented by integrability conditions so that if there is a gauge field background turned on the $`D9`$\- and $`\overline{D9}`$-branes, the tachyon could be a non-trivial field rather than just a constant.
Let us explain why we have equations for the tachyon field which are linear rather than quadratic in differentials. In fact, $`\kappa `$-invariance of a superstring theory in a background of SYM fields puts the latter on mass-shell, i.e. one gets second order differential (classical) equations (of motion) for the fields . Hence, the appearance of the first order differential equation for the tachyon field might seem suspicious. As we already mentioned, however, the $`\kappa `$-invariance is related to SUSY of the world-sheet theory. Also the mere presence of the tachyon field explicitly violates SUSY in the theory. Thus, we expect SUSY to be linearly realized only for some specific tachyon values. That is the reason we get BPS like linear differential equations for the tachyon field.
Our main interest is in a soliton which is of co-dimension four within the $`D9`$-branes. Only in such a situation is there a SUSY vacuum in the probe theory . In this case, the tachyon and background gauge field are functions of only four coordinates rather than ten. Then the integrability condition in question is just the self-duality equation for the gauge field. At large distances its instanton solution is represented as a pure gauge. The gauge matrix in the latter is equal to a non-trivial map (of a degree equal to the instanton charge) from $`S^3`$ at infinity to the group of Chan-Paton (CP) indexes. Specifically, we take the target of the map in question to be the diagonal $`USp(4k)`$ subgroup of the $`USp(4k)\times USp(4k)`$ group.
This choice could be clarified as follows. First, it should be stressed that we are considering a minimal construction of $`D5`$-branes from a $`D9`$-$`\overline{D9}`$-system. This construction is related to that due to Atiyah, Drinfeld, Hitchin and Manin (ADHM) for instantons in YM theories . In principle, one could study other situations which lead to non-minimal generalizations of the ADHM construction . For ”minimality” we consider a background gauge field on the $`D9\overline{D9}`$-brane system respecting only the $`USp(4k)\times USp(4k)\times SO(32)`$ subgroup of the largest possible group with the same number of CP indexes. Second, we consider the symplectic groups as factors because we are looking for a minimal construction which leads to $`k`$ type I $`D5`$-branes. As is established in ref. each of the branes should have two CP indexes, taking values in $`USp(2)SU(2)`$. Hence, $`k`$ type I $`D5`$-branes correspond to $`USp(2k)`$ group . Third, we consider $`4k`$ rather than just $`2k`$ indexes, because we have to embed two of the CP indexes of both $`USp(4k)`$’s into the tangent bundle of the target space (see for such a construction).
Now if the tachyon is covariantly constant in the background in question it contains the aforementioned map. Hence, the tachyon field is a rectangular matrix ($`[4k]\times [32+4k]`$) whose quadratic part ($`[4k]\times [4k]`$) is the map in question. This is exactly the tachyon value we expect to get for $`k`$ type I $`D5`$-branes to appear as remnants of the $`D9`$-brane annihilation .
In conclusion, we have a $`D9\overline{D9}`$-system with some excited Ramond-Ramond (RR) field corresponding to $`k`$ $`D5`$-branes before the annihilation. The RR field is encoded in terms of some gauge bundle on the $`D9`$-$`\overline{D9}`$-system . After the annihilation, the tachyon acquires a value which is covariantly constant on light-like surfaces in the gauge field background. In our case, such a tachyon value is proportional to the ADHM matrix , which corresponds to the ADHM construction of instantons for the $`SO(32)`$ group. Similarly as in ref. , this matrix defines a mass term for the fermionic fields on the probe $`D1`$-brane.
In this way, we obtain the probe theory which, on the level of massless modes, coincides with the non-linear $`\sigma `$-model for the ADHM construction . In other words, it flows in the IR limit to the same superconformal theory which describes the $`D5`$-brane background for the type I $`D1`$-brane theory as an instanton field of the $`SO(32)`$ group . This time it is the latter field which encodes the information about corresponding RR charge. The former gauge field from the diagonal subgroup of $`USp(4k)\times USp(4k)`$, being a pure gauge at low energies (due to the non-zero tachyon vacuum expectation value (VEV)), decouples after the IR limit is taken.
Thus, without knowing the tachyon potential we could fix tachyon values after the annihilation. This is our main result.
## 2 Twistor formulation of the probe theory and $`\kappa `$-invariance
We consider a phase of type I string theory containing $`32+4k`$ $`D9`$-branes and $`4k`$ $`\overline{D9}`$-branes. The D-brane world-volumes fill the entire ten-dimensional space-time. We probe the annihilation of the $`D9\overline{D9}`$-system by a $`D1`$-brane.
In the GS or twistor formalism the probe theory contains the following fields at low energies. First, there are low-energy modes of strings attached by both their ends to the $`D1`$-brane. The modes are ten bosons $`x_M,(M=0,..,9)`$ and ten-dimensional Majorana-Weyl fermions $`\psi _𝒜,(𝒜=1,\mathrm{},16)`$ . Second, there are low-energy modes of strings stretched between the $`D1`$\- and $`D9`$-branes. These modes are two-dimensional Majorana-Weyl fermions $`\lambda `$ . Third, there are also modes of strings stretched between the $`D1`$\- and $`\overline{D9}`$-branes. Correspondingly these modes are two-dimensional Majorana-Weyl fermions $`\chi `$ of opposite to $`\lambda `$ chirality . If it were not for the presence of $`\chi `$, the $`D1`$-brane would have the same quantum numbers as the Heterotic $`SO(32)`$ string .
We are going to work with the twistor formulation of the theory :
$`S={\displaystyle }d^2\sigma \{P_M^{}[e_a^{++}(^ax^M^a\psi ^𝒜\mathrm{\Gamma }_𝒜^M\psi ^{})\phi _{}\mathrm{\Gamma }^M\phi _{}]+`$
$`+\mathrm{Wess}\mathrm{Zumino}\mathrm{term}+`$
$`+\lambda ^p\left[e_{}^a\left(\delta ^{pq}_a_ax^MA_M^{pq}\left(x\right)\right)+{\displaystyle \frac{1}{4}}F_{ML}^{pq}\left(x\right)\mathrm{\Gamma }_𝒜^{ML}\psi ^𝒜\psi ^{}\right]\lambda ^q+`$
$`+\chi ^{\overline{p}}[e_{++}^a(\delta ^{\overline{p}\overline{q}}_a_ax^MB_M^{\overline{p}\overline{q}}\left(x\right))+{\displaystyle \frac{1}{4}}H_{ML}^{\overline{p}\overline{q}}\left(x\right)\mathrm{\Gamma }_𝒜^{ML}\psi ^𝒜\psi ^{}]\chi ^{\overline{q}}+\lambda ^p\chi ^{\overline{p}}T^{p\overline{p}}\left(x\right)\}.`$ (1)
Here $`P_M^a`$ and $`\phi _{}`$ are auxiliary fields. Their exact definition is not relevant for our further discussion and can be found in . These auxiliary fields should obey a Cartan-Penrose condition: $`P_{}^M=e_{}^4\phi _{}\mathrm{\Gamma }^M\phi _{}`$. Now it is easy to see how after the integration over $`P^M`$ one recovers the standard GS formulation of the Heterotic string if $`\chi `$ is absent .
Also in this formula $`e^a,a=1,2`$ is a zweibein; $`p=1,\mathrm{},32+4k`$ and $`\overline{p}=1,\mathrm{},4k`$; $`\mathrm{\Gamma }_M`$ ($`\mathrm{\Gamma }^{ML}=[\mathrm{\Gamma }^M,\mathrm{\Gamma }^L]`$) are ten-dimensional $`\gamma `$-matrices in the Majorana-Weyl representation; $`T`$ is a tachyon field which describes the lowest energy excitations of strings stretched between the $`D9`$\- and $`\overline{D9}`$-branes . It appears in (1) as an external field and transforms in the bi-fundamental representation under the $`USp(4k)\times SO(32)`$ and $`USp(4k)`$ groups. We choose the gauge fields $`A_M^{pq}`$ and $`B_M^{\overline{p}\overline{q}}`$ on the $`D9`$\- and $`\overline{D9}`$-branes (with the field strengths $`F_{ML}^{pq}`$ and $`H_{ML}^{\overline{p}\overline{q}}`$, correspondingly) respecting only this subgroup of the largest possible group with this number of CP indexes. These gauge fields couple to $`\lambda `$ and $`\chi `$ in the same way as a gauge field couples to Heterotic fermions . All other fields on the $`D9`$\- and $`\overline{D9}`$-branes are set to zero.
Hence, the theory we are starting with is a non-SUSY two-dimensional $`\sigma `$-model. It evolves via the RG flow to a superconformal theory in the IR if a proper background value of $`T`$ is standing in (1) . In fact, some of the $`D9`$\- and $`\overline{D9}`$-branes should annihilate leaving only the $`D1`$-brane in type I string theory, which contains only 32 $`D9`$-branes and, possibly, some non-trivial bundles on the latter. This $`D1`$-brane has the quantum numbers of the Heterotic string and its theory is superconformal. Thus, one can be sure that if the theory in question eventually evolved to a superconformal limit for some value of $`T`$, this value really corresponds to a minimum of the tachyon energy functional.
As we explained in the introduction, the $`\kappa `$-invariance could help find tachyon classical solutions respecting SUSY. Let us, hence, impose conditions on the invariance of the action (1) under $`\kappa `$-transformations. In the conformal gauge the transformations look as follows :
$`\delta P_M^a=0,\delta \phi _{}=0`$
$`\delta \psi ^𝒜=2iP_{}^M\mathrm{\Gamma }_M^𝒜\kappa _{++}`$
$`\delta x^M=i\delta \psi ^𝒜\mathrm{\Gamma }_𝒜^M\psi ^{}`$
$`\delta \left(A_M_{}x^M\right)=_{}\mathrm{\Lambda }_\kappa +[\mathrm{\Lambda }_\kappa ,A_M_{}x^M],\delta \lambda ^p=\left(\mathrm{\Lambda }_\kappa \lambda \right)^p`$ (2)
where $`\mathrm{\Lambda }_\kappa =\delta x^MA_M`$ if only the background gauge field is non-zero. Also, by analogy with the transformations of $`\lambda `$ we could choose a natural transformation law for $`\chi `$ to be:
$$\delta \left(B_M_{++}x^M\right)=_{++}\mathrm{\Lambda }_\kappa ^{}+[\mathrm{\Lambda }_\kappa ^{},B_M_{++}x^M],\delta \chi ^{\overline{p}}=\left(\mathrm{\Lambda }_\kappa ^{}\chi \right)^{\overline{p}},$$
(3)
and $`\mathrm{\Lambda }_\kappa ^{}=\delta x^MB_M`$. At the same time the tachyon field transforms under the $`\kappa `$-symmetry simply as follows:
$$\delta T\left(x\right)=_MT\left(x\right)\delta x^M$$
(4)
It is necessary to supplement these transformations by Virasoro and SYM constraints . In our case, the latter are equivalent to the classical SYM equations of motion. Moreover, for (1) to be invariant under (2) the tachyon field $`T`$ should obey the following equation:
$`\delta x^𝒜\widehat{𝒟}_𝒜T\left(x\right)=`$
$`=\delta x^𝒜\left\{\widehat{}_𝒜T^{q\overline{p}}\left(x\right)+T^{q\overline{q}}\left(x\right)\widehat{B}_𝒜^{\overline{q}\overline{p}}\left(x\right)\widehat{A}_𝒜^{qp}\left(x\right)T^{p\overline{p}}\left(x\right)\right\}=0`$
$`\mathrm{where}\delta x^𝒜=\widehat{P}_{}^{𝒜𝒞}\kappa _{𝒞++}\psi ^{}\mathrm{and}\widehat{P}_a^{𝒜𝒞}=P_a^M\mathrm{\Gamma }_M^{𝒜𝒞}`$ (5)
for any $`\kappa `$. Hence, this should be supplemented by integrability conditions:
$$[\widehat{P}\kappa _{(1)}\psi \widehat{𝒟},\widehat{P}\kappa _{(2)}\psi \widehat{𝒟}]=\delta x^{(1)}\delta x^{(2)}[𝒟^{(1)},𝒟^{(2)}]=0.$$
(6)
Note that the vector $`\widehat{P}\kappa \psi `$ has zero norm:$`(\widehat{P}\kappa \psi )^2P_{}^2ϵ_𝒜\psi ^𝒜\psi ^{}\delta ^{𝒞𝒟}\kappa _{𝒞++}\kappa _{𝒟++}=0`$ due to the anti-commutativity of $`\kappa `$. Hence, $`\widehat{P}\kappa \psi `$ is a constant, (independent of $`x_M`$) light-like vector in the ten-dimensional Minkowski space. Moreover, under variations of the parameter $`\kappa `$ (with $`P`$ and $`\psi `$ kept fixed) it sweeps an eight-dimensional hyperplane in the ten-dimensional space-time. In fact, as is well known, the matrix $`\widehat{P}\kappa `$ has eight rather than sixteen non-zero eigen-values . Hence, it defines eight real deformations of $`\kappa `$. Thus, there are eight varying components of the ten-vector in question, while the other two are fixed. This is the eight-dimensional hyperplane. At the same time under variations of $`\widehat{P}`$ and $`\psi `$ all eight-dimensional hyperplanes are swept.
## 3 Solutions to the constraints and $`D5`$-branes as remnants of the annihilation
As follows from (6) when the background $`A_M`$ and $`B_M`$ are zero, the tachyon field should be a constant up to a gauge transformation. Unfortunately we can not derive from our formulae what kind of constant it should be. However, as a warm up exercise let us try to guess it . Consider:
$$T_{[4k+32]\times [4k]}\left(D_{[4k]\times [4k]}0_{[32]\times [4k]}\right),$$
(7)
where $`D`$ is a diagonal matrix with all eigen-values of the order of the string scale. This tachyon value respects both $`SO(32)`$ and the diagonal $`USp(4k)`$ subgroups of the $`SO(32)\times USp(4k)\times USp(4k)`$ group. (See for discussion on this subject.) Gauge invariant expression for the tachyon VEV should be:
$$T^{\overline{p}p}T_p^{\overline{q}}=\delta ^{\overline{p}\overline{q}}.$$
(8)
Presumably the latter expression is related to the minimum of the tachyon potential: $`_TV(T)|_{T^2=1}=0`$, but its origin is not really important to us.
What is most important is that the formula (8) passes through the simplest check. In fact, consider the lowest energy excitations in the NS sector of the strings stretched between the $`D1`$-brane and the $`D9\overline{D9}`$-system . We denote these excitations as $`Q^p`$ and $`\stackrel{~}{Q}^{\overline{p}}`$, respectively. Their bare masses are equal to $`\frac{1}{2}`$ in string units . Also their interactions with the tachyon field are $`T^{p\overline{p}}T_{\overline{p}}^q\times Q_pQ_q`$ and $`T^{\overline{p}p}T_p^{\overline{q}}\times \stackrel{~}{Q}_{\overline{p}}\stackrel{~}{Q}_{\overline{q}}`$. Hence, when the tachyon acquires the VEV as in (7),(8) we have the proper number of the NS modes with mass $`\frac{1}{2}`$ to describe the $`D1`$-brane in the background of 32 $`D9`$-branes only. We are going to use eq. (8) later.
With the tachyon as in the eq. (7), the $`\chi ^{\overline{p}}`$ and $`\lambda ^{\overline{p}}`$ from (1) become massive, while $`\lambda ^n`$, $`n=1,\mathrm{},32`$ are left massless: here $`\lambda ^p=(\lambda ^{\overline{p}},\lambda ^n)`$. Because of IR effects in two dimensions, if such fields acquire masses there is no way for them to become massless. Hence, in the study of the RG evolution of the theory we can safely integrate these massive fields out, while leaving massless ones untouched. After the integration, we get the ordinary type I $`D1`$-brane theory. This theory describes the background of only 32 $`D9`$-branes and is superconformal . Thus, (7) is a proper VEV for the tachyon in this case. We just guessed the VEV in question but in the situation below we derive it.
Now let us study the case when there is a co-dimension four soliton left within the $`D9`$-branes after the annihilation. This corresponds to a solution of eq. (6) when $`A_M`$, $`B_M`$ and $`T`$ are functions of four coordinates. Say the latter are $`x_6,\mathrm{},x_9`$. In this case we expect a linear realization of SUSY in the IR limit<sup>2</sup><sup>2</sup>2We suppose, but can not prove, that in all other situations eq. (5) and (6) have trivial solutions..
In the presence of the soliton in question, the ten-dimensional Lorentz invariance is broken: $`SO(9,1)SO(1,1)\times SO(4)\times SO(4)`$. Here $`SO(1,1)`$ corresponds to rotations along the $`D1`$-brane (directions<sup>3</sup><sup>3</sup>3From now on we fix the light-cone gauge. $`0,1`$); one of the $`SO(4)`$’s is related to rotations along the soliton, but transverse to the $`D1`$-brane (directions $`2,\mathrm{},5`$), while another $`SO(4)`$ corresponds to rotations in the directions transversal to the soliton ($`6,\mathrm{},9`$).
Below we denote by $`\pm `$ the left and right chirality under the $`SO(1,1)`$ group. At the same time by $`\alpha `$ and $`\dot{\alpha }`$ we denote the indexes of the fundamental representation of the $`SU(2)_L`$ and $`SU(2)_R`$ subgroups of the aforementioned “transversal” $`SO(4)`$ group. Besides that we embed two among the CP indexes $`p`$ and $`\overline{p}`$ into the tangent bundle of the target space. Thus, $`p=(\alpha ,i,n)`$ and $`\overline{p}=(\dot{\alpha },i)`$, where $`n=1,\mathrm{},32`$ and $`i=1,\mathrm{},2k`$: hence, both $`USp(4k)`$’s are broken to $`SU(2)\times USp(2k)`$. In other words, the fermions on the $`D1`$-brane carry the following indexes $`\chi ^{\overline{p}}=\chi _{}^{\dot{\alpha }i}`$ and $`\lambda ^p=(\lambda _+^{\alpha i},\lambda _+^n)`$.
If $`A_M`$, $`B_M`$ and $`T`$ are functions of the four coordinates only, then (5) and (6) take the form:
$$d_{(1),(2)}^{\alpha \dot{\alpha }}𝒟_{\alpha \dot{\alpha }}T\left(x\right)=0\mathrm{and}d_{(1)}^{\alpha \dot{\alpha }}d_{(2)}^{\beta \dot{\beta }}F_{\alpha \dot{\alpha }\beta \dot{\beta }}\left(x\right)=0,$$
(9)
where $`F_{\mu \nu }=[𝒟_\mu ,𝒟_\nu ],F_{\alpha \dot{\alpha }\beta \dot{\beta }}=F_{\mu \nu }\tau _{\alpha \dot{\alpha }}^\mu \tau _{\beta \dot{\beta }}^\nu `$, $`\mu =6,\mathrm{},9`$. Also here $`d_{\alpha \dot{\alpha }}^{(1),(2)}`$ are constant (independent of $`x_\mu `$) vectors with zero norm in the complexified four-dimensional Euclidean space. They correspond to $`\widehat{P}\kappa \psi `$ with two different $`\kappa `$’s.
We consider complexification of the Euclidean space (use complex $`x_{\alpha \dot{\alpha }}`$ coordinates rather than real $`x_\mu `$) and complexify the gauge group to show that there are non-trivial solutions to the eq. (9). What is most important, if $`\widehat{P}`$ and $`\psi `$ are fixed, the vector $`d_{\alpha \dot{\alpha }}`$ sweeps a complex two-plane ($`\beta `$-plane in the notation of ref. ) in the complex four-dimensional space: In the case under study $`\widehat{P}\kappa `$ describes two complex deformations. At the same time, under variations of $`\widehat{P}`$ and $`\psi `$ the vector $`d_{\alpha \dot{\alpha }}`$ sweeps all light-like surfaces in the space.
Before going further we would like to remind that we are looking for a minimal construction of the $`D5`$-branes out of the $`D9`$-$`\overline{D9}`$-system. As we mentioned in the introduction, to construct $`k`$ type I $`D5`$-branes as a result of a $`D9`$-brane annihilation we need $`4k+32`$ $`D9`$-branes and $`4k`$ $`\overline{D9}`$-branes. Besides that we take $`A_\mu ^{nm}`$ to be a pure gauge. Then, the second equation in (9) is equivalent to the self-duality condition for the YM connection $`\stackrel{~}{A}_\mu ^{\overline{p}\overline{q}}=A_\mu ^{\overline{p}\overline{q}}B_\mu ^{\overline{p}\overline{q}}`$ from the diagonal subgroup of $`USp(4k)\times USp(4k)`$.
To find IR limit of the theory (1) we need to know the large distance behavior of the instanton solution to eq. (9). With the charge $`2k`$ the solution behaves at infinity as follows:
$`\stackrel{~}{A}_\mu ^{\overline{p}\overline{q}}\left(\left(_\mu \widehat{S}\right)\widehat{S}^1\right)_{\alpha \beta }^{ij}\mathrm{where}S_{\alpha \dot{\alpha }}^{ij}=\delta ^{ij}{\displaystyle \frac{\left(\widehat{x}_{\dot{\alpha }\alpha }\widehat{x}_{\dot{\alpha }\alpha }^{(i)}\right)}{\left|xx^{(i)}\right|}},|\widehat{x}|\mathrm{}.`$ (10)
Here $`\widehat{x}_{\alpha \dot{\alpha }}=x_\mu \tau _{\alpha \dot{\alpha }}^\mu `$ and $`x_{(i)}`$ are positions of the $`2k`$ instantons. This is not the most general behavior at infinity but we use it to clarify our idea.
The gauge field (10) defines a map $`\widehat{S}`$ of the order $`2k`$ from $`S^3USp(2)SU(2)`$ at space infinity to the diagonal subgroup of $`USp(4k)\times USp(4k)`$. In fact:
$$\left(\left(_\mu \widehat{S}\right)\widehat{S}^1\right)_{\alpha \beta }^{ij}=2\sigma _e^{\alpha \beta }\eta _{\mu \nu }^e\delta ^{ij}\frac{\left(xx^{(i)}\right)_\nu }{\left|xx^{(i)}\right|^2},e=1,2,3,$$
(11)
where $`\eta _{\mu \nu }^e`$ are t’Hooft symbols and $`\sigma _e`$ are generators of the $`SU(2)`$ group. This is just a singular instanton of the $`USp(4k)`$ group with the charge $`2k`$.
Plugging (10) into (9), we find that the tachyon field behaves at large distances as:
$$T_{p\dot{\alpha }}^j\left\{\frac{\left(\widehat{x}_{\dot{\alpha }\alpha }\widehat{x}_{\dot{\alpha }\alpha }^{(i)}\right)}{\left|xx^{(i)}\right|}\delta ^{ij}0_{\dot{\alpha }}^{in}\right\},|\widehat{x}|\mathrm{}.$$
(12)
This is the tachyon value found in . It describes $`k`$ singular $`D5`$-branes within type I string theory after the annihilation.
Now let us consider the more general situation, i.e. deform the gauge field (10) to:
$`\stackrel{~}{A}_\mu ^{\overline{p}\overline{q}}\left(\left(_\mu \widehat{S}\right)\widehat{S}^1\right)_{\alpha \beta }^{ij},\mathrm{where}S_{\alpha \dot{\alpha }}^{ij}=\mathrm{\Delta }^{il}\left(x\right)\left(\widehat{x}_{\dot{\alpha }\alpha }\delta ^{lj}\widehat{X}_{\dot{\alpha }\alpha }^{lj}\right),|\widehat{x}|\mathrm{}`$
$`\mathrm{and}\left(\mathrm{\Delta }^2\right)^{ij}=\left\{\left(x\delta ^{il}X^{il}\right)_\mu \left(x\delta ^{lj}X^{lj}\right)_\nu [X_\mu ,X_\nu ]^{ij}\right\}\tau ^\mu \tau ^\nu .`$ (13)
Here $`\widehat{S}`$ defines a most general (up to gauge transformations) map of the order $`2k`$ from $`S^3`$ at spatial infinity to the diagonal subgroup of $`USp(4k)\times USp(4k)`$. In this formula $`\widehat{X}_{\alpha \dot{\alpha }}^{ij}`$ is an arbitrary symplectic matrix from the diagonal subgroup of $`USp(2k)\times USp(2k)`$, which obeys the reality condition $`\widehat{X}^{\alpha \dot{\alpha }}=ϵ^{\alpha \beta }ϵ^{\dot{\alpha }\dot{\beta }}\widehat{X}_{\beta \dot{\beta }}^{}`$.
With this value of the gauge field substituted into eq. (9), we find that the tachyon behaves at infinity as :
$$T_{p\dot{\alpha }}^j\mathrm{\Delta }^{il}\left\{\left(\widehat{x}_{\dot{\alpha }\alpha }\delta ^{lj}\widehat{X}_{\dot{\alpha }\alpha }^{lj}\right)h_{\dot{\alpha }}^{ln}\right\},\mathrm{where}h^{\dot{\alpha }jn}=ϵ^{ij}ϵ^{\dot{\alpha }\dot{\beta }}\left(h_{i\dot{\beta }}^n\right)^{},|\widehat{x}|\mathrm{}.$$
(14)
It is a solution to eq. (9) up to the gauge transformation by the matrix $`\mathrm{\Delta }^{ij}`$ from the diagonal subgroup of $`USp(2k)\times USp(2k)`$. Here $`ϵ^{ij}`$ is $`USp(2k)`$ invariant tensor served to raise and lower $`i`$ indexes.
So far $`h_{\dot{\alpha }}^{in}`$ in eq. (14) is an arbitrary matrix, i.e. not fixed by the eq. (9). However, taking into account eq. (8), or $`T_{p\dot{\alpha }}^jT^{i\dot{\beta }p}\delta ^{ij}\delta _{\dot{\alpha }}^{\dot{\beta }}`$ in our case, the matrices $`X`$ and $`h`$ should obey the ADHM condition :
$`\left(ϵ^{\alpha \beta }\widehat{X}_{\alpha \dot{\alpha }}\widehat{X}^{\beta \dot{\beta }}\right)^{ij}+h_{\dot{\alpha }}^{in}h^{nj\dot{\beta }}=0,`$ (15)
with such a gauge choice as in eq. (14). Note that when $`h=0`$ we recover the situation of the singular $`D5`$-brane (12).
Let us now check whether or not we have found a proper tachyon value which minimizes its energy functional. Specifically we are going to check whether or not theory (1) flows to a superconformal limit in the IR with such a tachyon value.
Substituting the value (14) for the tachyon field into the action (1) we get:
$`=_{kin}(x,\psi ,\chi ,\lambda )+\chi _{}^{\alpha j}\mathrm{\Delta }^{jl}\left(\left(\widehat{x}_{\dot{\alpha }\alpha }\delta ^{li}\widehat{X}_{\dot{\alpha }\alpha }^{li}\right)\lambda _+^{\dot{\alpha }i}+h_\alpha ^{ln}\lambda _+^n\right).`$ (16)
Here we showed spinor indexes to present a close similarity of our theory to that considered in ref. . To understand how the theory (16) evolves under the RG flow, one must find massless fields among $`\lambda `$ and $`\chi `$. For this purpose it is necessary to look for a complete set of solutions to the equation :
$$T_{p\alpha }^j(x)v^{pn}(x)=0,$$
(17)
for a general $`X^{ij}`$ and $`h_{\dot{\alpha }}^{in}`$ obeying (15). Once we have found all $`32`$ solutions to these equations, it is possible to decompose $`\lambda _+`$ in their basis:
$$\lambda _+^p=\underset{n=1}{\overset{32}{}}v^{pn}\lambda _+^n.$$
(18)
Substituting this expression into (16) and integrating out massive modes, we get:
$$=(x,\psi )_{kin}+\lambda _+^n\left(_{}\delta ^{nm}+_{}\widehat{x}^{\dot{\alpha }\alpha }\widehat{A}_{\dot{\alpha }\alpha }^{nm}(x)\right)\lambda _+^m,$$
(19)
where $`\widehat{A}_{\dot{\alpha }\alpha }^{nm}(x)=\left(v_p^n\right)^1\frac{}{x^{\dot{\alpha }\alpha }}v_p^m`$. Note also that the gauge field (13), being a pure gauge at low energies (due to the non-zero tachyon VEV (8)), does not enter the IR Lagrangian (19). The theory (19) is superconformal .
Taking into account (14),(15) and (17) we see that $`\widehat{A}_{\alpha \dot{\alpha }}^{nm}`$ is the self-dual vector-potential corresponding to the ADHM “matrices” $`X`$ and $`h`$. This vector-potential describes $`k`$ type I $`D5`$-branes in some non-singular (but otherwise generic) point of their moduli space . So, choosing the value of $`T`$ as in (14) and (15), we arrive via RG at the superconformal theory (19). It describes the type I $`D1`$-brane in the background of $`k`$ $`D5`$-branes.
Now let us discuss a difference between our theory and that considered in . One immediately sees that the massive spectra of the two theories are different . First, because of the factor $`\mathrm{\Delta }^{ij}`$, masses of the massive fermions in (16) are different from those of fermions in the theory from ref. . Second, scalars which are present in and correspond (along with $`\chi _{}^{\alpha j}`$ and $`\lambda _+^{\dot{\alpha }j}`$) to the strings stretched between $`D1`$\- and $`D5`$-branes, are absent in eq. (16). These scalars are massive at a generic point of the instanton moduli space .
It is worth mentioning at this point that one should not expect the theory (16) to properly reproduce all massive modes. In fact, this theory is a low-energy one, because it depends on classical values of macroscopic fields such as $`T`$ and $`A`$ . Hence, the theory does not contain microscopic degrees of freedom.
## 4 Conclusions and Acknowledgments
Thus, the theories from eq. (16) and from ref. , while being different at high energies, flow to the same superconformal field theory in the IR. We believe in the existence of a microscopic theory underlying both of the theories in question . After the integration of one type of its massive modes, the microscopic theory should lead to the theory considered in as an intermediate step of the RG flow. However, the integration of another type of its massive modes leads the microscopic theory to the Lagrangian (16) as an intermediate step of the RG flow.
I would like to acknowledge discussions with D.Sorokin, D.Polyakov, A.Sen, N.Berkovits, E.Sezgin, A.Gorsky, K.Zarembo, A.Rosly and especially with A.Gerasimov. Also I would like to thank P.Horava for encouraging me to deal with the subject in question. This work was done under the support of NSERC NATO fellowship grant and under the partial support of grants INTAS-97-01-03 and RFBR 98-02-16575.
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# Prompt 𝐽/𝜓 Polarization at the Tevatron
## Acknowledgments
The author thanks Eric Braaten and Bernd A. Kniehl for their enjoyable collaboration on the subject discussed here. This work was supported in part by the Alexander von Humboldt Foundation.
## References
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# Reasoning with Axioms: Theory and PracticeThis paper appeared in the Proceedings of the Seventh International Conference on Priciples of Knowledge Representation and Reasoning (KR’2000).
## 1 MOTIVATION
Description Logics (DLs) form a family of formalisms which have grown out of knowledge representation techniques using frames and semantic networks. DLs use a class based paradigm, describing the domain of interest in terms of concepts (classes) and roles (binary relations) which can be combined using a range of operators to form more complex structured concepts \[BHH<sup>+</sup>91\]. A DL *terminology* typically consists of a set of asserted facts, in particular asserted subsumption (is-a-kind-of) relationships between (possibly complex) concepts.<sup>1</sup><sup>1</sup>1DLs can also deal with assertions about individuals, but in this paper we will only be concerned with *terminological* (concept based) reasoning.
One of the distinguishing characteristics of DLs is a formally defined semantics which allows the structured objects they describe to be reasoned with. Of particular interest is the computation of implied subsumption relationships between concepts, based on the assertions in the terminology, and the maintenance of a concept hierarchy (partial ordering) based on the subsumption relationship \[WS92\].
The problem of computing concept subsumption relationships has been the subject of much research, and sound and complete algorithms are now known for a wide range of DLs (for example \[HN90, BH91, Baa91, DMar, HST99\]). However, in spite of the fundamental importance of terminologies in DLs, most of these algorithms deal only with the problem of deciding subsumption between two concepts (or, equivalently, concept satisfiability), without reference to a terminology (but see \[BDS93, Cal96, DDM96, HST99\]). By restricting the kinds of assertion that can appear in a terminology, concepts can be syntactically expanded so as to explicitly include all relevant terminological information. This procedure, called *unfolding*, has mostly been applied to less expressive DLs. With more expressive DLs, in particular those supporting universal roles, it is often possible to encapsulate an arbitrary terminology in a single concept. This technique can be used with satisfiability testing to ensure that the result is valid with respect to the assertions in the terminology, a procedure called *internalisation*.
Although the above mentioned techniques suffice to demonstrate the theoretical adequacy of satisfiability decision procedures for terminological reasoning, experiments with implementations have shown that, for reasons of (lack of) efficiency, they are highly unsatisfactory as a practical methodology for reasoning with DL terminologies. Firstly, experiments with the Kris system have shown that integrating unfolding with the (tableaux) satisfiability algorithm (*lazy unfolding*) leads to a significant improvement in performance \[BFH<sup>+</sup>94\]. More recently, experiments with the FaCT system have shown that reasoning becomes hopelessly intractable when internalisation is used to deal with larger terminologies \[Hor98\]. However, the FaCT system has also demonstrated that this problem can be dealt with (at least for realistic terminologies) by using a combination of lazy unfolding and internalisation, having first manipulated the terminology in order to minimise the number of assertions that must be dealt with by internalisation (a technique called *absorption*).
It should be noted that, although these techniques were discovered while developing DL systems, they are applicable to a whole range of reasoning systems, independent of the concrete logic and type of algorithm. As well as tableaux based decision procedures, this includes resolution based algorithms, where the importance of minimising the number of terminological sentences has already been noted \[HS99\], and sequent calculus algorithms, where there is a direct correspondence with tableaux algorithms \[BFH<sup>+</sup>99\].
In this paper we seek to improve our theoretical understanding of these important techniques which has, until now, been very limited. In particular we would like to know exactly when and how they can be applied, and be sure that the answers we get from the algorithm are still correct. This is achieved by defining a formal framework that allows the techniques to be precisely described, establishing conditions under which they can be safely applied, and proving that, provided these conditions are respected, satisfiability algorithms will still function correctly. These results are then used to show that the procedures used in the FaCT system are correct<sup>2</sup><sup>2</sup>2Previously, the correctness of these procedures had only been demonstrated by a relatively ad-hoc argument \[Hor97\]. and, moreover, to show how efficiency can be significantly improved, while still retaining the guarantee of correctness, by relaxing the safety conditions for absorption. Finally, we identify several interesting directions for future research, in particular the problem of finding the “best” absorption possible.
## 2 PRELIMINARIES
Firstly, we will establish some basic definitions that clarify what we mean by a DL, a terminology (subsequently called a TBox), and subsumption and satisfiability with respect to a terminology, . The results in this paper are uniformly applicable to a whole range of DLs, as long as some basic criteria are met:
###### Definition 2.1 (Description Logic)
Let $`𝖫`$ be a DL based on infinite sets of atomic concepts $`\mathrm{𝖭𝖢}`$ and atomic roles $`\mathrm{𝖭𝖱}`$. We will identify $`𝖫`$ with the sets of its well-formed concepts and require $`𝖫`$ to be closed under boolean operations and sub-concepts.
An interpretation is a pair $`=(\mathrm{\Delta }^{},^{})`$, where $`\mathrm{\Delta }^{}`$ is a non-empty set, called the *domain* of $``$, and $`^{}`$ is a function mapping $`\mathrm{𝖭𝖢}`$ to $`2^\mathrm{\Delta }^{}`$ and $`\mathrm{𝖭𝖱}`$ to $`2^{\mathrm{\Delta }^{}\times \mathrm{\Delta }^{}}`$. With each DL $`𝖫`$ we associate a set $`\mathrm{𝖨𝗇𝗍}(𝖫)`$ of *admissible* interpretations for $`𝖫`$. $`\mathrm{𝖨𝗇𝗍}(𝖫)`$ must be closed under isomorphisms, and, for any two interpretations $``$ and $`^{}`$ that agree on $`\mathrm{𝖭𝖱}`$, it must satisfy $`\mathrm{𝖨𝗇𝗍}(𝖫)^{}\mathrm{𝖨𝗇𝗍}(𝖫)`$. Additionally, we assume that each DL $`𝖫`$ comes with a semantics that allows any interpretation $`\mathrm{𝖨𝗇𝗍}(𝖫)`$ to be extended to each concept $`C𝖫`$ such that it satisfies the following conditions:
* it maps the boolean combination of concepts to the corresponding boolean combination of their interpretations, and
* the interpretation $`C^{}`$of a compound concept $`C𝖫`$ depends only on the interpretation of those atomic concepts and roles that appear syntactically in $`C`$.
This definition captures a whole range of DLs, namely, the important DL $`𝒜𝒞`$ \[SS91\] and its many extensions. $`\mathrm{𝖨𝗇𝗍}(𝖫)`$ hides restrictions on the interpretation of certain roles like transitivity, functionality, or role hierarchies, which are imposed by more expressive DLs (e.g., \[HST99\]), as these are irrelevant for our purposes. In these cases, $`\mathrm{𝖨𝗇𝗍}(𝖫)`$ will only contain those interpretations which interpret the roles as required by the semantics of the logic, e.g., features by partial functions or transitively closed roles by transitive relations. Please note that various modal logics \[Sch91\], propositional dynamic logics \[DL94\] and temporal logics \[EH85\] also fit into this framework. We will use $`CD`$ as an abbreviation for $`\neg CD`$, $`CD`$ as an abbreviation for $`(CD)(DC)`$, and $``$ as a tautological concept, e.g., $`A\neg A`$ for an arbitrary $`A\mathrm{𝖭𝖢}`$.
A TBox consists of a set of axioms asserting subsumption or equality relations between (possibly complex) concepts.
###### Definition 2.2 (TBox, Satisfiability)
A *TBox* $`𝒯`$ for $`𝖫`$ is a finite set of axioms of the form $`C_1C_2`$ or $`C_1C_2`$, where $`C_i𝖫`$. If, for some $`A\mathrm{𝖭𝖢}`$, $`𝒯`$ contains one or more axioms of the form $`AC`$ or $`AC`$, then we say that $`A`$ is *defined* in $`𝒯`$.
Let $`𝖫`$ be a DL and $`𝒯`$ a TBox. An interpretation $`\mathrm{𝖨𝗇𝗍}(𝖫)`$ is a *model* of $`𝒯`$ iff, for each $`C_1C_2𝒯`$, $`C_1^{}C_2^{}`$ holds, and, for each $`C_1C_2𝒯`$, $`C_1^{}=C_2^{}`$ holds. In this case we write $`𝒯`$. A concept $`C𝖫`$ is *satisfiable* with respect to a TBox $`𝒯`$ iff there is an $`\mathrm{𝖨𝗇𝗍}(𝖫)`$ with $`𝒯`$ and $`C^{}\mathrm{}`$. A concept $`C𝖫`$ *subsumes* a concept $`D𝖫`$ w.r.t. $`𝒯`$ iff, for all $`\mathrm{𝖨𝗇𝗍}(𝖫)`$ with $`𝒯`$, $`C^{}D^{}`$ holds.
Two TBoxes $`𝒯,𝒯^{}`$ are called *equivalent* ($`𝒯𝒯^{})`$, iff, for all $`\mathrm{𝖨𝗇𝗍}(𝖫)`$, $`𝒯𝗂𝖿𝖿𝒯^{}`$.
We will only deal with concept satisfiability as concept subsumption can be reduced to it for DLs that are closed under boolean operations: $`C`$ subsumes $`D`$ w.r.t. $`𝒯`$ iff $`(D\neg C)`$ is not satisfiable w.r.t. $`𝒯`$.
For temporal or modal logics, satisfiability with respect to a set of formulae $`\{C_1,\mathrm{},C_k\}`$ asserted to be universally true corresponds to satisfiability w.r.t. the TBox $`\{C_1,\mathrm{},C_n\}`$.
Many decision procedures for DLs base their judgement on the existence of models or pseudo-models for concepts. A central rôle in these algorithms is played by a structure that we will call a *witness* in this paper. It generalises the notions of *tableaux* that appear in DL tableau-algorithms \[HNS90, BBH96, HST99\] as well as the *Hintikka-structures* that are used in tableau and automata-based decision procedures for temporal logic \[EH85\] and propositional dynamic logic \[VW86\].
###### Definition 2.3 (Witness)
Let $`𝖫`$ be a DL and $`C𝖫`$ a concept. A *witness* $`𝒲=(\mathrm{\Delta }^𝒲,^𝒲,^𝒲)`$ for $`C`$ consists of a non-empty set $`\mathrm{\Delta }^𝒲`$, a function $`^𝒲`$ that maps $`\mathrm{𝖭𝖱}`$ to $`2^{\mathrm{\Delta }^𝒲\times \mathrm{\Delta }^𝒲}`$, and a function $`^𝒲`$ that maps $`\mathrm{\Delta }^𝒲`$ to $`2^𝖫`$ such that the following properties are satisfied:
* there is some $`x\mathrm{\Delta }^𝒲`$ with $`C^𝒲(x)`$,
* there is an interpretation $`\mathrm{𝖨𝗇𝗍}(𝖫)`$ that *stems* from $`𝒲`$, and
* for each interpretation $`\mathrm{𝖨𝗇𝗍}(𝖫)`$ that *stems* from $`𝒲`$, it holds that $`D^𝒲(x)`$ implies $`xD^{}`$.
An interpretation $`=(\mathrm{\Delta }^{},^{})`$ is said to stem from $`𝒲`$ if it satisfies:
1. $`\mathrm{\Delta }^{}=\mathrm{\Delta }^𝒲`$,
2. $`^{}|_{\mathrm{𝖭𝖱}}=^𝒲`$, and
3. for each $`A\mathrm{𝖭𝖢}`$, $`A^𝒲(x)xA^{}`$ and $`\neg A^𝒲(x)xA^{}`$.
A witness $`𝒲`$ is called *admissible* with respect to a TBox $`𝒯`$ if there is an interpretation $`\mathrm{𝖨𝗇𝗍}(𝖫)`$ that stems from $`𝒲`$ with $`𝒯`$.
Please note that, for any witness $`𝒲`$, (W2) together with Condition 3 of “stemming” implies that, there exists no $`x\mathrm{\Delta }^𝒲`$ and $`A\mathrm{𝖭𝖢}`$, such that $`\{A,\neg A\}^𝒲(x)`$. Also note that, in general, more than one interpretation may stem from a witness. This is the case if, for an atomic concept $`A\mathrm{𝖭𝖢}`$ and an element $`x\mathrm{\Delta }^𝒲`$, $`^𝒲(x)\{A,\neg A\}=\mathrm{}`$ holds (because two interpretations $``$ and $`^{}`$, with $`xA^{}`$ and $`x\neg A^{^{}}`$, could both stem from $`𝒲`$).
Obviously, each interpretation $``$ gives rise to a special witness, called the *canonical witness*:
###### Definition 2.4 (Canonical Witness)
Let $`𝖫`$ be a DL. For any interpretation $`\mathrm{𝖨𝗇𝗍}(𝖫)`$ we define the *canonical witness* $`𝒲_{}=(\mathrm{\Delta }^𝒲_{},^𝒲_{},^𝒲_{})`$ as follows:
$`\mathrm{\Delta }^𝒲_{}`$ $`=\mathrm{\Delta }^{}`$
$`^𝒲_{}`$ $`=^{}|_{\mathrm{𝖭𝖱}}`$
$`^𝒲_{}`$ $`=\lambda x.\{D𝖫xD^{}\}`$
The following elementary properties of a canonical witness will be useful in our considerations.
###### Lemma 2.5
Let $`𝖫`$ be a DL, $`C𝖫`$, and $`𝒯`$ a TBox. For each $`\mathrm{𝖨𝗇𝗍}(𝖫)`$ with $`C^{}\mathrm{}`$,
1. each interpretation $`^{}`$ stemming from $`𝒲_{}`$ is isomorphic to $``$
2. $`𝒲_{}`$ is a witness for $`C`$,
3. $`𝒲_{}`$ is admissible w.r.t. $`𝒯`$ iff $`𝒯`$
###### Proof.
1. Let $`^{}`$ stem from $`𝒲_{}`$. This implies $`\mathrm{\Delta }^{^{}}=\mathrm{\Delta }^{}`$ and $`^{^{}}|_{\mathrm{𝖭𝖱}}=^{}|_{\mathrm{𝖭𝖱}}`$. For each $`x\mathrm{\Delta }^{}`$ and $`A\mathrm{𝖭𝖢}`$, $`\{A,\neg A\}^𝒲_{}(x)\mathrm{}`$, this implies $`^{^{}}|_{\mathrm{𝖭𝖢}}=^{}|_{\mathrm{𝖭𝖢}}`$ and hence $``$ and $`^{}`$ are isomorphic.
2. Properties (W1) and (W2) hold by construction. Obviously, $``$ stems from $`𝒲_{}`$ and from (1) it follows that each interpretation $`^{}`$ stemming from $`𝒲_{}`$ is isomorphic to $``$, hence (W3) holds.
3. Since $``$ stems from $`𝒲_{}`$, $`𝒯`$ implies that $`𝒲_{}`$ is admissible w.r.t. $`𝒯`$. If $`𝒲_{}`$ is admissible w.r.t. $`𝒯`$, then there is an interpretation $`^{}`$ stemming from $`𝒲_{}`$ with $`^{}𝒯`$. Since $``$ is isomorphic to $`^{}`$, this implies $`𝒯`$. ∎
As a corollary we get that the existence of admissible witnesses is closely related to the satisfiability of concepts w.r.t. TBoxes:
###### Lemma 2.6
Let $`𝖫`$ be a DL. A concept $`C𝖫`$ is satisfiable w.r.t. a TBox $`𝒯`$ iff it has a witness that is admissible w.r.t. $`𝒯`$.
###### Proof.
For the *only if*-direction let $`\mathrm{𝖨𝗇𝗍}(𝖫)`$ be an interpretation with $`𝒯`$ and $`C^{}\mathrm{}`$. From Lemma 2.5 it follows that the canonical witness $`𝒲_{}`$ is a witness for $`C`$ that is admissible w.r.t. $`𝒯`$.
For the *if*-direction let $`𝒲`$ be an witness for $`C`$ that is admissible w.r.t. $`𝒯`$. This implies that there is an interpretation $`\mathrm{𝖨𝗇𝗍}(𝖫)`$ stemming from $`𝒲`$ with $`𝒯`$. For each interpretation $``$ that stems from $`𝒲`$, it holds that $`C^{}\mathrm{}`$ due to (W1) and (W3). ∎∎
From this it follows that one can test the satisfiability of a concept w.r.t. to a TBox by checking for the existence of an admissible witness. We call algorithms that utilise this approach *model-building algorithms*.
This notion captures tableau-based decision procedures, \[HNS90, BBH96, HST99\], those using automata-theoretic approaches \[VW86, CDL99\] and, due to their direct correspondence with tableaux algorithms \[HS99, BFH<sup>+</sup>99\], even resolution based and sequent calculus algorithms.
The way many decision procedures for DLs deal with TBoxes exploits the following simple lemma.
###### Lemma 2.7
Let $`𝖫`$ be a DL, $`C𝖫`$ a concept, and $`𝒯`$ a TBox. Let $`𝒲`$ be a witness for $`C`$. If
$$\begin{array}{ccc}C_1C_2𝒯\hfill & & x\mathrm{\Delta }^𝒲.(C_1C_2^𝒲(x))\hfill \\ C_1C_2𝒯\hfill & & x\mathrm{\Delta }^𝒲.(C_1C_2^𝒲(x))\hfill \end{array}$$
then $`𝒲`$ is admissible w.r.t. $`𝒯`$.
###### Proof.
$`𝒲`$ is a witness, hence there is an interpretation $`\mathrm{𝖨𝗇𝗍}(𝖫)`$ stemming from $`𝒲`$. From (W3) and the fact that $`𝒲`$ satisfies the properties stated in 2.7 it follows that, for each $`x\mathrm{\Delta }^{}`$,
$$\begin{array}{ccc}C_1C_2𝒯\hfill & & C_1C_2^𝒲(x)\hfill \\ & & x(C_1C_2)^{}\hfill \\ C_1C_2𝒯\hfill & & C_1C_2^𝒲(x)\hfill \\ & & x(C_1C_2)^{}\hfill \end{array}$$
Hence, $`𝒯`$ and $`𝒲`$ is admissible w.r.t. $`𝒯`$. ∎∎
Examples of algorithms that exploit this lemma to deal with axioms can be found in \[DDM96, DL96, HST99\], where, for each axiom $`C_1C_2`$ ($`C_1C_2`$) the concept $`C_1C_2`$ ($`C_1C_2`$) is added to every node of the generated tableau.
Dealing with general axioms in this manner is costly due to the high degree of nondeterminism introduced. This can best be understood by looking at tableaux algorithms, which try to build witnesses in an incremental fashion. For a concept $`C`$ to be tested for satisfiability, they start with $`\mathrm{\Delta }^𝒲=\{x_0\}`$, $`^𝒲(x_0)=\{C\}`$ and $`^𝒲(R)=\mathrm{}`$ for each $`R\mathrm{𝖭𝖱}`$. Subsequently, the concepts in $`^𝒲`$ are decomposed and, if necessary, new nodes are added to $`\mathrm{\Delta }^𝒲`$, until either $`𝒲`$ is a witness for $`C`$, or an obvious contradiction of the form $`\{A,\neg A\}^𝒲(x)`$, which violates (W2), is generated. In the latter case, backtracking search is used to explore alternative non-deterministic decompositions (e.g., of disjunctions), one of which could lead to the discovery of a witness.
When applying Lemma 2.7, disjunctions are added to the label of each node of the tableau for each general axiom in the TBox (one disjunction for axioms of the form $`C_1C_2`$, two for axioms of the form $`C_1C_2`$). This leads to an exponential increase in the search space as the number of nodes and axioms increases. For example, with 10 nodes and a TBox containing 10 general axioms (of the form $`C_1C_2`$) there are already 100 disjunctions, and they can be non-deterministically decomposed in $`2^{100}`$ different ways. For a TBox containing large numbers of general axioms (there are 1,214 in the Galen medical terminology KB \[RNG93\]), this can degrade performance to the extent that subsumption testing is effectively non-terminating. To reason with this kind of TBox we must find a more efficient way to deal with axioms.
## 3 ABSORPTIONS
We start our considerations with an analysis of a technique that can be used to deal more efficiently with so-called primitive or acyclic TBoxes.
###### Definition 3.1 (Absorption)
Let $`𝖫`$ be a DL and $`𝒯`$ a TBox. An *absorption* of $`𝒯`$ is a pair of TBoxes $`(𝒯_u,𝒯_g)`$ such that $`𝒯𝒯_u𝒯_g`$ and $`𝒯_u`$ contains only axioms of the form $`AD`$ and $`\neg AD`$ where $`A\mathrm{𝖭𝖢}`$.
An absorption $`(𝒯_u,𝒯_g)`$ of $`𝒯`$ is called *correct* if it satisfies the following condition. For each witness $`𝒲`$, if, for each $`x\mathrm{\Delta }^𝒲`$,
$$\begin{array}{ccc}\hfill AD𝒯_uA^𝒲(x)& & D^𝒲(x)\hfill \\ \hfill \neg AD𝒯_u\neg A^𝒲(x)& & D^𝒲(x)\hfill \\ \hfill C_1C_2𝒯_g& & C_1C_2^𝒲(x)\hfill \\ \hfill C_1C_2𝒯_g& & C_1C_2^𝒲(x)\hfill \end{array}$$
then $`𝒲`$ is admissible w.r.t. $`𝒯`$. We refer to this properties by $`()`$. A witness that satisfies $`()`$ will be called *unfolded w.r.t. $`𝒯`$*.
If the reference to a specific TBox is clear from the context, we will often leave the TBox implicit and say that a witness is unfolded.
How does a correct absorption enable an algorithm to deal with axioms more efficiently? This is best described by returning to tableaux algorithms. Instead of dealing with axioms as previously described, which may lead to an exponential increase in the search space, axioms in $`𝒯_u`$ can now be dealt with in a deterministic manner. Assume, for example, that we have to handle the axiom $`AC`$. If the label of a node already contains $`A`$ (resp. $`\neg A`$), then $`C`$ (resp. $`\neg C`$) is added to the label; if the label contains neither $`A`$ nor $`\neg A`$, then *nothing* has to be done. Dealing with the axioms in $`𝒯_u`$ this way avoids the necessity for additional non-deterministic choices and leads to a gain in efficiency. A witness produced in this manner will be unfolded and is a certificate for satisfiability w.r.t. $`𝒯`$. This technique is generally known as *lazy unfolding* of primitive TBoxes \[Hor98\]; formally, it is justified by the following lemma:
###### Lemma 3.2
Let $`(𝒯_u,𝒯_g)`$ be a correct absorption of $`𝒯`$. For any $`C𝖫`$, $`C`$ has a witness that is admissible w.r.t. $`𝒯`$ iff $`C`$ has an unfolded witness.
###### Proof.
The *if*-direction follows from the definition of “correct absorption”. For the *only if*-direction, let $`C𝖫`$ be a concept and $`𝒲`$ a witness for $`C`$ that is admissible w.r.t. $`𝒯`$. This implies the existence of an interpretation $`\mathrm{𝖨𝗇𝗍}(𝖫)`$ stemming from $`𝒲`$ such that $`𝒯`$ and $`C^{}\mathrm{}`$. Since $`𝒯𝒯_u𝒯_g`$ we have $`𝒯_u𝒯_g`$ and hence the canonical witness $`𝒲_{}`$ is an unfolded witness for $`C`$. ∎∎
A family of TBoxes where absorption can successfully be applied are *primitive* TBoxes, the most simple form of TBox usually studied in the literature.
###### Definition 3.3 (Primitive TBox)
A TBox $`𝒯`$ is called *primitive* iff it consists entirely of axioms of the form $`AD`$ with $`A\mathrm{𝖭𝖢}`$, each $`A\mathrm{𝖭𝖢}`$ appears as at most one left-hand side of an axiom, and $`𝒯`$ is acyclic. Acyclicity is defined as follows: $`A\mathrm{𝖭𝖢}`$ is said to *directly use* $`B\mathrm{𝖭𝖢}`$ if $`AD𝒯`$ and $`B`$ occurs in $`D`$; *uses* is the transitive closure of “directly uses”. We say that $`𝒯`$ is *acyclic* if there is no $`A\mathrm{𝖭𝖢}`$ that uses itself.
For primitive TBoxes a correct absorption can easily be given.
###### Theorem 3.4
Let $`𝒯`$ be a primitive TBox, $`𝒯_g=\mathrm{}`$, and $`𝒯_u`$ defined by
$$𝒯_u=\{AD,\neg A\neg DAD𝒯\}.$$
Then $`(𝒯_u,𝒯_g)`$ is a correct absorption of $`𝒯`$.
###### Proof.
Trivially, $`𝒯𝒯_u𝒯_g`$ holds. Given an unfolded witness $`𝒲`$, we have to show that there is an interpretation $``$ stemming from $`𝒲`$ with $`𝒯`$.
We fix an arbitrary linearisation $`A_1,\mathrm{},A_k`$ of the “uses” partial order on the atomic concept names appearing on the left-hand sides of axioms in $`𝒯`$ such that, if $`A_i`$ uses $`A_j`$, then $`j<i`$ and the defining concept for $`A_i`$ is $`D_i`$.
For some interpretation $``$, atomic concept $`A`$, and set $`X\mathrm{\Delta }^{}`$, we denote the interpretation that maps $`A`$ to $`X`$ and agrees with $``$ on all other atomic concepts and roles by $`[AX]`$. For $`0ik`$, we define $`_i`$ in an iterative process starting from an arbitrary interpretation $`_0`$ stemming from $`𝒲`$ and setting
$$_i:=_{i1}[A_i\{x\mathrm{\Delta }^𝒲xD_i^{_{i1}}\}]$$
Since, for each $`A_i`$ there is exactly one axiom in $`𝒯`$, each step in this process is well-defined. Also, since $`\mathrm{𝖨𝗇𝗍}(𝖫)`$ may only restrict the interpretation of atomic roles, $`_i\mathrm{𝖨𝗇𝗍}(𝖫)`$ for each $`0ik`$. For $`=_k`$ it can be shown that $``$ is an interpretation stemming from $`𝒲`$ with $`𝒯`$.
First we prove inductively that, for $`0ik`$, $`_i`$ stems from $`𝒲`$. We have already required $`_0`$ to stem from $`𝒲`$.
Assume the claim was proved for $`_{i1}`$ and $`_i`$ does not stem from $`𝒲`$. Then there must be some $`x\mathrm{\Delta }^𝒲`$ such that either (i) $`A_i^𝒲(x)`$ but $`xA_i^_i`$ or (ii) $`\neg A_i^𝒲(x)`$ but $`xA_i^_i`$ (since we assume $`_{i1}`$ to stem from $`𝒲`$ and $`A_i`$ is the only atomic concept whose interpretation changes from $`_{i1}`$ to $`_i`$). The two cases can be handled dually:
* From $`A_i^𝒲(x)`$ it follows that $`D_i^𝒲(x)`$, because $`𝒲`$ is unfolded. Since $`_{i1}`$ stems from $`𝒲`$ and $`𝒲`$ is a witness, Property (W3) implies $`xD_i^{_{i1}}`$. But this implies $`xA_i^_i`$, which is a contradiction.
* From $`\neg A_i^𝒲(x)`$ it follows that $`\neg D_i^𝒲(x)`$ because $`𝒲`$ is unfolded. Since $`_{i1}`$ stems from $`𝒲`$ and $`𝒲`$ is an witness, Property (W3) implies $`x(\neg D_i)^{_{i1}}`$. Since $`(\neg D_i)^{_{i1}}=\mathrm{\Delta }^𝒲D_i^{_{i1}}`$ this implies $`xA_i^_i`$, which is a contradiction.
Together this implies that $`_i`$ also stems from $`𝒲`$.
To show that $`𝒯`$ we show inductively that $`_iA_jD_j`$ for each $`1ji`$. This is obviously true for $`i=0`$.
The interpretation of $`D_i`$ may not depend on the interpretation of $`A_i`$ because otherwise (I2) would imply that $`A_i`$ uses itself. Hence $`D_i^_i=D_i^{_{i1}}`$ and, by construction, $`_iA_iD_i`$. Assume there is some $`j<i`$ such that $`_i\vDash ̸A_jD_j`$. Since $`_{i1}A_jD_j`$ and only the interpretation of $`A_i`$ has changed from $`_{i1}`$ to $`_i`$, $`D_j^_iD_j^{_{i1}}`$ must hold because of (I2). But this implies that $`A_i`$ occurs in $`D_j`$ and hence $`A_j`$ uses $`A_i`$ which contradicts $`j<i`$. Thus, we have $`A_j=D_j`$ for each $`1jk`$ and hence $`𝒯`$. ∎∎
Lazy unfolding is a well-known and widely used technique for optimising reasoning w.r.t. primitive TBoxes \[BFH<sup>+</sup>94\]. So far, we have only given a correctness proof for this relatively simple approach, although one that is independent of a specific DL or reasoning algorithm. With the next lemma we show how we can extend correct absorptions and hence how lazy unfolding can be applied to a broader class of TBoxes. A further enhancement of the technique is presented in Section 5.
###### Lemma 3.5
Let $`(𝒯_u,𝒯_g)`$ be a correct absorption of a TBox $`𝒯`$.
1. If $`𝒯^{}`$ is an arbitrary TBox, then $`(𝒯_u,𝒯_g𝒯^{})`$ is a correct absorption of $`𝒯𝒯^{}`$.
2. If $`𝒯^{}`$ is a TBox that consists entirely of axioms of the form $`AD`$, where $`A\mathrm{𝖭𝖢}`$ and $`A`$ is not defined in $`𝒯_u`$, then $`(𝒯_u𝒯^{},𝒯_g)`$ is a correct absorption of $`𝒯𝒯^{}`$.
###### Proof.
In both cases, $`𝒯_u𝒯_g𝒯^{}𝒯𝒯^{}`$ holds trivially.
1. Let $`C𝖫`$ be a concept and $`𝒲`$ be an unfolded witness for $`C`$ w.r.t. the absorption $`(𝒯_u,𝒯_g𝒯^{})`$. This implies that $`𝒲`$ is unfolded w.r.t. the (smaller) absorption $`(𝒯_u,𝒯_g)`$. Since $`(𝒯_u,𝒯_g)`$ is a correct absorption, there is an interpretation $``$ stemming from $`𝒲`$ with $`𝒯`$. Assume $`\vDash ̸𝒯^{}`$. Then, without loss of generality,<sup>3</sup><sup>3</sup>3Arbitrary TBoxes can be expressed using only axioms of the form $`CD`$. there is an axiom $`DE𝒯^{}`$ such that there exists an $`xD^{}E^{}`$. Since $`𝒲`$ is unfolded, we have $`DE^𝒲(x)`$ and hence (W3) implies $`x(\neg DE)^{}=\mathrm{\Delta }^{}(D^{}E^{})`$, a contradiction. Hence $`𝒯𝒯^{}`$ and $`𝒲`$ is admissible w.r.t. $`𝒯𝒯^{}`$.
2. Let $`C𝖫`$ be a concept and $`𝒲`$ be an unfolded witness for $`C`$ w.r.t. the absorption $`(𝒯_u𝒯^{},𝒯_g)`$. From $`𝒲`$ we define a new witness $`𝒲^{}`$ for $`C`$ by setting $`\mathrm{\Delta }^𝒲^{}:=\mathrm{\Delta }^𝒲`$, $`^𝒲^{}:=^𝒲`$, and definig $`^𝒲^{}`$ to be the function that, for every $`x\mathrm{\Delta }^𝒲^{}`$, maps $`x`$ to the set
$$^𝒲(x)\{\neg AAD𝒯^{},A^𝒲(x)\}$$
It is easy to see that $`𝒲^{}`$ is indeed a witness for $`C`$ and that $`𝒲^{}`$ is also unfolded w.r.t. the absorption $`(𝒯_u𝒯^{},𝒯_g)`$. This implies that $`𝒲^{}`$ is also unfolded w.r.t. the (smaller) absorption $`(𝒯_u,𝒯_g)`$. Since $`(𝒯_u,𝒯_g)`$ is a correct absorption of $`𝒯`$, there exists an interpretation $``$ stemming from $`𝒲^{}`$ such that $`𝒯`$. We will show that $`𝒯^{}`$ also holds. Assume $`\vDash ̸𝒯^{}`$, then there is an axiom $`AD𝒯^{}`$ and an $`x\mathrm{\Delta }^{}`$ such that $`xA^{}`$ but $`xD^{}`$. By construction of $`𝒲^{}`$, $`xA^{}`$ implies $`A^𝒲^{}(x)`$ because otherwise $`\neg A^𝒲^{}(x)`$ would hold in contradiction to (W3). Then, since $`𝒲^{}`$ is unfolded, $`D^𝒲^{}(x)`$, which, again by (W3), implies $`xD^{}`$, a contradiction.
Hence, we have shown that there exists an interpretation $``$ stemming from $`𝒲^{}`$ such that $`𝒯_u𝒯^{}𝒯_g`$. By construction of $`𝒲^{}`$, any interpretation stemming from $`𝒲^{}`$ also stems from $`𝒲`$, hence $`𝒲`$ is admissible w.r.t. $`𝒯𝒯^{}`$. ∎
## 4 APPLICATION TO FaCT
In the preceeding section we have defined correct absorptions and discussed how they can be exploited in order to optimise satisfiability procedures. However, we have said nothing about the problem of how to find an absorption given an arbitrary terminology. In this section we will describe the absorption algorithm used by FaCT and prove that it generates correct absorptions.
Given a TBox $`𝒯`$ containing arbitrary axioms, the absorption algorithm used by FaCT constructs a triple of TBoxes $`(𝒯_g,𝒯_{\text{prim}},𝒯_{\text{inc}})`$ such that
* $`𝒯𝒯_g𝒯_{\text{prim}}𝒯_{\text{inc}}`$,
* $`𝒯_{\text{prim}}`$ is primitive, and
* $`𝒯_{\text{inc}}`$ consists only of axioms of the form $`AD`$ where $`A\mathrm{𝖭𝖢}`$ and $`A`$ is not defined in $`𝒯_{\text{prim}}`$.
We refer to these properties by $`()`$. From Theorem 3.4 together with Lemma 3.5 it follows that, for
$$𝒯_u:=\{AD,\neg A\neg DAD𝒯_{\text{prim}}\}𝒯_{\text{inc}}$$
($`𝒯_u`$,$`𝒯_g`$) is a correct absorption of $`𝒯`$; hence satisfiability for a concept $`C`$ w.r.t. $`𝒯`$ can be decided by checking for an unfolded witness for $`C`$.
In a first step, FaCT distributes axioms from $`𝒯`$ amongst $`𝒯_{\text{inc}}`$, $`𝒯_{\text{prim}}`$, and $`𝒯_g`$, trying to minimise the number of axioms in $`𝒯_g`$ while still maintaining $`()`$. To do this, it initialises $`𝒯_{\text{prim}},𝒯_{\text{inc}}`$, and $`𝒯_g`$ with $`\mathrm{}`$, and then processes each axiom $`X𝒯`$ as follows.
1. If $`X`$ is of the form $`AC`$, then
1. if $`A\mathrm{𝖭𝖢}`$ and $`A`$ is not defined in $`𝒯_{\text{prim}}`$ then $`X`$ is added to $`𝒯_{\text{inc}}`$,
2. otherwise $`X`$ is added to $`𝒯_g`$
2. If $`X`$ is of the form $`AC`$, then
1. if $`A\mathrm{𝖭𝖢}`$, $`A`$ is not defined in $`𝒯_{\text{prim}}`$ or $`𝒯_{\text{inc}}`$ and $`𝒯_{\text{prim}}\{X\}`$ is primitive, then $`X`$ is added to $`𝒯_{\text{prim}}`$,
2. otherwise, the axioms $`AC`$ and $`CA`$ are added to $`𝒯_g`$.
3. If $`X`$ is of the form $`CD`$, then add $`X`$ to $`𝒯_g`$
4. If $`X`$ is of the form $`CD`$, then add $`CD,DC`$ to $`𝒯_g`$.
It is easy to see that the resulting TBoxes $`𝒯_g,𝒯_{\text{prim}},𝒯_{\text{inc}}`$ satisfy $`()`$. In a second step, FaCT processes the axioms in $`𝒯_g`$ one at a time, trying to absorb them into axioms in $`𝒯_{\text{inc}}`$. Those axioms that are not absorbed remain in $`𝒯_g`$. To give a simpler formulation of the algorithm, each axiom $`(CD)𝒯_g`$ is viewed as a clause $`𝐆=\{D,\neg C\}`$, corresponding to the axiom $`CD`$, which is equivalent to $`CD`$. For each such axiom FaCT applies the following absorption procedure.
1. Try to absorb $`𝐆`$. If there is a concept $`\neg A𝐆`$ such that $`A\mathrm{𝖭𝖢}`$ and $`A`$ is not defined in $`𝒯_{\text{prim}}`$, then add $`AB`$ to $`𝒯_{\text{inc}}`$, where $`B`$ is the disjunction of all the concepts in $`𝐆\{\neg A\}`$, remove $`𝐆`$ from $`𝒯_g`$, and exit.
2. Try to simplify $`𝐆`$.
1. If there is some $`\neg C𝐆`$ such that $`C`$ is of the form $`C_1\mathrm{}C_n`$, then substitute $`\neg C`$ with $`\neg C_1\mathrm{}\neg C_n`$, and continue with step 2b.
2. If there is some $`C𝐆`$ such that $`C`$ is of the form $`(C_1\mathrm{}C_n)`$, then apply associativity by setting $`𝐆=𝐆\{C_1,\mathrm{},C_n\}\{(C_1\mathrm{}C_n)\}`$, and return to step 1.
3. Try to unfold $`𝐆`$. If, for some $`A𝐆`$ (resp. $`\neg A𝐆`$), there is an axiom $`AC`$ in $`𝒯_{\text{prim}}`$, then substitute $`A𝐆`$ (resp. $`\neg A𝐆`$) with $`C`$ (resp. $`\neg C`$) and return to step 1.
4. If none of the above were possible, then absorption of $`𝐆`$ has failed. Leave $`𝐆`$ in $`𝒯_g`$, and exit.
For each step, we have to show that $`()`$ is maintained. Dealing with clauses instead of axioms causes no problems. In the first step, axioms are moved from $`𝒯_g`$ to $`𝒯_{\text{inc}}`$ as long as this does not violate $`()`$. The second and the third step replace a clause by an equivalent one and hence do not violate $`()`$.
Termination of the procedure is obvious. Each axiom is considered only once and, for a given axiom, simplification and unfolding can only be applied finitely often before the procedure is exited, either by absorbing the axiom into $`𝒯_{\text{inc}}`$ or leaving it in $`𝒯_g`$. For simplification, this is obvious; for unfolding, this holds because $`𝒯_{\text{prim}}`$ is acyclic. Hence, we get the following:
###### Theorem 4.1
For any TBox $`𝒯`$, FaCT computes a correct absorption of $`𝒯`$.
## 5 IMPROVING PERFORMANCE
The absorption algorithm employed by FaCT already leads to a dramatic improvement in performance. This is illustrated by Figure 1, which shows the times taken by FaCT to classify versions of the Galen KB with some or all of the general axioms removed. Without absorption, classification time increased rapidly with the number of general axioms, and exceeded 10,000s with only 25 general axioms in the KB; with absorption, only 160s was taken to classify the KB with all 1,214 general axioms.
However, there is still considerable scope for further gains. In particular, the following definition for a *stratified* TBox allows lazy unfolding to be more generally applied, while still allowing for correct absorptions.
###### Definition 5.1 (Stratified TBox)
A TBox $`𝒯`$ is called *stratified* iff it consists entirely of axioms of the form $`AD`$ with $`A\mathrm{𝖭𝖢}`$, each $`A\mathrm{𝖭𝖢}`$ appears at most once on the left-hand side of an axiom, and $`𝒯`$ can be arranged monotonously, i.e., there is a disjoint partition $`𝒯_1\dot{}𝒯_2\dot{}\mathrm{}\dot{}𝒯_k`$ of $`𝒯`$, such that
* for all $`1j<ik`$, if $`A\mathrm{𝖭𝖢}`$ is defined in $`𝒯_i`$, then it does not occur in $`𝒯_j`$, and
* for all $`1ik`$, all concepts which appear on the right-hand side of axioms in $`𝒯_i`$ are monotone in all atomic concepts defined in $`𝒯_i`$.
A concept $`C`$ is monotone in an atomic concept $`A`$ if, for any interpretation $`\mathrm{𝖨𝗇𝗍}(𝖫)`$ and any two sets $`X_1,X_2\mathrm{\Delta }^{}`$,
$$X_1X_2C^{[AX_1]}C^{[AX_2]}.$$
For many DLs, a sufficient condition for monotonicity is *syntactic* monotonicity, i.e., a concept $`C`$ is syntactically monotone in some atomic concept $`A`$ if $`A`$ does no appear in $`C`$ in the scope of an odd number of negations.
Obviously, due to its acyclicity, every primitive TBox is also stratified and hence the following theorem is a strict generalisation of Theorem 3.4.
###### Theorem 5.2
Let $`𝒯`$ be a stratified TBox, $`𝒯_g=\mathrm{}`$ and $`𝒯_u`$ defined by
$$𝒯_u=\{AD,\neg A\neg DAD𝒯\}.$$
Then $`(𝒯_u,𝒯_g)`$ is a correct absorption of $`𝒯`$.
The proof of this theorem follows the same line as the proof of Theorem 3.4. Starting from an arbitrary interpretation $`_0`$ stemming from the unfolded witness, we incrementally construct interpretations $`_1,\mathrm{},_k`$, using a fixed point construction in each step. We show that each $`_i`$ stems from $`𝒲`$ and that, for $`1j<ik`$, $`_i𝒯_j`$, hence $`_k𝒯`$ and stems from $`𝒲`$.
Before we prove this theorem, we recall some basics of lattice theory. For any set $`𝒮`$, the powerset of $`𝒮`$, denoted by $`2^𝒮`$ forms a complete lattice, where the ordering, join and meet operations are set-inclusion $``$, union $``$, and intersection $``$, respectively. For any complete lattice $``$, its $`n`$-fold cartesian product $`^n`$ is also a complete lattice, with ordering, join, and meet defined in a pointwise manner.
For a lattice $``$, a function $`\mathrm{\Phi }:`$ is called monotone, iff, for $`x_1,x_2`$, $`x_1x_2`$ implies $`\mathrm{\Phi }(x_1)\mathrm{\Phi }(x_2)`$.
By Tarski’s fixed point theorem \[Tar55\], every monotone function $`\mathrm{\Phi }`$ on a complete lattice, has uniquely defined least and greatest fixed points, i.e., there are elements $`\overline{x},\underset{¯}{x}`$ such that
$$\overline{x}=\mathrm{\Phi }(\overline{x})\text{ and }\underset{¯}{x}=\mathrm{\Phi }(\underset{¯}{x})$$
and, for all $`x`$ with $`x=\mathrm{\Phi }(x)`$,
$$\underset{¯}{x}x\text{ and }x\overline{x}.$$
Proof of Theorem 5.2. $`𝒯_u𝒯_g𝒯`$ is obvious. Let $`𝒲=(\mathrm{\Delta }^𝒲,^𝒲,^𝒲)`$ be an unfolded witness. We have to show that there is an interpretation $``$ stemming from $`𝒲`$ with $`𝒯`$. Let $`𝒯_1,\mathrm{},𝒯_k`$ be the required partition of $`𝒯`$. We will define $``$ inductively, starting with an arbitrary interpretation $`_0`$ stemming from $`𝒲`$.
Assume $`_{i1}`$ was already defined. We define $`_i`$ from $`_{i1}`$ as follows: let $`\{A_1^iD_1^i,\mathrm{},A_m^iD_m^i\}`$ be an enumeration of $`𝒯_i`$. First we need some auxiliary notation: for any concept $`C𝖫`$ we define
$$C^𝒲:=\{x\mathrm{\Delta }^𝒲C^𝒲(x)\}.$$
Using this notation we define the function $`\mathrm{\Phi }`$ mapping subsets $`X_1,\mathrm{},X_m`$ of $`\mathrm{\Delta }^𝒲`$ to
$$\begin{array}{ccc}(\hfill & ((A_1^i)^𝒲(D_1^i)^{_{i1}(X_1,\mathrm{},X_m)})(\neg A_1^i)^𝒲,\hfill & \\ & \mathrm{},\hfill & \\ & ((A_m^i)^𝒲(D_m^i)^{_{i1}(X_1,\mathrm{},X_m)})(\neg A_m^i)^𝒲\hfill & )\hfill \end{array}$$
where
$$_{i1}(X_1,\mathrm{},X_m):=_{i1}[A_1^iX_1,\mathrm{},A_m^iX_m]$$
Since all of the $`D_j^i`$ are monotone in all of the $`A_m^i`$, $`\mathrm{\Phi }`$ is a monontone function. This implies that $`\mathrm{\Phi }`$ has a least fixed point, which we denote by $`(\underset{¯}{X}_1,\mathrm{},\underset{¯}{X}_m)`$. We use this fixed point to define $`_i`$ by
$$_i:=_{i1}[A_1^i\underset{¯}{X}_1,\mathrm{},A_m^i\underset{¯}{X}_m]$$
Claim 1: For each $`0ik`$, $`_i`$ stems from $`𝒲`$.
We show this claim by induction on $`i`$. We have already required $`_0`$ to stem from $`𝒲`$. Assume $`_{i1}`$ stems from $`𝒲`$. Since the only thing that changes from $`_{i1}`$ to $`_i`$ is the interpretation of the atomic concepts $`A_1^i,\mathrm{},A_m^i`$, we only have to check that $`A_j^i^𝒲(x)`$ implies $`x(A_j^i)^_i`$ and $`\neg A_j^i^𝒲(x)`$ implies $`x(A_j^i)^_i`$.
By definition of $`\mathrm{\Phi }`$, and because $`\{xA_j^i^𝒲(x)\}\{x\neg A_j^i^𝒲(x)\}=\mathrm{}`$, $`A_j^i^𝒲(x)`$ implies $`x(A_j^i)^_i`$. Also by the definition of $`\mathrm{\Phi }`$, $`\neg A_j^i^𝒲(x)`$ implies $`x(A_j^i)^_i`$. Hence, $`_i`$ stems from $`𝒲`$.
Claim 2: For each $`1jik`$, $`_i𝒯_j`$.
We prove this claim by induction over $`i`$ starting from $`0`$. For $`i=0`$, there is nothing to prove. Assume the claim would hold for $`_{i1}`$. The only thing that changes from $`_{i1}`$ to $`_i`$ is the interpretation of the atomic concepts $`A_1^i,\mathrm{}A_m^i`$ defined in $`𝒯_i`$. Since these concepts may not occur in $`𝒯_j`$ for $`j<i`$, the interpretation of the concepts in these TBoxes does not change, and from $`_{i1}𝒯_j`$ follows $`_i𝒯_j`$ for $`1ji1`$.
It remains to show that $`_i𝒯_i`$. Let $`A_j^iD_j^i`$ be an axiom from $`𝒯_i`$. From the definition of $`_i`$ we have
$$(A_j^i)^_i=((A_j^i)^𝒲(D_j^i)^_i)(\neg A_j^i)^𝒲.$$
(1)
$`𝒲`$ is unfolded, hence $`A_j^i^𝒲(x)`$ implies $`D_j^i^𝒲(x)`$ and, since $`_i`$ stems from $`𝒲`$, this implies $`x(D_j^i)^_i`$, thus
$$(A_j^i)^𝒲(D_j^i)^_i=(D_j^i)^_i$$
(2)
Furthermore, $`\neg A_j^i^𝒲(x)`$ implies $`\neg D_j^i^𝒲(x)`$ implies $`x(\neg D_j^i)^_i`$, thus
$$(D_j^i)^_i(\neg A_j^i)^𝒲=(D_j^i)^_i$$
(3)
Taking together (1), (2), and (3) we get
$$(A_j^i)^_i=(D_j^i)^_i,$$
and hence $`_iA_j^iD_j^i`$.
Together, Claim 1 and Claim 2 prove the theorem, since $`_k`$ is an interpretation that stems from $`𝒲`$ and satisfies $`𝒯`$. ∎
This theorem makes it possible to apply the same lazy unfolding strategy as before to cyclical definitions. Such definitions are quite natural in a logic that supports inverse roles. For example, an orthopaedic procedure might be defined as a procedure performed by an orthopaedic surgeon, while an orthopaedic surgeon might be defined as a surgeon who performs only orthopaedic procedures:<sup>4</sup><sup>4</sup>4This example is only intended for didactic purposes.
$$\begin{array}{ccc}\hfill \text{o-procedure}& & \text{procedure}(\text{performs}^{}.\text{o-surgeon})\hfill \\ \hfill \text{o-surgeon}& & \text{surgeon}(\text{performs}.\text{o-procedure})\hfill \end{array}$$
The absorption algorithm described in Section 4 would force the second of these definitions to be added to $`𝒯_g`$ as two general axioms and, although both axioms would subsequently be absorbed into $`𝒯_u`$, the procedure would result in a disjunctive term being added to one of the definitions in $`𝒯_u`$. Using Theorem 5.2 to enhance the absorption algorithm so that these kinds of definition are directly added to $`𝒯_u`$ reduces the number of disjunctive terms in $`𝒯_u`$ and can lead to significant improvements in performance.
This can be demonstrated by a simple experiment with the new FaCT system, which implements the $`𝒮𝒬`$ logic \[HST99\] and is thus able to deal with inverse roles. Figure 2 shows the classification time in seconds using the normal and enhanced absorption algorithms for terminologies consisting of between 5 and 50 pairs of cyclical definitions like those described above for o-surgeon and o-procedure. With only 10 pairs the gain in performance is already a factor of 30, while for 45 and 50 pairs it has reached several orders of magnitude: with the enhanced absorption the terminology is classified in 2–3 seconds whereas with the original algorithm the time required exceeded the 10,000 second limit imposed in the experiment.
It is worth pointing out that it is by no means trivially true that cyclical definitions can be dealt with by lazy unfolding. Even without inverse roles it is clear that definitions such as $`A\neg A`$ (or more subtle variants) force the domain to be empty and would lead to an incorrect absorption if dealt with by lazy unfolding. With converse roles it is, for example, possible to force the interpretation of a role $`R`$ to be empty with a definition such as $`AR.(R^{}.\neg A)`$, again leading to an incorrect absorption if dealt with by lazy unfolding.
## 6 OPTIMAL ABSORPTIONS
We have demonstrated that absorption is a highly effective and widely applicable technique, and by formally defining correctness criteria for absorptions we have proved that the procedure used by FaCT finds correct absorptions. Moreover, by establishing more precise correctness criteria we have demonstrated how the effectiveness of this procedure could be further enhanced.
However, the absorption algorithm used by FaCT is clearly sub-optimal, in the sense that changes could be made that would, in general, allow more axioms to be absorbed (e.g., by also giving special consideration to axioms of the form $`\neg AC`$ with $`A\mathrm{𝖭𝖢}`$). Moreover, the procedure is non-deterministic, and, while it is guaranteed to produce a correct absorption, its specific result depends on the order of the axioms in the original TBox $`𝒯`$. Since the semantics of a TBox $`𝒯`$ does not depend on the order of its axioms, there is no reason to suppose that they will be arranged in a way that yields a “good” absorption. Given the effectiveness of absorption, it would be desirable to have an algorithm that was guaranteed to find the “best” absorption possible for any set of axioms, irrespective of their ordering in the TBox.
Unfortunately, it is not even clear how to define a sensible optimality criterion for absorptions. It is obvious that simplistic approaches based on the number or size of axioms remaining in $`𝒯_g`$ will not lead to a useful solution for this problem. Consider, for example, the cyclical TBox experiment from the previous section. Both the original FaCT absorption algorithm and the enhanced algorithm, which exploits Theorem 5.2, are able to compute a complete absorption of the axioms ( i.e., a correct absorption with $`𝒯_g=\mathrm{}`$), but the enhanced algorithm leads to much better performance, as shown in Figure 2.
An important issue for future work is, therefore, the identification of a suitable optimality criterion for absorptions, and the development of an algorithm that is able to compute absorptions that are optimal with respect to this criterion.
### Acknowledgements
This work was partially supported by the DFG, Project No. GR 1324/3-1.
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# From CKM Matrix to MNS Matrix: A Model Based on Supersymmetric 𝑆𝑂(10)×𝑈(2)_𝐹 Symmetry
## Abstract
We construct a realistic model based on SUSY $`SO(10)`$ with $`U(2)`$ flavor symmetry. In contrast to the commonly used effective operator approach, $`126`$dimensional Higgses are used to construct the Yukawa sector. R-parity symmetry is thus preserved at low energies. The Dirac and right-handed Majorana mass matrices in our model have very small mixing, and they combine with the seesaw mechanism resulting in a large leptonic mixing. The symmetric mass textures arising from the left-right symmetry breaking chain of $`SO(10)`$ give rise to very good predictions; 15 masses (including 3 right-handed Majorana neutrino masses) and 6 mixing angles are predicted by 11 parameters. Both the vacuum oscillation and LOW solutions are favored for the solar neutrino problem.
preprint: COLO-HEP-445preprint: May 2000
The flavor problem with hierarchical fermion masses and mixing has attracted a great deal of attention especially since the advent of the atmospheric neutrino oscillation data from Super-Kamiokande SuperK:1998a indicating non-zero neutrino masses. The non-zero neutrino masses give support to the idea of grand unification based on $`SO(10)`$ in which all the 16 fermions (including the right-handed neutrinos) can be accommodated in one single spinor representation. Furthermore, it provides a framework in which seesaw mechanism arises naturally. Models based on $`SO(10)`$ (and some with $`E_6`$) combined with a continuous or discrete flavor symmetry group have been constructed to understand the flavor problem. Most of the recent ones have used asymmetric or ”lopsided” mass textures to account for the maximal mixing in the neutrino sector. Symmetric mass textures have less parameters and hence could lead to more predictive power. Naively one expects, for symmetric mass textures, six texture zeros in the quark sector. But it has been observed by Ramond, Roberts and Ross Ramond:fmass1993a that the highest number of texture zeros has to be five, and using phenomenological analyses, they were able to arrive at five sets of up- and down-quark mass matrices with five texture zeros. Our analysis with recent experimental data and using CP conserving real symmetric matrices indicates that only one set (labeled set (v) in Ramond:fmass1993a ) remains viable (see below). The aim of this paper is to construct a realistic model based on $`SO(10)`$ combined with $`U(2)`$ as the flavor group, utilizing this set of symmetric mass textures for charged fermions. We first discuss the viable phenomenology of mass textures followed by the model which accounts for it, and then the implications of the model for neutrino mixing are presented.
Mass Texture Analysis: Throughout this paper we consider CP conserving real mass matrices. We do not lose any generality in our results since the CP violating phases do not have significant contributions to other parameters. A more detailed analysis taking into account CP violating phases will be given elsewhere chen:fmass2000a .
We consider the following mass textures at the GUT scale for the up-quark, down-quark and charged lepton sectors Ramond:fmass1993a ,
$`M_u=\left(\begin{array}{ccc}0& 0& a\\ 0& b& c\\ a& c& 1\end{array}\right)d,`$ $`M_d=\left(\begin{array}{ccc}0& e& 0\\ e& f& 0\\ 0& 0& 1\end{array}\right)h`$ (7)
$`M_e=\left(\begin{array}{ccc}0& e& 0\\ e& 3f& 0\\ 0& 0& 1\end{array}\right)h`$ (11)
with $`abc1`$ and $`ef1`$. After diagonalizing $`M^{}M`$, one obtains the following non-negative mass eigenvalues and mass ratios:
$`m_u{\displaystyle \frac{a^2bd}{b+c^2}},`$ $`m_c(b+c^2)d,`$ $`m_td`$
$`m_d{\displaystyle \frac{e^2fh}{e^2+f^2}},`$ $`m_s{\displaystyle \frac{(2e^2f+f^3)h}{e^2+f^2}},`$ $`m_b=h`$
$`m_e{\displaystyle \frac{3e^2fh}{e^2+9f^2}},`$ $`m_\mu {\displaystyle \frac{6e^2fh+27f^3h}{e^2+9f^2}},`$ $`m_\tau =h`$ (12)
$`{\displaystyle \frac{m_d}{m_e}}{\displaystyle \frac{e^2+9f^2}{3(e^2+f^2)}}3+O({\displaystyle \frac{e^2}{f^2}})`$
$`{\displaystyle \frac{m_s}{m_\mu }}{\displaystyle \frac{2e^4+19e^2f^2+9f^4}{6e^4+33e^2f^2+27f^4}}{\displaystyle \frac{1}{3}}+O({\displaystyle \frac{e^2}{f^2}})`$
$`{\displaystyle \frac{m_b}{m_\tau }}=1`$ (13)
These analytic expressions are very good approximations to the exact eigenvalues. It can be easily seen that the phenomenologically favored Georgi-Jarlskog relations Georgi:fmass1979a are obtained
$$m_d3m_e,m_s\frac{1}{3}m_\mu ,m_b=m_\tau $$
(14)
As we will see later, these relations between the down-quark sector and the charged lepton sector can be naturally achieved in $`SO(10)`$.
In order to explain the smallness of the neutrino masses, we will adopt the type I seesaw mechanism Gell-Mann:fmass1979a which requires both Dirac and right-handed Majorana mass matrices to be present in the Lagrangian. The right-handed Majorana mass matrix $`M_{\nu _{RR}}`$ is at present an unknown sector. The only constraint is that it must be constructed in such a way that it gives a favored low energy Majorana neutrino mass matrix $`M_{\nu _{LL}}`$ via the seesaw mechanism. We first consider the low energy (left-handed) Majorana neutrino mass matrix to get some insights into the structure of $`M_{\nu _{RR}}`$. We adopt the hierarchical scenario: $`|m_{\nu _3}||m_{\nu _2}|,|m_{\nu _1}|`$, to accommodate the experimental neutrino oscillation data. One way to achieve a large mixing in the $`\nu _\mu \nu _\tau `$ sector and at the same time a large mass splitting between $`m_{\nu _2}`$ and $`m_{\nu _3}`$ is to consider
$$M_{\nu _{LL}}\left(\begin{array}{ccc}0& 0& t\\ 0& 1& 1\\ t& 1& 1\end{array}\right)\mathrm{\Lambda }$$
(15)
A generic feature of the mass matrix of this type is that it leads to a large mixing in both $`\nu _e\nu _\mu `$ and $`\nu _\mu \nu _\tau `$ sectors (the so-called bimaximal mixing) for a broad range of $`t`$, $`0t1`$. In order to have a large mass splitting, we require $`t1`$. The three eigenvalues of this mass matrix keeping only the dominant orders are given by, in units of $`\mathrm{\Lambda }`$,
$`|m_1|`$ $``$ $`{\displaystyle \frac{t}{\sqrt{2}}}{\displaystyle \frac{t^2}{8}}{\displaystyle \frac{3t^3}{64\sqrt{2}}}`$
$`|m_2|`$ $``$ $`{\displaystyle \frac{t}{\sqrt{2}}}+{\displaystyle \frac{t^2}{8}}{\displaystyle \frac{3t^3}{64\sqrt{2}}}`$
$`|m_3|`$ $``$ $`2+{\displaystyle \frac{t^2}{4}}`$ (16)
The diagonalization matrix up to order $`O(t^2)`$ is given by
$$U_{\nu _{LL}}=\left(\begin{array}{ccc}\frac{1}{\sqrt{2}}\frac{t}{16}\frac{17t^2}{256\sqrt{2}}& \frac{1}{2}\frac{5t}{16\sqrt{2}}\frac{23t^2}{512}& \frac{1}{2}\frac{3t}{16\sqrt{2}}\frac{25t^2}{512}\\ \frac{1}{\sqrt{2}}\frac{t}{16}+\frac{17t^2}{256\sqrt{2}}& \frac{1}{2}+\frac{5t}{16\sqrt{2}}\frac{23t^2}{512}& \frac{1}{2}+\frac{3t}{16\sqrt{2}}\frac{25t^2}{512}\\ \frac{t}{2\sqrt{2}}& \frac{1}{\sqrt{2}}\frac{3t^2}{16\sqrt{2}}& \frac{1}{\sqrt{2}}+\frac{3t^2}{16\sqrt{2}}\end{array}\right)$$
(17)
Note that what the neutrino mixing matrix really means is the mismatch between the charged lepton flavor basis and the neutrino flavor basis analogous to the Cabbibo-Kobayashi-Maskawa quark mixing matrix $`V_{CKM}`$, and is the Maki-Nakagawa-Sakata matrix $`U_{MNS}`$ defined as,
$`U_{MNS}U_{e_L}U_{\nu _{LL}}^{}`$
$`=\left(\begin{array}{ccc}U_{e\nu _1}& U_{e\nu _2}& U_{e\nu _3}\\ U_{\mu \nu _1}& U_{\mu \nu _2}& U_{\mu \nu _3}\\ U_{\tau \nu _1}& U_{\tau \nu _2}& U_{\tau \nu _3}\end{array}\right)`$ (21)
Since the mixing matrix in the charged lepton sector $`U_{e_L}`$ is almost diagonal, combining $`U_{\nu _{LL}}^{}`$ and $`U_{e_L}`$ results in a nearly bimaximal mixing pattern in the lepton mixing matrix $`U_{MNS}`$. The squared mass difference between $`m_{\nu _1}^2`$ and $`m_{\nu _2}^2`$ is of the order of $`O(t^3)`$ while the squared mass difference between $`m_{\nu _2}^2`$ and $`m_{\nu _3}^2`$ is of the order $`O(1)`$. It is clear that the mass matrix eq.(15) naturally leads to the phenomenologically favored result
$$|\mathrm{\Delta }m_{23}^2||\mathrm{\Delta }m_{12}^2|$$
(22)
Depending on the value of $`t`$, both vacuum oscillation (VO) solution and large angle MSW (LAMSW) solution are possible. The VO solution suggests that $`\frac{\mathrm{\Delta }m_{}^2}{\mathrm{\Delta }m_{atm}^2}10^7`$. Since $`\frac{\mathrm{\Delta }m_{}^2}{\mathrm{\Delta }m_{atm}^2}t^3`$, one can see immediately that $`t10^3`$. On the other hand, the LAMSW solution suggests that $`\frac{\mathrm{\Delta }m_{}^2}{\mathrm{\Delta }m_{atm}^2}10^2`$, and $`t`$ is then required to be $`10^1`$. Since the element $`U_{e\nu _3}`$ in $`U_{MNS}`$ is proportional to $`t`$, an accurate measurement of $`U_{e\nu _3}`$ thus could provide some hints to single out one of the solar oscillation solutions if the neutrino mixing pattern is indeed bimaximal.
We assume that the Dirac neutrino mass matrix has the same texture (that is, positions of the zeros) as the up-quark mass matrix
$$M_{\nu _{LR}}=\left(\begin{array}{ccc}0& 0& \alpha \\ 0& \beta & \gamma \\ \alpha & \gamma & 1\end{array}\right)\eta $$
(23)
with $`\alpha \beta \gamma 1`$. We see later that $`M_u`$ and $`M_{\nu _{LR}}`$ can be in fact identical in $`SO(10)`$. To achieve $`M_{\nu _{LL}}`$ of the form of eq.(15) one needs a right-handed neutrino Majorana mass matrix of the same texture as $`M_{\nu _{LR}}`$
$$M_{\nu _{RR}}=\left(\begin{array}{ccc}0& 0& \delta _1\\ 0& \delta _2& \delta _3\\ \delta _1& \delta _3& 1\end{array}\right)M_R$$
(24)
with
$$\delta _1,\delta _2,\delta _31.$$
$`\delta _1{\displaystyle \frac{\alpha ^2}{2\alpha 2\alpha \gamma +\gamma ^2t}},`$ $`\delta _2{\displaystyle \frac{\beta ^2t}{2\alpha 2\alpha \gamma +\gamma ^2t}}`$
$`\delta _3{\displaystyle \frac{\alpha (\gamma \beta )+\beta \gamma t}{2\alpha 2\alpha \gamma +\gamma ^2t}}`$ (25)
After seesaw mechanism takes place,
$$M_{\nu _{LL}}=M_{\nu _{LR}}^TM_{\nu _{RR}}^1M_{\nu _{LR}}$$
(26)
$`M_{\nu _{LL}}`$ of the form eq.(15) results. It is interesting to see that the matrix operation in eq.(26) is form invariant. That is to say, $`M_{\nu _{LL}}`$ has the same texture as that of $`M_{\nu _{LR}}`$ and $`M_{\nu _{RR}}`$. Since $`\delta _1,\delta _2`$ and $`\delta _3`$ are much smaller than $`1`$, the mixing in the right-handed neutrino mass matrix $`M_{\nu _{RR}}`$ is generally small. Since the mixings in both $`M_{\nu _{LR}}`$ and $`M_{\nu _{RR}}`$ are small, our model falls into the category that the large neutrino mixing is purely due to the matrix operations in the seesaw mechanism given that charged lepton, Dirac neutrino and right-handed neutrino mixings are small, as classified in Ref.Barr:fmass2000a . We note that with the structure of $`M_{\nu _{LR}}`$ in eq.(23) we find it hard, though not impossible, to accommodate the small angle MSW solution in our model chen:fmass2000a .
We emphasize that, in the neutrino sector, we have been able to get a large mixing and at the same time a large mass splitting by using a symmetric Dirac mass matrix and a hierarchical right-handed mass matrix with very small mixing; the latter gives three superheavy hierarchical right-handed neutrino masses. Asymmetric Dirac mass matrices have been used before to get a large mixing and a large mass splitting Altarelli:1998nx ; Babu:1998wi ; Albright:1998vf ; Berezhiani:fmass1999a ; U1A .
$`𝐔(\mathrm{𝟐})`$ as a Flavor Symmetry: A prototype scenario which produces hierarchy in the fermion mass matrices is the Froggatt-Nielsen mechanism Froggatt:fmass1979a . It simply says that the heaviest matter fields acquire their masses through tree level interactions with the Higgs fields while masses of lighter matter fields are produced by higher dimensional interactions involving, in addition to the regular Higgs fields, exotic vector-like pairs of matter fields and the so-called flavons (flavor Higgs fields). After integrating out superheavy $`(M)`$ vector-like matter fields, the mass terms of the light matter fields get suppressed by a factor of $`\frac{<\theta >}{M}`$, where $`<\theta >`$ is the VEVs of the flavons and $`M`$ is the UV-cutoff of the effective theory above which the flavor symmetry is exact. We assume $`MM_{GUT}`$. We choose $`U(2)`$ as the flavor symmetry group Barbieri:1996uv which has two attractive features: (i) it gives rise to the degeneracies between 1-2 families needed to suppress the supersymmetric FCNC in the squark sector, and (ii) a multi-step breaking of $`U(2)`$ gives rise to the observed inter-family hierarchy naturally. Unlike models based on the most commonly used $`U(1)`$ symmetry, in which one has the freedom in choosing $`U(1)`$ charges for various matter fields, a $`U(2)`$ flavor symmetry appears to be a much more constrained framework for constructing realistic models. The basic idea is very simple. The three families of matter fields transform under a $`U(2)`$ flavor symmetry as
$$\psi _a\psi _3=21$$
(27)
where $`a=1,2`$ and the subscripts refer to family indices. In the symmetric limit, only the third family of matter fields have non-vanishing Yukawa couplings. This can be understood easily since the third family of matter fields have much higher masses compared to the other two families of matter fields. $`U(2)`$ breaks down in two steps:
$$U(2)\stackrel{ϵM}{}U(1)\stackrel{ϵ^{}M}{}nothing$$
(28)
with $`ϵ^{}ϵ1`$ and $`M`$ is the UV cut-off of the effective theory mentioned before. These small parameters $`ϵ`$ and $`ϵ^{}`$ are the ratios of the vacuum expectation values of the flavon fields to the cut-off scale. Note that since
$`\psi _3\psi _31_S,\psi _3\psi _a2`$
$`\psi _a\psi _b22=1_A3`$ (29)
the only relevant flavon fields are in the $`1_A,2`$ and $`3`$ dimensional representations of $`U(2)`$, namely,
$$A^{ab}1_A,\varphi ^a2,S^{ab}3$$
(30)
Because we are confining ourselves to symmetric mass textures, we use only $`\varphi ^a`$ and $`S^{ab}`$. Since all the $`16`$ observed matter fields of each family fall nicely into a $`16`$dimensional spinor representation of $`SO(10)`$, the most general superpotential that generates fermion masses for a $`SO(10)\times U(2)`$ model has the following very simple form
$$W=H(\psi _3\psi _3+\psi _3\frac{\varphi ^a}{M}\psi _a+\psi _a\frac{S^{ab}}{M}\psi _b)$$
(31)
In a specific $`U(2)`$ basis,
$$\frac{\varphi }{M}O\left(\begin{array}{c}ϵ^{}\\ ϵ\end{array}\right),\frac{S^{ab}}{M}O\left(\begin{array}{cc}ϵ^{}& ϵ^{}\\ ϵ^{}& ϵ\end{array}\right)$$
(32)
Here we have indicated the VEVs all the flavon fields could acquire for symmetry breaking in eq.(28). The mass matrix would take the following form
$$MO\left(\begin{array}{ccc}ϵ^{}& ϵ^{}& ϵ^{}\\ ϵ^{}& ϵ& ϵ\\ ϵ^{}& ϵ& 1\end{array}\right)$$
(33)
In $`SO(10)`$, at the renormalizable level, only three types of Higgs fields can couple to fermions,
$$1616=10_S120_A126_S$$
(34)
namely, $`10`$, $`120_A`$, and $`\overline{126}_S`$, where the subscripts $`S`$ and $`A`$ refer to the symmetry property under interchanging two family indices in the Yukawa couplings $`𝒴_{ab}`$. That is,
$$𝒴_{ab}^{10}=𝒴_{ba}^{10},𝒴_{ab}^{120}=𝒴_{ba}^{120},𝒴_{ab}^{\overline{126}}=𝒴_{ba}^{\overline{126}}$$
(35)
$`\varphi ^a`$ and $`S^{ab}`$ can couple to only $`10_S`$ and $`\overline{126}_S`$; $`120_A`$ has no role in giving rise to mass textures. Note that $`SO(10)`$ can break down to SM through many different breaking chains. Different breaking chains give rise to different mass relations among the up-quark, down-quark, charged lepton and neutrino sectors. Since we are interested in symmetric mass textures, a natural choice is the left-right symmetric route, that is,
$$\begin{array}{ccc}SO(10)\hfill & \hfill & SU(4)\times SU(2)_L\times SU(2)_R\hfill \\ & \hfill & SU(3)\times SU(2)_L\times SU(2)_R\times U(1)_{BL}\hfill \\ & \hfill & SU(3)\times SU(2)_L\times U(1)_Y\hfill \\ & \hfill & SU(3)\times U(1)_{EM}\hfill \end{array}$$
(36)
We have the up-quark sector related to the neutrino sector, and the down-quark sector to the charged lepton sector. A Clebsch-Gordon coefficient $`(3)`$ appears in the lepton sectors when the $`SU(4)\times SU(2)_L\times SU(2)_R`$ components $`(15,2,2)`$ in $`\overline{126}`$ are involved in the Yukawa couplings. This factor of $`(3)`$ is very crucial for obtaining the Georgi-Jarlskog relations as we have seen in the previous section. The general fermion Dirac mass matrices are thus given schematically by
$`M_u𝒴_{ab}^{10}10^++𝒴_{ab}^{\overline{126}}\overline{126}^+`$
$`M_d𝒴_{ab}^{10}10^{}+𝒴_{ab}^{\overline{126}}\overline{126}^{}`$
$`M_e𝒴_{ab}^{10}10^{}3𝒴_{ab}^{\overline{126}}\overline{126}^{}`$
$`M_{\nu _{LR}}𝒴_{ab}^{10}10^+3𝒴_{ab}^{\overline{126}}\overline{126}^+`$ (37)
and general Majorana mass matrices are given by
$$M_{\nu ,RR}𝒴_{ab}^{\overline{126}}\overline{126}^0$$
(38)
$$M_{\nu ,LL}𝒴_{ab}^{\overline{126}}\overline{126}^+$$
(39)
where various VEVs are those of the neutral components of $`SO(10)`$ representations as indicated below (with subscripts referring to the symmetry groups on the r.h.s. of eq.(36); and $`+/0/`$ referring to the sign of the hypercharge Y).
$$\begin{array}{c}10^+:(1,0)_{31}(1,2,1)_{321}(1,2,2,0)_{3221}(1,2,2)_{422}10\hfill \\ 10^{}:(1,0)_{31}(1,2,1)_{321}(1,2,2,0)_{3221}(1,2,2)_{422}10\hfill \end{array}$$
(40)
$$\begin{array}{c}\overline{126}^+:(1,0)_{31}(1,2,1)_{321}(1,2,2,0)_{3221}(15,2,2)_{422}\overline{126}\hfill \\ \overline{126}^{}:(1,0)_{31}(1,2,1)_{321}(1,2,2,0)_{3221}(15,2,2)_{422}\overline{126}\hfill \\ \overline{126}^0:(1,0)_{31}(1,1,0)_{321}(1,1,3,2)_{3221}(10,1,3)_{422}\overline{126}\hfill \\ \overline{126}^+:(1,0)_{31}(1,3,2)_{321}(1,3,1,2)_{3221}(\overline{10},3,1)_{422}\overline{126}\hfill \end{array}$$
(41)
A remark is in order here. Some models avoid the use of $`\overline{126}`$ dimensional Higgses by introducing nonrenormlaizable operators of the form $`f_af_b(16)_h(16)_h`$. Such models appear to be less constrained due to the inclusion of nonrenormalizable operators. Also, a discrete symmetry, the R-parity symmetry, must be imposed by hand to avoid dangerous Baryon number violating terms in the effective potential at low energies which otherwise could lead to fast proton decay rate. Here we use $`\overline{126}`$ dimensional representation of Higgses which has the advantage that R-parity symmetry is automatic Mohapatra:rparity1986a . The $`\overline{126}`$ representation has been used in model building before Aulakh:1999cd . It is to be noted that the contribution of the $`\overline{126}`$-dimensional representation to the $`\beta `$-function makes the model nonperturbative (with the onset of the Landau pole) above the unification scale $`M_{GUT}`$. One could view our model as an effective theory valid below this scale where coupling constants are perturbative.
Other breaking chains of $`SO(10)`$ have been considered resulting in various interesting mass textures and thus mass relations. For example, Albright:1998vf considers $`SO(10)`$ breaking through $`SU(5)`$ to SM and obtains the so-called ”lopsided” mass textures due to the fact that $`SU(5)`$ gives the relation $`M_d=M_e^T`$. A large lepton mixing $`(U_{MNS})`$ arises in this class of models from a large left-handed charged lepton mixing which relates to a large mixing in the right-handed down-quark sector.
A Model Based on $`\mathrm{𝐒𝐎}(\mathrm{𝟏𝟎})\times 𝐔(\mathrm{𝟐})`$: We now demonstrate how the above phenomenology emerges from a model based on $`SO(10)\times U(2)`$. Here we only present the Yukawa sector. A more complete account of the Higgs potential including symmetry breaking sector and the doublet-triplet splitting sector will be given elsewhere chen:fmass2000a .
In order to uniquely specify the Yukawa superpotential without any unwanted interaction terms, we need to introduce $`Z_2\times Z_2\times Z_2`$ discrete symmetry. The fields needed are indicated below
Matter fields:
$`\psi _a(16,2)^{++}`$ $`(a=1,2)`$
$`\psi _3(16,1)^{+++}`$ (42)
Higgs fields for the mass matrices:
$`(10,1):`$ $`T_1^{+++},T_2^+,T_3^+`$
$`T_4^{},T_5^+`$
$`(\overline{126},1):`$ $`\overline{C}^{},\overline{C}_1^{+++},\overline{C}_2^{++}`$ (43)
Flavon fields:
$`(1,2):`$ $`\varphi _{(1)}^{++},\varphi _{(2)}^{++},\mathrm{\Phi }^+`$
$`(1,3):`$ $`S_{(1)}^+,S_{(2)}^{},\mathrm{\Sigma }^{++}`$ (44)
Note that, the entries in the parenthesis indicate the $`SO(10)`$ and $`U(2)`$ representations respectively. The superscript $`+/`$ indicates the charges under $`Z_2\times Z_2\times Z_2`$ symmetry. Various Higgs fields acquire VEVs in the following directions
$`T_1:`$ $`10_1^+,10_1^{}`$
$`T_2,T_3,T_4:`$ $`10_{2,3,4}^+`$
$`T_5:`$ $`10_5^{}`$
$`\overline{C}:`$ $`\overline{126}^{}`$
$`\overline{C}_1,\overline{C}_2:`$ $`\overline{126}_{1,2}^0`$ (45)
and
$`10_1^+=10_3^+,10_1^{}=10_5^{}`$
$`\overline{126}_1^{}_{}{}^{}0=\overline{126}_2^{}_{}{}^{}0`$ (46)
(Note that, with a $`\overline{126}_H`$ acquiring VEV, there must be a conjugate $`126_H`$ acquiring VEV to cancel the D-term. Since $`126_H`$ does not couple to $`16_i`$, it has no role in the construction of the Yukawa sector.) The needed flavon VEVs are given by
$`\varphi _{(1)}=\left(\begin{array}{c}ϵ^{}\\ 0\end{array}\right),\varphi _{(2)}=\left(\begin{array}{c}0\\ ϵ\end{array}\right)`$ (51)
$`S_{(1)}=\left(\begin{array}{cc}0& ϵ^{}\\ ϵ^{}& 0\end{array}\right),S_{(2)}=\left(\begin{array}{cc}0& 0\\ 0& ϵ\end{array}\right)`$ (56)
$`\mathrm{\Phi }=\left(\begin{array}{c}\delta _1\\ \delta _3\end{array}\right),\mathrm{\Sigma }=\left(\begin{array}{cc}0& 0\\ 0& \delta _2\end{array}\right)`$ (61)
Our $`(Z_2)^3`$ charge assignments give rise to a unique superpotential:
$$W=W_{Dirac}+W_{\nu _{RR}}$$
(62)
$`W_{Dirac}=\psi _3\psi _3T_1+{\displaystyle \frac{1}{M}}\psi _3\psi _a\left(T_2\varphi _{(1)}+T_3\varphi _{(2)}\right)`$
$`+{\displaystyle \frac{1}{M}}\psi _a\psi _b\left(T_4+\overline{C}\right)S_{(2)}+{\displaystyle \frac{1}{M}}\psi _a\psi _bT_5S_{(1)}`$
$`W_{\nu _{RR}}=\psi _3\psi _3\overline{C}_1+{\displaystyle \frac{1}{M}}\psi _3\psi _a\mathrm{\Phi }\overline{C}_2+{\displaystyle \frac{1}{M}}\psi _a\psi _b\mathrm{\Sigma }\overline{C}_2`$ (63)
The mass matrices then can be read from the superpotential to be
$`M_{u,\nu _{LR}}`$ $`=`$ $`\left(\begin{array}{ccc}0& 0& 10_2^+ϵ^{}\\ 0& 10_4^+ϵ& 10_3^+ϵ\\ 10_2^+ϵ^{}& 10_3^+ϵ& 10_1^+\end{array}\right)`$ (67)
$`=`$ $`\left(\begin{array}{ccc}0& 0& r_2ϵ^{}\\ 0& r_4ϵ& ϵ\\ r_2ϵ^{}& ϵ& 1\end{array}\right)M_U`$ (71)
$`M_{d,e}`$ $`=`$ $`\left(\begin{array}{ccc}0& 10_5^{}ϵ^{}& 0\\ 10_5^{}ϵ^{}& (1,3)\overline{126}^{}ϵ& 0\\ 0& 0& 10_1^{}\end{array}\right)`$ (75)
$`=`$ $`\left(\begin{array}{ccc}0& ϵ^{}& 0\\ ϵ^{}& (1,3)pϵ& 0\\ 0& 0& 1\end{array}\right)M_D`$ (79)
where
$$M_U10_1^+,M_D10_1^{}$$
(80)
$`r_210_2^+/10_1^+,`$ $`r_410_4^+/10_1^+`$
$`p\overline{126}^{}/10_1^{}`$ (81)
The right-handed neutrino mass matrix is
$`M_{\nu _{RR}}`$ $`=`$ $`\left(\begin{array}{ccc}0& 0& \overline{126}_2^{}_{}{}^{}0\delta _1\\ 0& \overline{126}_2^{}_{}{}^{}0\delta _2& \overline{126}_2^{}_{}{}^{}0\delta _3\\ \overline{126}_2^{}_{}{}^{}0\delta _1& \overline{126}_2^{}_{}{}^{}0\delta _3& \overline{126}_1^{}_{}{}^{}0\end{array}\right)`$ (85)
$`=`$ $`\left(\begin{array}{ccc}0& 0& \delta _1\\ 0& \delta _2& \delta _3\\ \delta _1& \delta _3& 1\end{array}\right)M_R`$ (89)
with $`M_R\overline{126}_1^{}_{}{}^{}0`$. We have thus arrived at the mass matrices shown in eq.(7), (23) and (24).
RGE Analysis and Results: In order to obtain the input parameters at the GUT scale, first we need to know various Yukawa couplings (the diagonal elements) and mixing angles at the GUT scale. We use the expressions derived from 1-loop RGEs given by Arason:rge1992a ; Berezhiani:fmass1999a :
$`m_u=Y_u^0R_u\eta _uB_t^3v_u,`$ $`m_c=Y_c^0R_u\eta _cB_t^3v_u`$
$`m_t=Y_c^0R_uB_t^6v_u`$
$`m_d=Y_d^0R_d\eta _dv_d,`$ $`m_s=Y_s^0R_d\eta _sv_d`$
$`m_b=Y_b^0R_d\eta _bB_tv_d`$
$`m_e=Y_e^0R_ev_d,`$ $`m_\mu =Y_\mu ^0R_ev_d`$
$`m_\tau =Y_\tau ^0R_ev_d`$ (90)
$$V_{ij}=\{\begin{array}{ccc}V_{ij}^0,\hfill & & ij=ud,us,cd,cs,tb\hfill \\ V_{ij}^0B_t^1,\hfill & & ij=ub,cb,td,ts.\hfill \end{array}$$
(91)
where $`V_{ij}`$ are CKM matrix elements; quantities with superscript $`0`$ are evaluated at GUT scale, and all the $`m_f`$ and $`V_{ij}`$ are the experimental values PDG:exp1998a . We will assume $`\mathrm{tan}\beta =\frac{v_u}{v_d}=10`$ and $`v=\sqrt{v_u^2+v_d^2}=\frac{246}{\sqrt{2}}GeV.`$ The running factor $`\eta _f`$ includes QCD + QED contributions: For $`f=b,c`$, $`\eta _f`$ is for the range $`m_f`$ to $`m_t`$, and for $`f=u,d,s`$, $`\eta _f`$ is for the range $`1GeV`$ to $`m_t`$;
$`\eta _u=\eta _d=\eta _s=2.38_{0.19}^{+0.24}`$
$`\eta _c=2.05_{0.11}^{+0.13}`$
$`\eta _b=1.53_{0.04}^{+0.03}.`$
$`R_{u,d,e}`$ are contributions of the gauge-coupling constants running from weak scale $`M_z`$ to the SUSY breaking scale, taken to be $`m_t`$, with the SM spectrum, and from $`m_t`$ to the GUT scale with MSSM spectrum;
$$R_u=3.53_{0.07}^{+0.06},R_d=3.43_{0.06}^{+0.07},R_e=1.50.$$
$`B_t`$ is the running induced by large top-quark Yukawa coupling defined by
$$B_t=\mathrm{exp}\left[\frac{1}{16\pi ^2}_{\mathrm{ln}M_{SUSY}}^{\mathrm{ln}M_{GUT}}Y_t^2(\mu )d(\mathrm{ln}\mu )\right]$$
(92)
which varies from $`0.7`$ to $`0.9`$ corresponding to the perturbative limit $`Y_t^03`$ and the lower limit $`Y_t^00.5`$ imposed by the top-pole mass.
In order to have a good fit to these values, we first obtain the following approximate analytic expressions
$`a\sqrt{{\displaystyle \frac{Y_u^0Y_c^0}{Y_t^0(Y_c^0+c^2Y_t^0)}}},`$ $`bc^2+{\displaystyle \frac{Y_c^0}{Y_t^0}}`$
$`dY_t^0`$
$`e\sqrt{{\displaystyle \frac{Y_e^0}{Y_\mu ^02Y_e^0}}}{\displaystyle \frac{Y_\mu ^0Y_e^0}{Y_\tau ^0}},`$ $`f{\displaystyle \frac{Y_\mu ^0Y_e^0}{3Y_\tau ^0}}`$
$`h=Y_\tau ^0`$ (93)
With the GUT scale values of $`Y_e^0,Y_\mu ^0`$ and $`Y_\tau ^0`$, the three parameters $`e,f`$, and $`h`$ in the down-quark and charged lepton sectors are uniquely determined. With the GUT scale values of $`Y_u^0,Y_c^0`$ and $`Y_t^0`$, these relations reduce the number of parameters in the up-quark and Dirac neutrino sectors from four down to one, the parameter $`c`$. Using the GUT scale value of the Cabbibo angle, $`V_{us}`$, the value of $`c`$ is determined. At the GUT scale which is taken to be $`M_{GUT}=2.39\times 10^{16}GeV`$, with $`g_1=g_2=g_3=0.7530`$, our input parameters are chosen to be:
$`a=\alpha =0.00226,`$ $`b=\beta =0.00381`$
$`c=\gamma =0.0328,`$ $`d=\eta =0.572`$
$`e=0.00403,`$ $`f=0.0195`$
$`h=0.0678`$
$`\delta _1=0.00116,`$ $`\delta _2=3.32\times 10^5`$
$`\delta _3=0.0152`$
$`M_R=1.32\times 10^{14}GeV`$ (94)
These parameters are related to $`ϵ`$ and $`ϵ^{}`$ given in eq.(51) and Higgs VEVs and their ratios given in eq.(80) chen:fmass2000a . $`\delta _i`$’s could be obtained using $`t=1\times 10^3`$ in eq.(From CKM Matrix to MNS Matrix: A Model Based on Supersymmetric $`SO(10)\times U(2)_F`$ Symmetry); $`\mathrm{\Lambda }`$ is expressible in terms of $`\frac{\eta ^2}{M_R}`$ due to eq.(26). The Yukawa couplings in the down-quark and charged lepton sectors are then given by
$`Y_d^0=0.00005441,`$ $`Y_s^0=0.001374`$
$`Y_b^0=0.06779`$
$`Y_e^0=0.00001880,`$ $`Y_\mu ^0=0.003979`$
$`Y_\tau ^0=0.06779`$ (95)
and various ratios are given by
$$\frac{Y_d^0}{Y_e^0}=2.895,\frac{Y_s^0}{Y_\mu ^0}=\frac{1}{2.895},\frac{Y_b^0}{Y_\tau ^0}=1$$
(96)
which agree with Georgi-Jarlskog relations.
Having determined the GUT scale values of these elements, we then numerically solve the one-loop RGEs for the MSSM spectrum with three right-handed neutrinos Babu:rge1993a from GUT scale to the effective right-handed neutrino mass scale, $`M_R1.32\times 10^{14}GeV`$. At $`M_R`$, seesaw mechanism is implemented. We then run the MSSM RGEs Arason:rge1992a from $`M_R`$ down to the SUSY breaking scale $`m_t176GeV`$, and then the SM RGEs from $`m_t`$ to $`M_z=91.187GeV`$. The light neutrino RGEs Babu:rge1993a are also used from $`M_R`$ to $`M_z`$. Predictions obtained at $`M_z`$ are summarized in Table 1, taking into account the SUSY threshold corrections Hall:rge1994a
$$\mathrm{\Delta }_s=0.10,\mathrm{\Delta }_b=0.25$$
They are to be compared with the values at $`M_z`$ calculated from the experimental values by the authors of Fusaoka:exp1998a .
The quark mixing matrix $`V_{CKM}`$ at $`M_z`$ is predicted to be
$`\left|V_{CKM,predict}\right|=|V_{u_L}V_{d_L}^{}|`$
$`=\left(\begin{array}{ccc}0.9751& 0.2215& 0.003541\\ 0.2215& 0.9745& 0.03695\\ 0.004735& 0.03681& 0.9993\end{array}\right)`$ (100)
They are to be compared with the experimental results extrapolated to $`M_z`$ Fusaoka:exp1998a
$$\left|V_{CKM,exp}\right|=\left(\begin{array}{ccc}0.97450.9757& 0.2190.224& 0.0020.005\\ 0.2180.224& 0.97360.9750& 0.0360.046\\ 0.0040.014& 0.0340.046& 0.99890.9993\end{array}\right)$$
(101)
Our model predicts the three light Majorana neutrino masses to be
$`m_{\nu _1}=2.0052\times 10^4eV`$
$`m_{\nu _2}=2.0123\times 10^4eV`$
$`m_{\nu _3}=0.05574eV`$ (102)
and the resulting squared mass differences are
$`\mathrm{\Delta }m_{23}^2=3.11\times 10^3eV^2`$
$`\mathrm{\Delta }m_{12}^2=2.87\times 10^{10}eV^2`$ (103)
The lepton mixing matrix is given by
$`\left|U_{MNS,predict}\right|=\left|U_{e_L}U_{\nu _{LL}}^{}\right|`$
$`=\left(\begin{array}{ccc}0.6710& 0.7396& 0.0527\\ 0.5410& 0.4397& 0.7169\\ 0.5070& 0.5096& 0.6952\end{array}\right)`$ (107)
This translates into
$`\mathrm{sin}^22\theta _{atm}4|U_{\mu \nu _3}|^2(1|U_{\mu \nu _3}|^2)=0.9992`$
$`\mathrm{sin}^22\theta _{}4|U_{e\nu _2}|^2(1|U_{e\nu _2}|^2)=0.9912.`$ (108)
These values agree with the Super-Kamiokande atmospheric neutrino oscillation data SuperK:1998a ; LP99atm , and the solar VO solution LP99solar . And the $`(1,3)`$ element of $`U_{MNS}`$ is given by $`|U_{e\nu _3}|=0.0527`$ which is far below the bound by the CHOOZ experiment $`|U_{e\nu _3}|0.16`$ Apollonio:1999ae . The three eigenvalues of the right-handed neutrino Majorana mass matrix are given by
$`M_{RR_1}2.963\times 10^7GeV`$
$`M_{RR_2}2.643\times 10^{10}GeV`$
$`M_{RR3}1.319\times 10^{14}GeV`$ (109)
We can have the LOW solution (– a LAMSW solution with $`\mathrm{\Delta }m_{12}^210^610^7eV^2`$) with
$$\delta _1=0.001147,\delta _2=0.0002354,\delta _3=0.01675$$
(110)
$$M_R=1.615\times 10^{13}GeV$$
These change the predictions of $`m_{u,c,t}`$ by less than $`1\%`$ but have no observable effects on down-quark and charged lepton masses, and the CKM matrix remains essentially the same chen:fmass2000a . In the neutrino sector, we get
$`m_{\nu _1}=0.001626eV`$
$`m_{\nu _2}=0.001650eV`$
$`m_{\nu _3}=0.06303eV`$ (111)
and the squared mass differences are
$`\mathrm{\Delta }m_{23}^2=3.973\times 10^3eV^2`$
$`\mathrm{\Delta }m_{12}^2=1.298\times 10^7eV^2`$ (112)
The lepton mixing matrix is given by
$`\left|U_{MNS,predict}\right|=|U_{e_L}U_{\nu _{LL}}^{}|`$
$`=\left(\begin{array}{ccc}0.6665& 0.7418& 0.07428\\ 0.5511& 0.4231& 0.7192\\ 0.5021& 0.5202& 0.6909\end{array}\right)`$ (116)
The element $`|U_{e\nu _3}|`$ is predicted to be 0.07428, which is less than the experimental upper bound. The three right-handed neutrino eigenvalues are predicted to be
$`M_{RR_1}9.558\times 10^7GeV`$
$`M_{RR_2}8.453\times 10^8GeV`$
$`M_{RR_3}1.615\times 10^{13}GeV`$ (117)
It is also possible to have the LAMSW solution with
$$\delta _1=0.001082,\delta _2=0.0009870,\delta _3=0.02238$$
(118)
$$M_R=2.415\times 10^{12}GeV$$
The predictions in the quark and the charged lepton sectors remain the same. In the neutrino sector, we get
$`m_{\nu _1}=0.01089eV`$
$`m_{\nu _2}=0.01206eV`$
$`m_{\nu _3}=0.09999eV`$ (119)
and the squared mass differences are
$`\mathrm{\Delta }m_{23}^2=9.851\times 10^3eV^2`$
$`\mathrm{\Delta }m_{12}^2=2.752\times 10^5eV^2`$ (120)
The lepton mixing matrix is given by
$`\left|U_{MNS,predict}\right|=|U_{e_L}U_{\nu _{LL}}^{}|`$
$`=\left(\begin{array}{ccc}0.6439& 0.7486& 0.1580\\ 0.6045& 0.3712& 0.7049\\ 0.4690& 0.5494& 0.6915\end{array}\right)`$ (124)
The element $`|U_{e\nu _3}|`$ is predicted to be 0.1580 which is right at the experimental bound $`|U_{e\nu _3}|0.16`$ Apollonio:1999ae . The three right-handed neutrino eigenvalues are given by
$`M_{RR_1}5.732\times 10^6GeV`$
$`M_{RR_2}1.177\times 10^9GeV`$
$`M_{RR_3}2.417\times 10^{12}GeV`$ (125)
We note that a $`|U_{e\nu _3}|`$ value of less than 0.1580 would lead to $`\mathrm{\Delta }m_{23}^2>10^2eV^2`$ leading to the elimination of the LAMSW solution in our model. This is a characteristic of the LAMSW solution with $`\mathrm{\Delta }m_{12}^210^5eV^2`$.
Note added: The form invariance of eq.(26) – $`M_{\nu _{LL}}`$ having the same texture as that of $`M_{\nu _{LR}}`$ and $`M_{\nu _{RR}}`$ – also occurs in a model of neutrino mixing Fritzsch:fmass1999a which uses different symmetric mass textures.
Summary: We have constructed a realistic model based on SUSY $`SO(10)`$ combined with $`U(2)`$ flavor symmetry. The up-quark sector is related to the Dirac neutrino sector, and the down-quark sector is related to the charged lepton sector via $`SO(10)`$ symmetry. The inter-family hierarchy is achieved via $`U(2)`$ symmetry. In contrast to the commonly used effective operator approach, we use $`126`$-dim Higgses to construct the Yukawa sector. R-parity symmetry is thus automatically preserved at low energies. In our model, the Dirac and right-handed Majorana neutrino mass matrices which have very small mixing combine with the seesaw mechanism resulting in a large mixing in the lepton sector. The symmetric mass textures arising from the left-right symmetry breaking chain of $`SO(10)`$ which we have considered give rise to very good predictions; 15 masses (including the right-handed neutrino masses) and 6 mixing angles are predicted by 11 parameters. Our model favors the vacuum oscillation and LOW solutions to the solar neutrino problem.
###### Acknowledgements.
We thank K.S. Babu, C. Carone, N. Irges and S. Oh for useful communications. This work was supported, in part, by the US Department of Energy Grant No. DE FG03-05ER40894.
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# Relativity of spatial scale and of the Hubble flow: The logical foundations of relativity and cosmology
## I Introduction
Traditional reasoning holds that a simultaneous expansion of the observer’s own body and instruments would render the cosmological expansion unobservable (MTW, , p719) (Rindler, , p197), and that accordingly, the observable expansion cannot include small scales. It has been lately shown that while the expansion could occur on all scales Anderson1995 , the unobservability argument does not get challenged, as even over planetary distances, the effects would be extremely small Cooperstock1998 .
I examine here the converse question: shouldn’t, then, a *contraction* of the observer, not driven by the metric, be considered at least as a mathematically equivalent picture of the Hubble flow? I show that this would not only reproduce the Hubble redshift in a static universe, but would *exactly* predict both the cosmological constant $`\mathrm{\Lambda }`$IV) and the Pioneer anomaly (§V), also obviating, in the first case, the current speculations of a large-scale repulsion. The hypothesised contraction also turns out to be consistent with unresolved evidence of a past expansion of the earth Wesson1973 IV), which is irreconcilable with any of the past theories.
In view of the erroneous objections raised by various referees, I need to point out that the prior relativistic notion of scale, even in the Brans-Dicke theories, is by definition defined by the relativistic metric, and signifies only a general order of magnitude of the observed entities (Wald, , p98), i.e. a *virtual* scale, in which the observer’s own neighbourhood, including its physical unit referents of scale, remains unscaled. This makes the existing framework of general relativity *logically incomplete*, as there is no way to account for mechanisms that might cause our unit referents to vary differently from the underlying metric, and the absence of such mechanisms cannot be considered *a priori* to be a law of physics. I have separately described the complete logical foundations of quantum mechanics Prasad2000b , and shown that the constancy of Planck’s constant $`h`$ does not suffice to establish equality of spatial scale from either quantum interactions or thermal equilibrium. Formalising the dependence on referents not only leads to a more precise interpretation of relativistic curvature (§III), but also to a fundamental and appropriately simple proof of the inherent consistency with quantum mechanics (§III, Appendix C), and to a purely logical derivation of the relativity postulates.
A second clarification is also clearly in order, concerning the obvious difficulty of accounting for a *source*-dependent redshift purely by the *observer*’s physics, which is why cosmological expansion and “tired”-light theories have been the only explanations conceived of. These theories are, however, incomplete for the same reason, viz that variability of the spatial scale of the observer’s instruments was not recognised in the prior formulation of quantum mechanics, which had evolved by matching conjectural notions with empirical grounds, rather than of sound reasoning from the fundamental definitions of mechanics. The issue is therefore unaddressed by relativistic field theory, which accounts only for the relativistic metric, and not possible non-metrical variation of referents. As particularly established now in the classical thermodynamic derivation of Planck’s law Prasad2000a , the radiation quantum corresponds to an incremental change of an antinodal lobe in a stationary mode of the receiver, and the energy of a lobe *per se* is independent of its wavelength, meaning that quantisation does not require the equality of spatial scale between interacting entities. The packetisation of incoming electromagnetic energy thus depends *entirely* on the thermodynamics of the observer.
## II Relativity of Hubble’s law
It is remarkable that despite the unobservability argument cited above and Feynman’s “hot-plate” model suggesting the variability of material referents (Feynman, , II-42-1), no formal consideration of the implied dependence on local referents exists in the founding concepts of relativity (EinsteinMeaning, , ch.1), or in the subsequent literature. The dependence is definitional, because even to conceive of a physical quantity $`S`$, unless it be a dimensionless ratio, we need reference to some material referent $`R`$. For example, in citing the wavelength of a Lyman-$`\alpha `$ line, we implicitly refer to the standard metre, itself defined as so many wavelengths of an atomic transition. In attributing the redshifts of stellar Lyman spectra to Doppler and gravitational effects only, we implicitly assume that terrestrial and stellar matter are otherwise inherently identical. Although this looks entirely reasonable, our physics can be neither logically complete nor precise without formally accounting for this circularity and examining how it could break our assumption. This formalisation is readily obtained from the simple observation that the numerical value of a quantity $`S`$ is necessarily a ratio
$$n=f(S)/f(R),$$
(1)
where $`f`$ denotes the “notch-counts” obtained in the physical course of measurement, and $`f(R)`$ signifies the calibration.
The possibility of virtual Hubble flow, which does break our past assumption, now arises as follows. The traditional derivation of Hubble’s law (Wald, , p98),
$$v(S)=\frac{dS}{dt}\frac{S}{a_v}\frac{da_v}{dt}=HS,H\frac{\dot{a}_v}{a_v},$$
(2)
concerns a “virtual” form of scale not relating to material referents $`R`$ or changes thereof (hence my subscript $`v`$), i.e. $`R`$ has been implicitly assumed to be constant in the past. Recalling from Dirac’s Large Number Hypothesis (LNH) Dirac1937 that even the constancy of the absolute constants cannot be assumed on this time scale, there can be justification for assuming $`R`$ to be constant on this scale either. Incorporating eq. (1) into eq. (2), and allowing $`R`$ to vary, we get
$$\begin{array}{cc}\hfill n(v(S))& \frac{d}{dt}\frac{f(S)}{f(R)}=\frac{1}{f(R)}\frac{df(S)}{dt}\frac{f(S)}{f(R)^2}\frac{df(R)}{dt}\hfill \\ & =\frac{1}{f(R)}\frac{df(S)}{dt}\frac{f(S)}{f(R)}\left[\frac{1}{f(R)}\frac{df(R)}{dt}\right],\hfill \end{array}$$
(3)
more succinctly expressible as
$$\begin{array}{cc}\hfill v(S)=\frac{dS}{dt}S\frac{da}{dt}=\frac{S}{a_v}\frac{da_v}{dt}\frac{S}{a}\frac{da}{dt}& =H_tS,\hfill \\ \hfill \text{where }H_t& =HH_r,\text{ and }H_r\dot{a}/a\hfill \end{array}$$
(4)
for the observed (total) redshift $`H_t`$, where $`H`$ continues to represent actual expansion; the unsubscripted $`a`$, a scale factor expressly referencing our “real” referent $`R`$; and $`H_r`$ quantifies the impact of $`R`$’s expansion on the observation. This is a clear and correct formalisation of the relativistic argument, as by setting $`\dot{a}/a=\dot{a}_v/a_v`$, $`a_v`$ denoting the scale factor in the the Friedmann-Robertson-Walker (FRW) formalism, we do get $`H_t=0`$ as expected.
Eq. (4) also admits, however, the possibility that *any local mechanism capable of causing a continuous variation of $`R`$ ($`\dot{a}0`$) could introduce redshifts obeying Hubble’s law*. For example, the atomic wavelengths used to define the metre Petley would necessarily change due to incremental gravitational redshift if the earth’s mean radius $`r_e`$ were to be changing for some reason, which has been hypothesised before NarliKem1988 Wesson1973 . By the circularity described above, the change would be undetectible on ground, as well as by satellites, as their orbital heights are dependent on the same ground referents. The circularity is broken only by the deep space missions, which is why the deep space anomaly acquires a fundamental significance to be described in §V.
## III Gravitational shrinkage
As the constancy of $`R`$ is a deeply ingrained belief (Dirac, , p3), I need to first show that such a variation is produced by relativity as a result of gravitational redshift. The basic idea is that if the internal scale $`a`$ were indeed independent of the electromagnetic wavelengths, the structure of matter should become incongruous with radiation in the presence of gravity. This is particularly suggested by the Huygens picture of the gravitational bending of light; specifically, the redshift is traditionally interpreted as scale-preserving, but the notion depends on assuming that $`c`$ varies Einstein1911 . To appreciate this and the present result, recall that in special relativity, a light clock essentially relates the local scales of length and time, and that there is no *a priori* reason for assuming either to remain absolutely constant. We could take an equivalent view in which $`c`$ is preserved and the length scale varies instead, as follows.
Consider two observers bearing light clocks, the first, subscripted $`1`$, stationed within a gravitational well and the other, subscribted $`2`$, located outside, their light clocks being of lengths $`l_i`$ and periods $`t_i`$, $`i=1,2`$, respectively. In terms of our more precise formalism of eq. (1), special relativity postulates that $`c`$ is locally preserved, i.e.
$$\frac{f_1(l_1)}{f_1(t_1)}=\frac{f_2(l_2)}{f_2(t_2)}=c,$$
(5)
where the $`f_i`$ denote notch-counts by the $`i`$-th observer. In Einstein’s variable-$`c`$ perspective, clock $`2`$ runs slower than clock $`1`$, meaning $`f_i(t_2)>f_i(t_1)`$. This merely implies $`f_2(l_2)/f_2(t_2)<f_2(l_2)/f_2(t_1)`$ and is insufficient to establish
$$\frac{f_2(l_2)}{f_2(t_2)}c<c_{1,2}\frac{f_2(l_1)}{f_2(t_1)},$$
(6)
$`c_{i,j}`$ denoting the speed of light at $`i`$-th location measured by observer $`j`$, for which we also need to assume $`f_2(l_2)=f_2(l_1)`$. There is no physical basis for this further assumption: the only *physical* relation we have is already given by eq. (5), with which all physical laws are presumably consistent, so $`c_{i,j}`$ cannot bear independent significance to physics. It should be entirely equivalent, therefore, for us to also take
$$c_{i,j}=c\text{ so that }f_i(l_2)>f_i(l_1).$$
(7)
This is a remarkable result, as illustrated by application to Schwartzschild geometry. Say our unit rod is a solid $`N`$ lattice constants in length; by eq. (7), its atoms should become more densely packed when the rod is taken into the earth. Direct measurement of its diameter by lining up unit rods should thus yield a larger number than that obtainable from the surface area or circumference, which would be measured by rods left on the surface. A simpler explanation of gravitational curvature is thus achieved, which is not only logically more precise than prior theory, but particularly brings out the conservation of quantised properties like the number density $`N`$. The result represents *covariance* of the quantum scale, and is clearly a less restrictive interpretation than constancy and the only one required by the laws of physics. Further support for this inherent quantum consistency is given in Appendix C, where the very postulates of relativity are logically derived from these notions.
The reality of the *virtual* Hubble flow, $`H_r`$, is also easily established by this equivalence. If observer $`1`$ is subject to an increasing acceleration $`\dot{g}`$, eq. (7) says that its local rods and clocks must shrink at rate $`H_r\dot{a}/a<0`$; all objects will then appear to be receding at velocities proportional to their distance, *replete with Doppler shift resembling the cosmological expansion*. This conflicts with our usual notions, since, from the perspective of observer $`2`$, identical photons from different sources, bearing no information of the respective distances, appear to be magically received by observer $`1`$ at varying redshifts Elmegreen1999pvt . It is, however, as should be, because the absorbed photon eigenfunctions must correspond to the stationary states of the receiver in the receiver’s reference frame, and those of observer $`1`$ cannot be stationary with respect to $`2`$. As observer $`1`$’s stationarity involves a decreasing $`c`$, relative to $`2`$, its atoms cannot “know” of this decrease, let alone that it is a local artifact, and must therefore perceive the optical path lengths as inexplicably increasing. The applicable eigenfunctions have long been described by Parker in another context Parker1968 Parker1969 and do display the distance-dependence. $`\mathrm{}`$
## IV Variable earth implications
A changing gravitation $`\dot{g}`$ on earth, for instance due to changing $`G`$ or $`r_e`$, could thus contribute to the Hubble redshift. Certain geological and paleomagnetic evidence have indicated a past expansion of the earth at about $`0.4`$-$`0.6`$ mm/y Runcorn1965 , which is still about two orders of magnitude too small to be directly verified. Incidentally, the indicated expansion is difficult to reconcile with known physics, and even variable-$`G`$ theories can account for only a fraction of this rate Wesson1973 . It is also of the wrong sign to account for the Hubble redshift, but fortunately too small to be significant: from the incremental redshift, we find $`\dot{a}t=\delta \mathrm{\Phi }/c^2gr_e/c^2`$, so that the full Hubble flow corresponds to a contraction rate of $`dr_e/dtc^2\dot{a}/gc^2Hr_e/g128.43`$ km/s, ruling out significant variable-$`G`$ and $`\dot{r}_e`$ contributions. Moreover, the circularity of scale prohibits a *uniform* expansion of all matter on earth from being observable.
It does not, however, rule out an ongoing contraction of *surface matter*, which would include our solid referents, due to a non-relativistic cause, as that could be just right, via eq. (1), to cause the earth to *appear* to have expanded at the empirically indicated rate. The possibility is valid as the evidence equivalently indicates that the sialic masses contracted *relative* to $`r_e`$ Wesson1973 , but it also implies a proportional recession of the moon and the planets, as well as of the distant stars. The past expansion indeed happens to be of the same order as $`H`$ MacDougall1963 ($`0.4\text{mm/y}65\text{km/s-Mpc}`$), and even the lunar recession and existing planetary range data are consistent with the small-scale Hubble flow, (see Appendix A). Additionally, an upper bound of $`z=2`$ exists for gravitational redshift (Wald, , §6.3), which, by our analysis in §III, would impose an upper bound on the redshift of the incoming photons; this too is clearly represented in the geological data, as $`r_e`$ appears to have at most doubled since the birth of the solar system Runcorn1965 . As mentioned, the contraction also exactly predicts $`\mathrm{\Lambda }`$, as follows.
Hubble’s law $`\dot{r}=H_tr`$ implies an intrinsic acceleration as a particle initially at distance $`r`$ must pick up the speed increment $`\delta \dot{r}H_t\delta r`$ by the time it reaches $`r+\delta r`$. This intrinsic acceleration $`\ddot{r}=\dot{H}_tr+H_t\dot{r}\dot{H}_tr+H_t^2r`$ already represents the deceleration factor
$$q\frac{\ddot{a}}{a}\frac{1}{H_t^2}=\frac{1+\dot{H}_t/H_t^2}{a}=(1+\dot{H}_t/H_t^2).$$
(8)
Now, $`\dot{H}_t\dot{H}\dot{H}_r`$, and per FRW theory, $`\dot{H}`$ has many possible variations depending on the undetermined curvature constant $`k`$ and the radiation pressure (Wald, , p98) (MTW, , p772-774), allowing for the well known gamut of values for the matter ($`\mathrm{\Omega }_M`$) and energy ($`\mathrm{\Omega }_\mathrm{\Lambda }`$) densities in the universe Reiss1998 Garnavich1998a . However, eq. (4) and its underlying semantics (eq. 1) bear no dependency on the past values of $`H_r`$, meaning that $`H_r`$ would always appear to be constant at the instant of measurement, i.e. $`\dot{H}_r=0`$ regardless of the past values of $`H_r`$. Thus, if the observed Hubble flow $`H_t`$ were indeed due to $`H_r`$ alone, we would get $`q=1`$ identically, which is precisely the value consistently indicated by Type Ia supernovae Reiss1998 Leibundgut1998 Garnavich1998b . $`\mathrm{}`$
## V Deep space confirmation
As mentioned in §II, measurements from deep space would be independent of the inherent circularity of scale, and should therefore be able to test the theory. Because of the immensely small order involved ($`H10^{18}`$ s), only the six deep space missions that employed spin-stablisation could provide ranging data of the necessary precision Anderson1998 , and *all six* have displayed an anomaly of precisely this order Anderson1998 . Only a constant part, $`h_cH`$, is actually to be accounted for by $`H_r`$, as the variations between missions and along the trajectories are conceivably due to mundane mechanisms, given that the temporal variations of $`\dot{g}\dot{r}/r^2`$ cannot account for its apparently oscillatory character Turyshev1999 . The inferred acceleration $`\delta g=r^1c^2\delta z`$ is stated to be quantitatively equivalent to the time dilation Anderson1998
$$h_c=c^1\delta g2.8\times 10^{18}\text{ s}^186\text{ km/s-Mpc}$$
(9)
signifying an expansion of onboard clocks *relative* to those on earth. Unlike the relativistic $`H`$, which is negligible at planetary range Cooperstock1998 , $`H_r`$ should produce a cosmological time dilation (CTD) as a uniform dilation of *all* clocks in the universe *relative* to our own, implying that the anomaly is indeed the CTD expected from $`H_r`$. This has two fundamental implications, first, that onboard the spacecraft, the observable Hubble flow must be $`H_tHh_c0`$, i.e. *the Hubble flow must be invisible from deep space*; and second, that the relativistic $`H`$, from FRW theory, must be zero as $`H_t`$ is almost the same for the distant stars, i.e. *the Hubble flow must be entirely due to terrestrial contraction.*
As a check, we may again write Hubble’s law as $`\dot{r}=hr`$, predicting a continually increasing unmodelled contribution in the signal path $`rct`$ to the spacecraft, with rate of change once again resembling an acceleration Rosales1998
$$\delta g\ddot{r}=\frac{d}{dt}\dot{r}=\frac{d}{dt}\{hct\}=hc,$$
(10)
identical to NASA’s computation of the equivalent time dilation, eq. (9). Eq. (10) does not suffice to prove planetary Hubble flow, as we cannot deduce $`\dot{r}=hr`$ from it. But the consistency of $`h_c`$ across all six missions that were at all equipped to measure it Anderson1998 , our inability to explain the constant residual part Turyshev1999 Anderson1999a , its consistency with lunar and planetary range data (see Appendix A) and $`\mathrm{\Lambda }`$IV), and the *total* absence of evidence to the contrary, together seem to be compelling indication that this is the case.
Eq. (10) is not the complete picture because unlike ordinary radar, the ranging procedure also involves a significant onboard segment in the signal path, comprising frequency conversion and phase-locked loops Bender1989 Vincent1990 Anderson1993 Anderson2000 . The effective length of this segment cannot possibly be constant to within the magnitude of the anomaly, viz $`O(10^{18})`$ s<sup>-1</sup>, as plastic flow under centrifugal action alone could cause expansion of at least this order, and its modulation due to the varying orbit conditions could account for the reported variations; in particular, this seems to explain the perihelion increase ($`\delta g12\times 10^{10}`$ m/s<sup>2</sup> at $`1.3`$ AU, in the case of Ulysses). The mechanism should likewise produce contraction under the compressive stress of earth’s gravity, and the tidal factors seem to correctly scale as well, as will be described separately, in strong support of the present theory.
## Conclusion
As mentioned, Feynman’s “hot-plate” model anticipates the formalism introduced here. Its power, demonstrated by the logical derivation of the relativity postulates and proof of the inherent consistency with quantum mechanics (Appendix C), reflects the improved logical precision in the treatment of physical variables. It should also be clear that the cosmological constant $`\mathrm{\Lambda }`$ and the Pioneer anomaly constitute two null indications, viz $`\dot{H}_r=0`$ and $`H_t=0`$ in space, respectively, in analogy to the Michelson-Morley result, that favour the present theory.
Admittedly, we do not as yet know how or if other successes of the standard model, such as the cosmic microwave background (CMB) and Olbers’ paradox, would be accomodated in the present theory, but the concern seems to be outweighed by the geological and deep space data near at hand. Some of these questions might already be answered in the “scale-expansion cosmology” (SEC) theory Masreliez1998 , as the latter’s conjecture of homogeneous expansion of scale is exactly provided by $`H_r`$. The predicted invisibility of the Hubble flow in deep space also remains to be verified.
###### Acknowledgements.
Many thanks are owed to Bruce G Elmegreen, A Joseph Hoane and Gyan Bhanot for substantial criticisms, help and guidance leading up to this work.
## Appendix A Evidence of planetary flow
I now show that the inference $`H_tH_rH0`$ is consistent with planetary, lunar and terrestrial data. Anderson *et al.* contend that the anomaly is missing in planetary ranging data, but they have sought only an actual acceleration in the planetary orbits, using the Viking data to set the upper bound of $`10^{11}`$ m/s<sup>2</sup>Anderson1998 . The absence of orbital influences invalidates the acceleration theory, but not the Hubble’s law prediction given in §V. The observability of $`H`$ on planetary and shorter scales has already been ruled out Cooperstock1998 ; what we need to examine is that of $`H_r`$, which is non-relativistic and larger than $`H`$ by over $`30`$ orders. Using eq. (4), in place of the FRW equations, we get recession rates of only $`1.6`$ $`\mu `$m/s for Jupiter ($`5`$ AU) and $`12.5`$ $`\mu `$m/s for Neptune or Pluto ($`40`$ AU), far too small for Doppler measurement. Even the yearly recession of Venus or Mars with respect to earth would be only about $`2`$-$`6`$ m, less than a tenth of the precision $`100`$-$`150`$ m cited by Anderson. We would have to have measurements spanning at least a $`100`$ years to notice the cumulative effect at all. Any ranging technique would also be inherently characterised by a relative precision $`\delta rh_\mu r`$, i.e. *the ranging error also follows Hubble’s law*. The anomaly was detected, therefore, only because increased sensitivity by two orders of magnitude intended to test general relativity Bender1989 Vincent1990 .
A more subtle problem is that successive sets of measurement tend to be treated as improvements, not comparisons, so that cumulative displacements, if at all detected, would have been attributed to “systematic error”, as described by Slade *et al.* for the Goldstone radar observations of Mercury Slade1998 . A discrepancy of $`+11`$ km is in fact reported for Jupiter in 1992 from the predicted orbital radius by the PEP740 ephemeris Harmon1994 , likely calibrated from 1970s data; the discrepancy amounts to a recession of $`611`$ m/y ($`H770`$ km/s-Mpc), about $`10`$ times too large and enough to completely mask the recession.
Lunar and geophysical measurements are more precise and just as much in favour. We would have been contradicted by the lunar recession if the latter were less than or almost equal to our “Hubble’s law” prediction ($`2.57`$ cm/y), leaving no room to accomodate a purely frictional component. However, the measured value is $`3.84`$ cm/y Lunar1994 , almost $`50`$% higher. The slowing of the earth’s rotation does not seem to be an independent means for estimating the friction, as it is already a factor in all the indications of the past expansion of the earth Runcorn1965 , and in any case, the plastic flow mechanism (§V) appears to be part of the tidal friction. The indictated value $`0.4`$-$`0.6`$ mm/y ($`H61`$-$`92`$ km/s-Mpc MacDougall1963 ), is once again two orders smaller than the available precision using the GPS, which should, because of the circularity of scale, be inherently incapable of detecting the expansion (§II, §IV). $`\mathrm{}`$
## Appendix B Verification
It has been suggested Elmegreen1999pvt that the theory could be verified on earth by measuring the relative expansion, or “aging”, of a sufficiently strong laser pulse reflected back and forth in an interferometer, traversing a total optical path length $`L`$. The resolving power of $`\delta \varphi `$ fringes would be matched within
$$\delta t=\frac{\lambda \delta \varphi }{H_rL}=\frac{L}{c},\mathrm{or}L=\sqrt{\frac{c\lambda \delta \varphi }{H_r}}.$$
(11)
For example, a shift of $`10^6`$ fringes at $`500`$ nm requires $`L=7317`$ km, using the Pioneer anomaly as the indication of $`H_r`$. With mirror reflectivity 99.9%, a 20 dB budget for mirror loss would allow 4600 reflections, requiring the cavity to be $`1.6`$ km in length. The phase shift would be distinguishable from the measurement uncertainty in $`L`$ from its $`\lambda ^{1/2}`$ dispersive character and linear growth during the measurement, given by
$$d\varphi /dt=H_rL/\lambda .$$
(12)
Any aging found should be entirely due to terrestrial contraction, as FRW contribution should be zero on earth. The method still presents several difficulties: for instance, the spectral spread of the pulse would also be amplified by the same ratio, and a refractive medium cannot be used to shorten the path, as there would be no way to distinguish the effects of the medium from true aging. It might be possible to overcome some of these by using matter waves instead of light, as $`\lambda `$ could be made smaller by several tens of orders.
## Appendix C Logical deduction of relativity
I show below that the formalism of §II also suffices for logically deducing the postulates and both theories of relativity, implying thereby that the relativity of scale, eq. (1), is the hitherto missing *logical* basis of the postulates of relativity. Accordingly, I deduce the postules solely on considerations of relative scale, viz that observers separated by space or time cannot possibly share their physical referents, and that sharing is impossible even between coexisting, colocated observers if they happen to be moving with respect to one another, so that the only way to compare their observations is by exchanging the numerical values from their measurements. This leads to the Lorentz transform by considering two such observers $`O`$ and $`O^{}`$, as follows.
$`O`$ cannot assume that $`O^{}`$ will arrive at the same numerical value $`n^{}f^{}(S)/f^{}(R^{})`$ when measuring a space-time interval $`S`$, as it cannot *a priori* assume that their “notch-counts” will match, i.e. in general, $`f^{}(R^{})f(R)`$ and $`f^{}(S)f(S)`$. We thus have four unknowns, $`R`$, $`R^{}`$, $`f`$ and $`f^{}`$ for determining the relative scale, and one common variable, the relative velocity $`𝐯`$, whose magnitude must be match between their perspectives. Each observer projects the other’s measurements on its own space-time plot, in particular, obtaining a mapping $``$ of the other’s coordinate axes onto its own. As $`𝐯`$ is the only physical aspect distinguishing the observers, $``$ must depend only on $`𝐯`$, and must be linear, as a more complex transformation would require additional distinguishing parameters, and would imply different physical conditions between their respective neighbourhoods. The conditions of coexistence and colocation mean that the origins can be trivially made to coincide, i.e. $`_𝐯(0,0)=(0,0)`$, and linearity means that the two sets of axes must be mutually inclined. There are only two possibilities for this, by a real or an imaginary rotation, and the first is ruled out because it could confuse the semantic distinction between space and time. $``$ must be an imaginary rotation, therefore, and further, it cannot depend on $`𝐯`$ directly, as we have not yet established the transformation of scales. Accordingly, we must have $`(\beta ),\beta 𝐯/c`$, where $`c`$ has the same dimensions as $`𝐯`$. We have thus deduced the *form* of the Lorentz transformation, and it remains for us to establish the constancy of $`c`$, which presumably relates to $`f`$ and $`f^{}`$, since the transformation must relate to the cross-measurements of the referents, $`f^{}(R)`$ and $`f(R^{})`$.
To derive the special relativity postulates, we first observe that $``$ cannot be of fundamental significance unless the physical interactions used in the measurement, $`f`$ and $`f^{}`$, are themselves fundamental. It would negate our purpose to *a priori* assume electromagnetism, or the strong or weak nuclear forces, to be fundamental, but such assumptions are unnecessary in the formalism, as the relativity postulates follow logically from the notion of fundamentality:
1. The “notch-counting functions” must be *analytic*, i.e. continuous in a path-independent manner, with respect to displacement of the *observer* relative to the measured object $`S`$. As is well known from partial differential theory (Sokolnikoff, , §62,65,66), analyticity over any two dimensions $`x`$ and $`t`$ is expressed by the Cauchy-Riemann (CR) equations
$$\frac{f_x}{x}+\frac{1}{c}\frac{f_t}{t}=0\text{ and }\frac{f_t}{x}\frac{1}{c}\frac{f_x}{t}=0,$$
(13)
where $`f_x`$ and $`f_t`$ are measuring functions along $`x`$ and $`t`$, the constant $`c`$ is simply a scale factor relating the dimensions. In our context, as $`x`$ is identifiable with space and $`t`$ with time, this $`c`$ has the requisite dimensions of speed and is identifiable with the Lorentz scale factor. Furthermore, the CR conditions yield the wave equation
$$\frac{^2𝐟}{x^2}\frac{1}{c^2}\frac{^2𝐟}{t^2}=0,\text{ where }𝐟=f_x,f_t$$
(14)
showing that the interactions $`f`$ and $`f^{}`$ could be used to communicate, and that the communication would be limited by the speed $`c`$.
2. Given a multitude of such physical means that however differ in speed, we would eliminate all but the fastest interactions, as the information returned by a slow process would be rendered redundant if a faster process could be employed at the same time. As $`f`$ and $`f^{}`$ are physical interactions, this also applies to the transmission of physical effects. For example, sound would be eliminated because it is slower than electromagnetism, and in every case we can employ sound, electromagnetism is also involved in the form of intermolecular forces. Thus, $`c`$ must be the maximum of all physical speeds.
3. (B) also implies that all fundamental physical means of interaction must have the same speed, as a fundamental interaction cannot be redundant. We would infer, for instance, that electromagnetic and gravitational waves must travel at the same speed $`c`$, if we knew they were both fundamental.
4. (B) appears to favour an infinite value for $`c`$, but that would also mean, by (C), that all fundamental interactions would operate instantaneously. It would be impossible to construct any physical clock whatsoever, because *no finite sequence* of fundamental physical interactions, meaning no realisable process, could yield a delayed action. An infinite speed of interaction is also functionally equivalent to colocation, hence *for the dimensions of space and time to exist at all, the fundamental forces must operate at finite speed*. This proves $`c<\mathrm{}`$.
5. Eqs. (13-14) clearly hold for quantum wavefunctions, which, by definition, are analytic and represent stationary modes of the waves in eq. (14). Moreover, the wavefunctions are amplitudes of probabilities that are equivalent to information in Shannon theory, so that the stationary modes literally represent physical information that would become available on measurement. We do not need empirical authority, therefore, beyond that which already led to the non-relativistic quantum theory, in order to identify our $`c`$ with that in quantum field theory, and, as a special case, with the speed of light, also establishing inherent consistency between relativity and quantum mechanics (§III).
General relativity concerns the complementary case of possibly *comoving* observers unable to share their referents because of a spatial or temporal separation $`x^\mu `$, which means that the physical conditions in their neighbourhoods can no longer be assumed to be identical, so that linearity cannot hold. Analyticity is still a valid assumption, yielding the general curvilinear transformation which is the basis of general relativity,
$$𝒢(x^\mu ):dx_{}^{}{}_{}{}^{\nu }=g_\nu ^\mu dx^\kappa ,ds_{}^{}{}_{}{}^{2}=g_{\nu \kappa }(x^\mu )dx_{}^{}{}_{}{}^{\nu }dx_{}^{}{}_{}{}^{\kappa }.$$
(15)
The relativity of scale is thus adequate logical foundation for both theories of relativity. $`\mathrm{}`$
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# I Introduction
## I Introduction
This paper is devoted to the study of entanglement of quantum states, which is one of the most decisively non-classical features in quantum theory. The question of quantifying entanglement in the case of mixed quantum states represented by density operators on finite dimensional Hilbert spaces has recently been studied extensively in the context of quantum information theory, see, e.g., and references therein.
An *entanglement measure* is a real-valued function defined on the set of density operators on some tensor product Hilbert space subject to further physically motivated conditions, see, e.g., and below. A number of entanglement measures have been discussed in the literature, such as the von Neumann reduced entropy, the relative entropy of entanglement , the entanglement of distillation and the entanglement of formation . Several authors proposed physically motivated postulates to characterize entanglement measures, see below. These postulates (although they vary from author to author in the details) have in common that they are based on the concepts of the operational formulation of quantum mechanics . We shall discuss one version of these *operational characterizations* of entanglement measures in Section IV.
In this paper we introduce new entanglement measures based on the greatest cross norm on the tensor product of the sets of trace class operators on Hilbert space (see Sections V and VI). We shall show that the measures introduced in this work satisfy all the basic requirements for entanglement measures. These include convexity, invariance under local unitary transformations, and non-increase under procedures composed of local quantum operations and classical communication.
Throughout this paper the set of trace class operators on some Hilbert space $``$ is denoted by $`𝒯()`$ and the set of bounded operators on $``$ by $`()`$. A density operator is a positive trace class operator with trace one.
## II Preliminaries
In this section we collect some basic definitions and results which are used in the course of this paper.
In the present paper we restrict ourselves mainly to the situation of a composite quantum system consisting of two subsystems with Hilbert space $`_1_2`$ where $`_1`$ and $`_2`$ denote the Hilbert spaces of the subsystems (except in Section VI). The states of the system are identified with the density operators on $`_1_2`$.
###### Definition 1
Let $`_1`$ and $`_2`$ be two Hilbert spaces of arbitrary dimension. A density operator $`\varrho `$ on the tensor product $`_1_2`$ is called *separable* or *disentangled* if there exist a family $`\left\{\omega _i\right\}`$ of positive real numbers, a family $`\left\{\rho _i^{(1)}\right\}`$ of density operators on $`_1`$ and a family $`\left\{\rho _i^{(2)}\right\}`$ of density operators on $`_2`$ such that
$$\varrho =\underset{i}{}\omega _i\rho _i^{(1)}\rho _i^{(2)},$$
(1)
where the sum converges in trace class norm.
The set of states is a convex set and its extreme points, which are also called *pure states*, are the projection operators. Every pure state obviously correspond to a unit vector $`\psi `$ in $`_1_2`$. We denote the projection operator onto the subspace spanned by the unit vector $`\psi `$ by $`P_\psi `$.
The Schmidt decomposition is of central importance in the characterization and quantification of entanglement associated with pure states.
###### Lemma 2
Let $`_1`$ and $`_2`$ be Hilbert spaces of arbitrary dimension and let $`\psi _1_2`$. Then there exist a family of non-negative real numbers $`\{p_i\}_i`$ and orthonormal bases $`\{a_i\}_i`$ and $`\{b_i\}_i`$ of $`_1`$ and $`_2`$ respectively such that
$$\psi =\underset{i}{}\sqrt{p_i}a_ib_i.$$
The family of positive numbers $`\{p_i\}_i`$ is called the family of *Schmidt coefficients* of $`\psi `$. For pure states the family of Schmidt coefficients of a state completely characterizes the amount of entanglement of that state. A pure state $`\psi `$ is separable if and only if $`\psi =ab`$ for some $`a_1`$ and $`b_2`$.
The *von Neumann reduced entropy* for density operators $`\sigma `$ on a tensor product Hilbert space $`_1_2`$ is defined as
$$S_{\mathrm{vN}}(\sigma ):=\mathrm{Tr}__1(\mathrm{Tr}__2\sigma \mathrm{ln}(\mathrm{Tr}__2\sigma )),$$
(2)
where $`\mathrm{Tr}__1`$ and $`\mathrm{Tr}__2`$ denote the partial traces over $`_1`$ and $`_2`$ respectively. In the case of pure states $`\sigma =P_\psi `$, it can be shown that $`\mathrm{Tr}__1(\mathrm{Tr}__2P_\psi \mathrm{ln}(\mathrm{Tr}__2P_\psi ))=`$
$`\mathrm{Tr}__2(\mathrm{Tr}__1P_\psi \mathrm{ln}(\mathrm{Tr}__1P_\psi ))=_ip_i\mathrm{ln}p_i`$ where $`\{p_i\}_i`$ denotes the family of Schmidt coefficients of $`\psi `$. However, for a general mixed state $`\sigma `$ we have $`\mathrm{Tr}__1(\mathrm{Tr}__2\sigma \mathrm{ln}(\mathrm{Tr}__2\sigma ))\mathrm{Tr}__2(\mathrm{Tr}__1\sigma \mathrm{ln}(\mathrm{Tr}__1\sigma ))`$.
## III Effects and Operations
In this section we recall some of the fundamental concepts and definitions in the operational approach to quantum theory and in quantum measurement theory . The quantum mechanical *state* of a quantum system is described by a density operator $`\varrho `$ on the system’s Hilbert space $``$, i.e., by a positive trace class operator with trace one. Let $`𝒦`$ be another Hilbert space. An *operation* is a positive linear map $`T:𝒯()𝒯(𝒦)`$ such that $`T`$ is trace non-increasing for positive trace class operators, i.e., $`0\mathrm{Tr}(T(\sigma ))\mathrm{Tr}(\sigma )`$ for all positive $`\sigma 𝒯().`$ Following we adopt the point of view that allowed operations in a laboratory are (O1) adding an ancilla, (O2) tracing out part of the system, (O3) performing unitary operations, and (O4) performing possibly selective yes-no experiments. It can be shown (for a detailed proof see, e.g., ) that the class of operations $`T:𝒯()𝒯(𝒦)`$ composed out of operations of the form (O1)-(O4) coincides with the class of trace non-increasing *completely positive* operations, i.e., has the property that for all $`n0`$ the map $`T_n`$ on $`𝒯(^n)`$ defined by $`T_n:=T1_n`$, where $`1_n(^n)`$ denotes the unit matrix, is positive. For a further physical motivation of the requirement of complete positivity see, e.g., . In the sequel it is always understood that all operations are completely positive. If $``$ and $`𝒦`$ are both finite dimensional Hilbert spaces, then it follows from the Choi-Kraus representation theorem for operations that for every operation $`T:𝒯()𝒯(𝒦)`$ there exists a family of bounded operators $`\{A_k:𝒦\}_k`$ with $`_kA_kA_k^{}1_𝒦`$ such that $`T`$ can be expressed as
$$T(\sigma )=\underset{k}{}A_k^{}\sigma A_k$$
(3)
for all $`\sigma 𝒯()`$. If $`=𝒦`$, the Choi-Kraus representation is also valid for infinite dimensional Hilbert spaces (all sums converge in trace class norm). The family $`\{A_k\}_k`$ is not unique. However, the operator $`E:=_kA_kA_k^{}=T^{}(1)`$ is independent of the family $`\{A_k\}_k`$ chosen and is called the *effect* corresponding to the operation $`T`$ and its associated yes-no measurement ($`T^{}`$ denotes the adjoint of $`T`$, ). Generally, an operator $`E`$ is called an *effect operator* if $`E`$ is bounded and Hermitean and if $`0E1`$. *Effect operator valued measures* are then the most general *observables* in the theory . They are also called *positive operator valued (POV) measures*. A *Lüders-von Neumann operation* is an operation of the form $`T_L(\sigma )=_kP_k\sigma P_k`$ where $`\{P_k\}_k`$ is a set of mutually orthogonal projection operators on $``$. Lüders-von Neumann operations are repeatable. In the case of a general operation, it is possible to view the terms in its Choi-Kraus representation as representing different possible measurement outcomes. In the terminology of operational quantum theory the individual terms in the Choi-Kraus representation (3) form a set of operations corresponding to coexistent effects, see : two effect operators $`E_1`$ and $`E_2`$ are called *coexistent* if there exist effect operators $`F,G,H`$ with $`F+G+H1`$ such that $`E_1=F+G`$ and $`E_2=F+H`$ (in general $`F,G`$ and $`H`$ will not be unique however). Therefore in general two coexistent effect operators $`E_1`$ and $`E_2`$ do not correspond to mutually complementary measurement outcomes but instead may have some ‘overlap’ represented by the operator $`F`$ even if $`E_1+E_21`$. Coexistent effect operators need not commute.
## IV Entanglement measures
An entanglement measure is a functional $`E`$ defined on the set of density operators on the Hilbert space of a composite quantum system measuring the degree of entanglement of every given density operator. Every measure of entanglement $`E`$ should satisfy the following requirements
* An entanglement measure is a positive real-valued functional $`E`$ which for any given two systems is well-defined on the set $`𝒟(_1_2)`$ of density operators on the tensor product $`_1_2`$ of the Hilbert spaces of the two systems. Moreover, $`E`$ is *expansible*, i.e., whenever $`\rho 𝒟(_1_2)𝒟(𝖧_1𝖧_2)`$ with embeddings $`_1𝖧_1`$ and $`_2𝖧_2`$ of Hilbert spaces $`_1`$ and $`_2`$ into larger Hilbert spaces $`𝖧_1`$ and $`𝖧_2`$ respectively, then $`E|_{_1_2}(\rho )=E|_{𝖧_1𝖧_2}(\rho )`$.
* If $`\sigma `$ is separable, then $`E(\sigma )=0`$.
* Local unitary transformations leave $`E`$ invariant, i.e.,
$$E(\sigma )=E\left((U_1U_2)\sigma (U_1^{}U_2^{})\right)$$
for all unitary operators $`U_1`$ and $`U_2`$ acting on $`_1`$ or $`_2`$ respectively.
* Entanglement cannot increase under procedures consisting of local operations on the two quantum systems and classical communication. If $`T`$ is an operation which is trace-preserving on positive operators and can be realized by means of local operations and classical communication, i.e., is composed out of local operations of the form (O1) - (O4) and classical communication, then
$$E(T(\sigma ))E(\sigma )$$
(4-a)
for all $`\sigma 𝒟(_1_2)`$. It is clear that every procedure acting on an individual single quantum system $`_1_2`$ composed only of local operations and classical communication can formally be represented as a finite sequence of operations of the form $`T_1T_2`$, where $`T_1`$ and $`T_2`$ are local operations on $`_1`$ and $`_2`$ respectively. The requirement that entanglement cannot increase under local operations and classical communication is thus equivalent to
$$E((T_1T_2)(\sigma ))E(\sigma ),$$
(4-b)
for all $`\sigma 𝒟(_1_2)`$.
###### Remark 3
Equation (4-b) stipulates that local operations cannot increase entanglement. In the quantum information literature most authors replace Equations (4-a) and (4-b) by the stronger requirement
$$\underset{i}{}p_iE(\sigma _i)E(\sigma ).$$
(5)
Equation 5 stipulates that after the measurement the entanglement (as measured by $`E`$) averaged over the possible output states $`\sigma _i`$ is less than or equal to the original entanglement. Here $`p_i`$ denotes the probability that the final state $`\sigma _i`$ occurs. In the literature Equation 5 is normally taken as the formal expression for the paradigm that it is impossible to create or increase entanglement by performing procedures composed of local quantum operations and classical communication alone. A disadvantage of Equation (5) is that it makes sense only in measurement situations and that the ‘possible output states’ corresponding to a given operation $`T`$ are not uniquely defined. Mathematically this corresponds to the fact that the Choi-Kraus representation of an operation $`T`$ is in general not unique. The difference between Equations (4-a) and (5) is that Equation (5) stipulates that entanglement cannot increase *on average* under local operations and classical communications (for a detailed discussion of this point see ). In contrast Equation (4-a) says that entanglement cannot increase for any operation which acts on individual systems and is composed of local operations and classical communication. If one takes up the former (ensemble) point of view of Equation (5), then Equation (4-b) does no longer represent the most general condition because from the ensemble point of view the most general operations composed out of local operations and classical communications can contain correlations between terms of the Choi-Kraus representations of subsequent local operations. A precise definition can be found, e.g., in . Some authors consider also other classes of local operations, most prominently the class of *separable* operations considered in .
* Mixing of states does not increase entanglement, i.e., $`E`$ is convex
$$E(\lambda \sigma +(1\lambda )\tau )\lambda E(\sigma )+(1\lambda )E(\tau )$$
for all $`0\lambda 1`$ and all $`\sigma ,\tau 𝒟(_1_2)`$.
Apart from the requirements (E0) - (E4) on entanglement measures many authors add further requirements to the definition of entanglement measures but we are not going to discuss them in this paper. For a details, see . In the sequel we exclude the trivial functional $`E0`$ which also satisfies (E0) - (E4).
###### Remark 4
Postulate (E2) is an immediate consequence of (E3).
###### Remark 5
It has been argued in that the entanglement of distillation $`E_D`$ introduced in does vanish for certain non-separable states (so called bound entangled states). Therefore it has been pointed out in that replacing (E1) by the stronger requirement that for every entanglement measure $`E(\sigma )=0`$ if and only if $`\sigma `$ is separable might exclude interesting entanglement measures. For more information the reader is referred to the references.
###### Example 6
Post selection of a subensemble means selecting a (non-normalized) output state of a quantum operation and normalizing its trace to 1. This procedure can lead to an increase in entanglement. This can be seen by considering a very simple example. Consider a composite quantum system composed of two 3-level quantum systems and the state
$$\rho _ϵ=(1ϵ)|0000|+\frac{ϵ}{2}(|12|21)(12|21|).$$
For $`ϵ`$ small it is intuitively obvious that this state does not contain ‘much’ entanglement and every entanglement measure should reflect this. Indeed, consider for example the relative entropy of entanglement introduced in defined by
$$E_S(\sigma ):=\underset{\rho }{inf}\left(\mathrm{Tr}(\sigma \mathrm{ln}\sigma \sigma \mathrm{ln}\rho )\right)$$
(6)
where the infimum runs over all separable states $`\rho `$ for which $`\mathrm{Tr}(\sigma \mathrm{ln}\rho )`$ is well-defined and finite. Elementary estimates using the results of show that
$$E_S(\rho _ϵ)ϵ\mathrm{ln}2.$$
If we subject the system to an operation testing whether or not the system is in the state $`|00`$ and select after the measurement the subensemble corresponding to the negative outcome (system is not in the state $`|00`$), then clearly the final state after the operation and post selection is given by $`\frac{1}{2}(|12|21)(12|21|)`$. Notice that this operation can be achieved by local operations and classical communication. We find
$$E_S\left(\frac{1}{2}(|12|21)(12|21|)\right)=\mathrm{ln}2>E_S(\rho _ϵ).$$
Similarly it can be shown that the entanglement measure $`_\gamma `$ to be introduced below may increase under post selection of subensembles. Therefore we see that we must not replace the operation in Equation (4-a) by some normalized non-linear operation $`\rho \frac{T(\rho )}{\mathrm{Tr}(T(\rho ))}`$ corresponding to post selection of a subensemble.
## V A new class of entanglement measures
Consider the situation that the two Hilbert spaces $`_1`$ and $`_2`$ are both finite dimensional and consider the spaces $`𝒯(_1)`$ and $`𝒯(_2)`$ of trace class operators on $`_1`$ and $`_2`$ respectively. Both spaces are Banach spaces when equipped with the trace class norm $`_1^{(1)}`$ or $`_1^{(2)}`$ respectively, see, e.g., Schatten . In the sequel we shall drop the superscript and write $`_1`$ for both norms, slightly abusing the notation; it will always be clear from the context which norm is meant. The algebraic tensor product $`𝒯(_1)_{\mathrm{alg}}𝒯(_2)`$ of $`𝒯(_1)`$ and $`𝒯(_2)`$ is defined as the set of all finite linear combinations of elementary tensors $`uv`$, i.e., the set of all finite sums $`_{i=1}^nu_iv_i`$ where $`u_i𝒯(_1)`$ and $`v_i𝒯(_2)`$ for all $`i`$.
It is well known that we can define a cross norm on $`𝒯(_1)_{\mathrm{alg}}𝒯(_2)`$ by
$$t_\gamma :=inf\{\underset{i=1}{\overset{n}{}}u_i_1v_i_1|t=\underset{i=1}{\overset{n}{}}u_iv_i\},$$
(7)
where $`t𝒯(_1)_{\mathrm{alg}}𝒯(_2)`$ and where the infimum runs over all finite decompositions of $`t`$ into elementary tensors. It is also well known that $`_\gamma `$ majorizes any subcross seminorm on $`𝒯(_1)_{\mathrm{alg}}𝒯(_2)`$. We denote the completion of $`𝒯(_1)_{\mathrm{alg}}𝒯(_2)`$ with respect to $`_\gamma `$ by $`𝒯(_1)_\gamma 𝒯(_2)`$. $`𝒯(_1)_\gamma 𝒯(_2)`$ is a Banach algebra .
As both $`_1`$ and $`_2`$ are finite dimensional, $`𝒯(_1)=(_1)`$ and $`𝒯(_2)=(_2)`$ and $`(_1)_{\mathrm{alg}}(_2)=(_1_2)`$, see, e.g., , Example 11.1.6. In finite dimensions all Banach space norms on $`(_1_2)`$, in particular the operator norm $``$, the trace class norm $`_1`$, and the norm $`_\gamma `$, are equivalent, i.e., generate the same topology on $`(_1_2)`$.
For later reference we compute the value of $`P_\psi _\gamma `$ for one dimensional projection operators $`P_\psi =|\psi \psi |`$ on $`_1_2`$ in terms of the coefficients in the Schmidt representation of $`|\psi `$. In this section we make extensive use of the Dirac bra-ket notation.
###### Proposition 7
Let $`|\psi _1_2`$ be a unit vector and $`|\psi =_i\sqrt{p_i}|\varphi _i|\chi _i`$ its Schmidt representation, where $`\{|\varphi _i\}_i`$ and $`\{|\chi _i\}_i`$ are orthonormal bases of $`_1`$ and $`_2`$ respectively and where $`p_i0`$ and $`_ip_i=1`$. Let $`P_\psi `$ denote the one dimensional projection operator onto the subspace spanned by $`|\psi `$. Then
$$P_\psi _\gamma =\underset{ij}{}\sqrt{p_ip_j}=\left(\underset{i}{}\sqrt{p_i}\right)^2.$$
*Proof*: Without loss of generality we assume that $`_1=_2`$ which can always be achieved by possibly suitably enlarging one of the two Hilbert spaces. Further, we identify $`_1=_2`$ with $`^n`$, where $`n=dim_1`$, i.e., we fix an orthonormal basis in $`_1`$ which we identify with the canonical real basis in $`^n`$. With respect to this canonical real basis in $`^n`$ we can define complex conjugates of elements of $`_1`$ and the complex conjugate as well as the transpose of a linear operator acting on $`_1`$. From the Schmidt decomposition it follows that
$$P_\psi =|\psi \psi |=\underset{ij}{}\sqrt{p_ip_j}|\varphi _i\varphi _j||\chi _i\chi _j|.$$
(8)
From the definition of $`_\gamma `$ it is thus obvious that $`P_\psi _\gamma _{ij}\sqrt{p_ip_j}`$. Now consider the Hilbert space $``$ of Hilbert-Schmidt operators on $`_1_2`$ equipped with the Hilbert-Schmidt inner product $`f|g=\mathrm{Tr}(f^{}g)`$. Equation (8) induces an operator $`𝔄_\psi `$ on $``$ as follows. Every element $`\zeta `$ in $``$ can be written $`\zeta =_kx_ky_k`$ where $`x_k`$ and $`y_k`$ are trace class operators on $`_1`$ and $`_2`$ respectively. Then $`𝔄_\psi `$ is defined on $`\zeta `$ as $`𝔄_\psi (\zeta ):=_{ijk}\sqrt{p_ip_j}\chi _i^{}|x_k|\chi _j^{}|\varphi _i\varphi _j|y_k`$ where $`|\chi _i^{}`$ denotes the complex conjugate of the vector $`|\chi _i`$ with respect to the canonical real basis in $`^n`$. Proposition 11.1.8 in implies that $`𝔄_\psi (\zeta )`$ is independent of the representation of $`\zeta `$. Consider a representation $`P_\psi =_{i=1}^ru_iv_i`$ of $`P_\psi `$ as sum over simple tensors. Denote the transpose of $`v_i`$ by $`v_i^T`$. Then the operator defined by
$$𝒜_\psi (\zeta ):=\underset{i,k=1}{\overset{r}{}}\mathrm{Tr}(v_i^Tx_k)u_iy_k$$
(9)
is equal to $`𝔄_\psi `$ (by virtue of Proposition 11.1.8 in ). We denote the trace on $`𝒯()`$ by $`\tau ()`$. The operator $`𝔄_\psi `$ is of trace class and the right hand side of Equation 8 is the so-called polar representation of $`𝔄_\psi `$ which implies $`\tau (𝔄_\psi )=_{ij}\sqrt{p_ip_j}`$, see . $`𝔄_\psi `$ admits also many other representations $`𝔄_\psi =_if_ig_i`$ with families of operators $`\{f_i\}`$ and $`\{g_i\}`$ acting on $`_1`$ and $`_2`$ respectively. It is well known that
$$\tau (𝔄_\psi )=inf\{\underset{i}{}f_i_2g_i_2|𝔄_\psi =\underset{i}{}f_ig_i\}P_\psi _\gamma ,$$
where $`_2`$ denotes the Hilbert-Schmidt norm and where the latter inequality follows from $`z_2z_1`$ and from the fact that each decomposition of $`𝔄_\psi `$ corresponds in an obvious one-to-one fashion to a decomposition of $`P_\psi `$. This proves the proposition. $`\mathrm{}`$
###### Corollary 8
Let $`\rho `$ be a density operator on $`_1_2`$, where $`_1`$ and $`_2`$ are finite dimensional Hilbert spaces. If $`\rho =_{ij}a_{ij}|\varphi _i\varphi _j||\chi _i\chi _j|`$, then $`\rho _\gamma =_{ij}|a_{ij}|.`$
An immediate corollary of Proposition 7 is that a pure state $`|\psi _1_2`$ is separable if and only if $`P_\psi _\gamma =1`$. In it has been proven that more generally all separable density matrices can be characterized by $`_\gamma `$.
###### Theorem 9
Let $`_1`$ and $`_2`$ be finite dimensional Hilbert spaces and $`\varrho `$ be a density operator on $`_1_2`$. Then $`\varrho `$ is separable if and only if $`\varrho _\gamma =1`$.
In it has been tentatively suggested that $`_\gamma `$ can be considered as a quantitative measure of entanglement. In the present work we substantiate this claim by proving
###### Proposition 10
The function
$$E_\gamma (\sigma ):=\sigma _\gamma \mathrm{log}\sigma _\gamma $$
satisfies the criteria (E0) - (E4) for entanglement measures.
*Proof*: (E1) is an immediate consequence of Theorem 9 and (E0) and (E2) are clear. (E3): Let $`T`$ be an operation composed of local operations, and classical communication. As we have argued above every such $`T`$ can be realized as a sequence of operations of the form $`T_1T_2`$ where $`T_1`$ and $`T_2`$ are local operations on system 1 and 2 respectively. We show that $`(T_1T_2)(\sigma )_\gamma \sigma _\gamma `$. By linearity of $`T_1T_2`$ every decomposition of $`\sigma `$ into finite sums of simple tensors $`\sigma =_{i=1}^rx_iy_i`$, where $`x_i`$ and $`y_i`$ are trace class operators on $`_1`$ and $`_2`$ respectively, induces a decomposition of $`(T_1T_2)(\sigma )`$ into a sum of simple tensors $`(T_1T_2)(\sigma )=_{i=1}^rT_1(x_i)T_2(y_i)`$. Thus
$`(T_1T_2)(\sigma )_\gamma `$ $`=`$ $`inf\{{\displaystyle \underset{i=1}{\overset{r}{}}}X_i_1Y_i_1|(T_1T_2)(\sigma )={\displaystyle \underset{i=1}{\overset{r}{}}}X_iY_i\}`$
$``$ $`inf\{{\displaystyle \underset{i=1}{\overset{r}{}}}T_1(x_i)_1T_2(y_i)_1|\sigma ={\displaystyle \underset{i=1}{\overset{r}{}}}x_iy_i\}`$
$``$ $`T_1T_2inf\{{\displaystyle \underset{i=1}{\overset{r}{}}}x_i_1y_i_1|\sigma ={\displaystyle \underset{i=1}{\overset{r}{}}}x_iy_i\}`$
$``$ $`inf\{{\displaystyle \underset{i=1}{\overset{r}{}}}x_i_1y_i_1|\sigma ={\displaystyle \underset{i=1}{\overset{r}{}}}x_iy_i\}`$
$`=`$ $`\sigma _\gamma ,`$
where we have used that both $`T_1`$ and $`T_2`$ are bounded maps on the space of trace class operators on $`_1`$ and $`_2`$ respectively and that
$$T_i=sup\{\mathrm{Tr}(T_i(\rho ))|\rho 𝒯(_i),\rho 0\mathrm{and}\mathrm{Tr}(\rho )=1\}1,$$
see, e.g., Lemma 2.2.1 in . (E3) follows from the fact that $`[1,\mathrm{}[ss\mathrm{log}s`$ is monotone. Finally, (E4) follows from the facts that $`_\gamma `$ is subadditive and that $`[1,\mathrm{}[ss\mathrm{log}s`$ is monotone and convex. $`\mathrm{}`$
###### Remark 11
It follows from the proof of Proposition 10 that if $`f`$ is a convex, monotonously increasing function on $`[1,\mathrm{}[`$ with $`f(1)=0`$, then
$$E_f(\sigma ):=f(\sigma _\gamma )$$
is an entanglement measure satisfying the requirements (E0) - (E4). A possible choice is $`f_1(x)=x1`$ leading to the entanglement measure $`E_{f_1}(\sigma )=\sigma _\gamma 1`$. This shows that indeed (as claimed in ) the function $`\sigma _\gamma 1`$ is an entanglement measure on the space of density operators. Other possible choices for $`f`$ are $`f_2(x)=x\mathrm{ln}xx+1`$, $`f_3(x)=e^{a(x1)},a>0`$ and so forth.
###### Corollary 12
The entanglement measures constructed in Remark 11 (including the measure $`E_\gamma `$ from Proposition 10) satisfy that $`E_f(\sigma )=0`$ if and only if $`\sigma `$ is separable.
*Proof*: This is an immediate consequence of Theorem 9. $`\mathrm{}`$
###### Proposition 13
Let $`T_1`$ and $`T_2`$ be two trace-preserving Lüders-von Neumann operations on finite dimensional Hilbert spaces $`_1`$ and $`_2`$ respectively and let $`T_L=T_1T_2`$ denote the corresponding Lüders-von Neumann operation acting locally on $`_1_2`$. Let $`T_1(\sigma _1)=_iP_i\sigma _1P_i`$ and $`T_2(\sigma _2)=_jQ_j\sigma _2Q_j`$ be Choi-Kraus representations of $`T_1`$ and $`T_2`$ respectively in terms of families $`\{P_i\}`$ and $`\{Q_j\}`$ of, respectively, mutually orthogonal projection operators. Then the entanglement measure $`E_{f_1}=_\gamma 1`$ as in Remark 11 satisfies
$$\underset{ij}{}p_{ij}\left(\sigma _{ij}_\gamma 1\right)\sigma _\gamma 1$$
where $`p_{ij}:=\mathrm{Tr}((P_iQ_j)\sigma (P_iQ_j))`$ and $`\sigma _{ij}=\frac{(P_iQ_j)\sigma (P_iQ_j)}{p_{ij}}`$ and where $`\sigma `$ is a density operator on $`_1_2`$.
*Proof*: Let $`P`$ and $`P^{}`$ be orthogonal projection operators. Then $`PxP+P^{}yP^{}_1=PxP_1+P^{}yP^{}_1`$ for all operators $`x,y`$. This follows from considering the spectral resolutions of $`PxP`$ and $`P^{}yP^{}`$. Hence $`_iP_ix_kP_i_1=_iP_ix_kP_i_1x_k_1`$. A similar argument shows that $`_jQ_j\stackrel{~}{z}Q_j_1\stackrel{~}{z}_1`$ for all $`\stackrel{~}{z}(_2)`$. Hence
$`{\displaystyle \underset{ij}{}}p_{ij}\left(\sigma _{ij}_\gamma 1\right)`$ $``$ $`inf\{{\displaystyle \underset{ijk}{}}P_ix_kP_i_1Q_jy_kQ_j_1|\sigma ={\displaystyle \underset{k}{}}x_ky_k\}1`$
$``$ $`inf\{{\displaystyle \underset{k}{}}x_k_1y_k_1|\sigma ={\displaystyle \underset{k}{}}x_ky_k\}1`$
$`=`$ $`\sigma _\gamma 1.\mathrm{}`$
It is known that some physically interesting entanglement measures coincide with the von Neumann reduced entropy on pure states, for instance the relative entropy of entanglement . However, it follows immediately from Proposition 7 that $`E_\gamma `$ does not coincide with the von Neumann reduced entropy on pure states: it follows from that the entropy of entanglement for a pure state of the form $`|\varphi =\alpha |00+\beta |11`$ is equal to $`|\alpha |^2\mathrm{ln}|\alpha |^2|\beta |^2\mathrm{ln}|\beta |^2`$, whereas it follows from Proposition 7 that $`E_\gamma \left(|\varphi \varphi |\right)=2\left(|\alpha |+|\beta |\right)^2\mathrm{ln}\left(|\alpha |+|\beta |\right).`$ Therefore we have explicitly constructed an entanglement measure satisfying a physically reasonable set of requirements which is not equal to the von Neumann reduced entropy on pure states. We have proven
###### Proposition 14
$`E_\gamma `$ and $`S_{\mathrm{vN}}`$ do not coincide on pure states.
In necessary and sufficient conditions for an entanglement measure to coincide with $`S_{\mathrm{vN}}`$ on pure states were derived. It is easy to see that, e.g., $`E_\gamma `$ does not satisfy the additivity condition (P4) considered in .
## VI Higher Tensor Product Hilbert Spaces
So far we restricted ourselves to tensor product Hilbert spaces of two finite dimensional Hilbert spaces. It is straightforward, however, to generalize our results to the situation of tensor products of more than two, but at most finitely many, finite dimensional Hilbert spaces. To this end consider the tensor product $`=`$ $`_1\mathrm{}_n`$ of $`n`$ finite dimensional Hilbert spaces $`_1,\mathrm{},_n`$. The obvious generalization of the definition of $`_\gamma `$ is
$$t_\gamma ^{(n)}:=inf\{\underset{i=1}{\overset{r}{}}u_i^{(1)}_1\mathrm{}u_i^{(n)}_1|t=\underset{i=1}{\overset{r}{}}u_i^{(1)}\mathrm{}u_i^{(n)}\},$$
(10)
where $`t`$ is a trace class operator on $``$.
It is straightforward to generalize the main result of to $`n`$-fold tensor product Hilbert spaces $`=`$ $`_1\mathrm{}_n`$
###### Definition 15
Let $`_1,\mathrm{},_n`$ be Hilbert spaces of arbitrary dimension. A density operator $`\varrho `$ on the tensor product $`_1\mathrm{}_n`$ is called *disentangled* or *separable* (with respect to $`_1,\mathrm{},_n`$) if there exist a family $`\left\{\omega _i\right\}`$ of positive real numbers, and families $`\left\{\rho _i^{(k)}\right\}`$ of density operators on $`_k`$ respectively, where $`1kn`$, such that
$$\varrho =\underset{i}{}\omega _i\rho _i^{(1)}\mathrm{}\rho _i^{(n)},$$
(11)
where the sum converges in trace class norm.
###### Theorem 16
Let $`_1,\mathrm{},_n`$ be finite dimensional Hilbert spaces and $`\varrho `$ be a density operator on $`=`$ $`_1\mathrm{}_n`$. Then $`\varrho `$ is separable if and only if $`\varrho _\gamma ^{(n)}=1`$.
We now consider the situation that $``$ is the $`m`$-fold tensor product of $`_1_2`$ with two finite dimensional Hilbert spaces $`_1`$ and $`_2`$. The functional $`E_\gamma `$ from Proposition 10 admits an obvious extension
$$E_\gamma (\sigma ):=\sigma _\gamma ^{(n)}\mathrm{ln}\sigma _\gamma ^{(n)}$$
(12)
for all trace class operators $`\sigma `$ on $``$.
###### Proposition 17
The functional defined by Equation (12) satisfies the criteria (E0)-(E4) for entanglement measures.
## VII Conclusion
To conclude, in this paper we have introduced a new class of entanglement measures on the space of density operators on tensor product Hilbert spaces. Our entanglement measures are based on the greatest cross norm $`_\gamma `$ on the set of trace class operators on the tensor product Hilbert space. We showed that our entanglement measures satisfy a number of physically desirable requirements, in particular that they do not increase under local quantum operations.
Acknowledgements
I thank Matthew J. Donald, Michał Horodecki and Michael Wolf for their comments on a previous version of this paper and Armin Uhlmann for pointing out Corollary 8 to me.
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# Emergence of Quantum Chaos in Quantum Computer Core and How to Manage It
## I Introduction
During the last decade, a remarkable progress has been achieved in the fundamental understanding of the main elements necessary for the creation of a quantum computer. Indeed, as stressed by Feynman , classical computers have tremendous problems to simulate very common quantum systems, since the computation time grows exponentially with the number of quantum particles. Therefore for such problems it is natural to envision a computer composed from quantum elements (qubits) which operate according to the laws of quantum mechanics. In any case, such devices will be in a sense unavoidable since the technological progress leads to chips of smaller and smaller size which will eventually reach the quantum scale. At present a quantum computer is viewed as a system of $`n`$ qubits (two-level quantum systems), with the possibility of switching on and off a coupling between them (see the detailed reviews in ). The operation of such computers is based on reversible unitary transformations in the Hilbert space whose dimension $`N_H=2^n`$ is exponentially large in $`n`$. It has been shown that all unitary operations can be realized with two-qubit transformations . This makes necessary the existence of a coupling between qubits. Any quantum algorithm will be a sequence of such fundamental transformations, which form the basis of a new quantum logic.
An important next step was the discovery of quantum algorithms which can make certain computations much faster than on a classical computer. The most impressive is the problem of factorization of large numbers in prime factors, for which Shor constructed a quantum algorithm which is exponentially faster than the classical ones. It was also shown by Grover that the searching of an item in a long list is parametrically much faster on a quantum computer. The recent development of error-correcting codes showed that a certain amount of noise due to external coupling could be tolerable in the operation of a quantum computer.
All these exciting developments motivated a great body of experimental proposals to effectively realize such a quantum computer. They include ion traps , nuclear magnetic resonance systems , nuclear spins with interaction controlled electronically or by laser pulses , quantum dots , Cooper pair boxes , optical lattices and electrons floating on liquid helium . As a result, a two-qubit gate has been experimentally realized with cold ions , and the Grover algorithm has been performed for three qubits made from nuclear spins in a molecule . However, to have a quantum computer competitive with a classical one will require a much larger number of qubits. For example, the minimal number of qubits for which Shor’s algorithm will become useful is of the order of $`n=1000`$ . As a result, a great experimental effort is still needed on the way to quantum computer realization.
A serious obstacle to the physical realization of such computers is the quantum decoherence due to the couplings with the external world which gives a finite lifetime to the excited state of a given qubit. This question has been discussed by several groups for different experimental qubit realizations . The effects of decoherence and laser pulse shape broadening were numerically simulated in the context of Shor’s algorithm , and shown to be quite important for the operability of the computer. However, in a number of physical proposals, for example nuclear spins in two-dimensional semiconductor structures, the relaxation time due to this decoherence process can be many orders of magnitude larger than the time required for the gates operation , so that there are hopes to manage this obstacle.
Here we will focus on a different obstacle to the physical realization of quantum computers that was not stressed up to now. This problem arises even if the decoherence time is infinite and the system is isolated/decoupled from the external world. Indeed, even in the absence of decoherence there are always imperfections in physical systems. Due to that the spacing between the two states of each qubit will fluctuate in some finite detuning interval $`\delta `$. Also, some residual static interaction $`J`$ between qubits will be unavoidably present (we remind that an inter-qubit coupling is required to operate the gates). Extensive studies of many-body interacting systems such as nuclei, complex atoms, quantum dots and quantum spin glasses have shown that generically in such systems the interaction leads to quantum chaos characterized by ergodicity of the eigenstates and level spacing statistics as in Random Matrix Theory (RMT) . In a sense the interaction leads to dynamical thermalization without coupling to an external thermal bath. If the quantum computer were in such a regime, its operability would be effectively destroyed since the noninteracting multi-qubit states representing the quantum register states will be eliminated by quantum ergodicity.
In this respect, it is important to stress that unavoidably the residual interaction $`J`$ will be much larger than the energy spacing $`\mathrm{\Delta }_n`$ between adjacent eigenstates of the quantum computer. Indeed the residual interaction $`J`$ is relatively small so that all $`N_H`$ computer eigenenergies are distributed in an energy band of size $`\mathrm{\Delta }En\mathrm{\Delta }_0`$, where $`\mathrm{\Delta }_0`$ is the average energy distance between the two levels of one qubit and $`n`$ is the total number of qubits in the computer. As a consequence, the spacing between multi-qubit states is $`\mathrm{\Delta }_n\mathrm{\Delta }E/N_Hn\mathrm{\Delta }_02^n\mathrm{\Delta }_0`$. Let us consider a realistic estimate for $`\mathrm{\Delta }_n`$ and $`J`$ for the case with $`n=1000`$ as required for Shor’s algorithm to be useful. For $`\mathrm{\Delta }_01`$ K, which corresponds to the typical one-qubit spacing in the experimental proposals , the multi-qubit spacing becomes $`\mathrm{\Delta }_n10^3\times 2^{10^3}\mathrm{\Delta }_010^{298}`$ K. This value will definitely be much smaller than any physical residual interaction. In the case of the proposal , for example, with a distance between donors of $`r=200`$ Å and an effective Bohr radius of $`a_B=30`$ Å ( Eq.(2) of ), the coupling between qubits (spin-spin interaction) is $`J\mathrm{\Delta }_01`$ K. By changing the electrostatic gate potential, the effective electron mass can be modified up to a factor of two. Since $`J(r/a_B)^{5/2}\mathrm{exp}(2r/a_B)/a_B`$, and $`a_B`$ is inversely proportional to the effective mass, this gives a minimal residual spin-spin interaction of $`J10^5`$ K $`\mathrm{\Delta }_n`$. In this situation, one would naturally/naively expect that level mixing, quantum ergodicity of eigenstates and chaos are unavoidable since the interaction is much bigger than the energy spacing between adjacent levels ($`J\mathrm{\Delta }_n`$).
In spite of this natural expectation, it was shown recently in that in the quantum computer the quantum chaos sets in only for couplings $`J`$ exponentially stronger than $`\mathrm{\Delta }_n`$. In fact, it was shown that the critical coupling $`J_c`$ for the transition to quantum chaos decreases only linearly with the number of qubits $`n`$ (for short-range inter-qubit coupling): $`J_c\mathrm{\Delta }_0/n`$. This result opens a broad parameter region where a quantum computer can be operated below the quantum chaos border, when noninteracting multi-qubit states are very close to the exact quantum computer eigenstates. For example, at $`n=1000`$ and $`\mathrm{\Delta }_01`$ K, the critical residual interaction is $`J_c1`$ mK, compatible with the proposal discussed above .
In the present paper, we study in more details the transition to chaos and how it affects the time evolution of the system. The effects of residual interaction in the presence or absence of fine fluctuations of individual qubit energy spacing are analyzed in great detail. The paper is composed as follows. In the next section we describe the standard generic quantum computer (SGQC) model, introduced in . In section III, we present the result of numerical and analytical studies of eigenenergies and eigenstate properties of this model. Section IV is devoted to the analysis of the time evolution of this system, and the typical time scales for the development of quantum chaos are presented as a function of the system parameters. We end by some concluding remarks in the last section.
## II Standard generic quantum computer model
In the standard generic quantum computer (SGQC) model was introduced to describe a system of $`n`$ qubits containing imperfections which generate a residual inter-qubit coupling and fluctuations in the energy spacings between the two states of one qubit. The Hamiltonian of this model reads:
$$H=\underset{i}{}\mathrm{\Gamma }_i\sigma _i^z+\underset{i<j}{}J_{ij}\sigma _i^x\sigma _j^x,$$
(1)
where the $`\sigma _i`$ are the Pauli matrices for the qubit $`i`$ and the second sum runs over nearest-neighbor qubit pairs on a two-dimensional lattice with periodic boundary conditions applied. The energy spacing between the two states of a qubit is represented by $`\mathrm{\Gamma }_i`$ randomly and uniformly distributed in the interval $`[\mathrm{\Delta }_0\delta /2,\mathrm{\Delta }_0+\delta /2]`$. The detuning parameter $`\delta `$ gives the width of the distribution near the average value $`\mathrm{\Delta }_0`$ and may vary from $`0`$ to $`\mathrm{\Delta }_0`$. Fluctuations in the values of $`\mathrm{\Gamma }_i`$ appear generally as a result of imperfections. For example, in the frame of the experimental proposals , the detuning $`\delta `$ will appear for nuclear spin levels as a result of local magnetic fields and density fluctuations. For electrons floating on liquid helium , it will appear due to fluctuations of the electric field near the surface. The couplings $`J_{ij}`$ represent the residual static interaction between qubits which is always present for reasons explained in the introduction. They can originate from spin-exciton exchange , Coulomb interaction , dipole-dipole interaction , etc… To catch the general features of the different proposals, we chose $`J_{ij}`$ randomly and uniformly distributed in the interval $`[J,J]`$. We note that a similar Hamiltonian, but without coupling/detuning fluctuations, was discussed for a quantum computer based on optical lattices . This SGQC model describes the quantum computer hardware, while the gate operation in time should include additional time-dependent terms in the Hamiltonian (1) and will be studied separately. At $`J=0`$ the noninteracting eigenstates of the SGQC model can be presented as $`|\psi _i>=|\alpha _1,\mathrm{},\alpha _n>`$ where $`\alpha _k=0,1`$ marks the polarization of each individual qubit. These are the ideal eigenstates of a quantum computer, and we will call them quantum register states. For $`J0`$, these states are no longer eigenstates of the Hamiltonian, and the new eigenstates are now linear combinations of different quantum register states. We will use the term multi-qubit states to denote the eigenstates of the SGQC model with interaction but also for the case $`J=0`$.
While in the main studies were concentrated on the case where $`\delta `$ is relatively large and comparable to $`\mathrm{\Delta }_0`$, here we will focus on the case $`\delta \mathrm{\Delta }_0`$, which corresponds to the situation where fluctuations induced by imperfections are relatively weak. In this case, the unperturbed energy spectrum of (1) (corresponding to $`J=0`$) is composed of $`n+1`$ well separated bands, with interband spacing $`2\mathrm{\Delta }_0`$. An example of the density of multi-qubit states $`\rho _n=1/\mathrm{\Delta }_n`$ in this situation is presented in Fig.1. Since the $`\mathrm{\Gamma }_i`$ randomly fluctuate in an interval of size $`\delta `$, each band at $`J=0`$, except the extreme ones, have a Gaussian shape with width $`\sqrt{n}\delta `$. The number of states inside a band is approximately $`N_H/n`$, so that the energy spacing between adjacent multi-qubit states inside one band is exponentially small ($`\delta _nn^{3/2}2^n\delta `$), in line with the general estimate in Section I.
In the presence of a residual interaction $`J\delta `$, the spectrum will still have the above band structure with exponentially large density of states. For $`J\delta \mathrm{\Delta }_0`$, the interband coupling is very weak and can be neglected. In this situation, the SGQC Hamiltonian (1) is to a good approximation described by the renormalized Hamiltonian $`H_P=\mathrm{\Sigma }_{k=1}^{n+1}\widehat{P_k}H\widehat{P_k}`$ where $`\widehat{P_k}`$ is the projector on the $`k^{th}`$ band, so that qubits are coupled only inside one band. We will thereafter concentrate our studies on the band nearest to $`E=0`$. For an even $`n`$ this band is centered exactly at $`E=0`$, while for odd $`n`$ there are two bands centered at $`E=\pm \mathrm{\Delta }_0`$, and we will use the one at $`E=\mathrm{\Delta }_0`$. Such a band corresponds to the highest density of states, and in a sense represents the quantum computer core. It is clear that quantum chaos and ergodicity will first appear in this band, which will therefore set the limit for operability of the quantum computer. Inside this band, the system is described by a renormalized Hamiltonian $`H_P`$ which depends only on the number of qubits $`n`$ and the dimensionless coupling $`J/\delta `$.
## III Quantum Computer Eigenenergies and Eigenstates
The first investigations in showed that the quantum chaos border in the SGQC model (1) corresponds to a critical interaction $`J_c`$ given by:
$$J_c\frac{C\delta }{n},$$
(2)
where $`C`$ is a numerical constant. This border is exponentially larger than the energy spacing between adjacent multi-qubit states $`\mathrm{\Delta }_n`$. The physical origin of this difference is due to the fact that the interaction is of a two-body nature. As a result, one noninteracting multi-qubit state $`|\psi _i>`$ has nonzero coupling matrix elements only with $`2n`$ other multi-qubit states. In the basis of quantum register states $`|\psi _i>`$, the Hamiltonian is represented by a very sparse nondiagonal matrix with only $`2n+1`$ nonzero matrix elements by line of length $`N_H=2^n`$. For $`\delta \mathrm{\Delta }_0`$ all these transitions take place in an energy interval $`B`$ of width of order $`6\mathrm{\Delta }_0`$. Therefore the energy spacing between directly coupled states is $`\mathrm{\Delta }_cB/2n3\mathrm{\Delta }_0/n`$. According to the studies of quantum chaos in many-body systems , the transition to chaos takes place when the matrix elements become larger than the energy spacing between directly coupled states. This gives $`J>\mathrm{\Delta }_c`$ which leads to the relation (2). For the case $`\delta \mathrm{\Delta }_0`$ on which we focus here, still in the renormalized Hamiltonian $`H_P`$ the number of nonzero matrix elements in one line is of the order of $`n`$, and $`B\delta `$, so that $`\mathrm{\Delta }_c\delta /n`$, that leads to the result (2) .
The transition to quantum chaos and ergodicity can be clearly seen in the change of the spectral statistics of the system. One of the most convenient is the level spacing statistics $`P(s)`$, which gives the probability to find two adjacent levels whose spacing is in $`[s,s+ds]`$. Here $`s`$ is the energy spacing measured in units of average level spacing. It is well known that while the average density of states is not sensitive to the presence or absence of chaos, the fluctuations of the energy spacings between adjacent levels around the mean value, determined by $`P(s)`$, are sensitive to it. In the presence of chaos, eigenstates are ergodic, overlap of wavefunctions gives a finite coupling matrix element between nearby states and the spectral statistics $`P(s)`$ follows the Wigner-Dyson (WD) distribution $`P_W(s)=(\pi s/2)\mathrm{exp}(\pi s^2/4)`$ typical for random matrices. This distribution $`P_W(s)`$ shows level repulsion at small $`s`$, due to the fact that overlap matrix elements between adjacent levels tend to move them away from each other. On the contrary, in the integrable case at $`JJ_c`$, the overlap coupling matrix element between nonergodic states is very small. As a result, energy levels are uncorrelated and $`P(s)`$ follows the Poisson distribution $`P_P(s)=\mathrm{exp}(s)`$ known to be valid for integrable one-particle systems .
In the SGQC model, we expect a transition from $`P_P(s)`$ at small $`J`$ to $`P_W(s)`$ above the quantum chaos border (2). An example of such a transition is shown in Fig.2. To decrease the statistical fluctuations we averaged over several independent realizations of the $`\mathrm{\Gamma }_i`$ and $`J_{ij}`$ in (1), which is the standard procedure used in Random Matrix Theory . We used up to $`N_D=5\times 10^4`$ realizations so that the total statistics $`1.5\times 10^5N_S>1.2\times 10^4`$. It is interesting to note that in the limit $`J/\delta \mathrm{}`$ ($`\delta J\mathrm{\Delta }_0`$) the system remains in the regime of quantum chaos with WD statistics , as is illustrated in Fig.3. This means that in the absence of individual qubit energy fluctuations, the residual coupling alone leads to chaotic eigenstates.
To characterize the variation of $`P(s)`$ from one limiting distribution to another it is convenient to use the parameter $`\eta =_0^{s_0}(P(s)P_W(s))𝑑s/_0^{s_0}(P_P(s)P_W(s))𝑑s`$ , where $`s_0=0.4729\mathrm{}`$ is the intersection point of $`P_P(s)`$ and $`P_W(s)`$. In this way $`P_P(s)`$ corresponds to $`\eta =1`$, and $`P_W(s)`$ to $`\eta =0`$. The studies of different systems has already shown that this parameter characterizes well the transition from one statistics to the other . Indeed, according to the data of Fig.4, $`\eta `$ changes from $`1`$ at small $`J`$ to $`\eta 0`$ at large $`J`$. To characterize this transition, we chose the critical value $`J_c`$ by the condition $`\eta (J_c)=0.3`$. The dependence of $`\eta `$ on the rescaled coupling strength $`J/J_c`$ shows that the transition becomes sharper and sharper when $`n`$ increases (Fig.4).
The dependence of the critical coupling strength $`J_c`$ on the number of qubits $`n`$ is shown on Fig.5. It clearly shows that this critical strength decreases linearly with $`n`$ and follows the theoretical border (2) with $`C3`$. For comparison on the same figure we also show the dependence of the multi-qubit spacing $`\mathrm{\Delta }_n`$ (computed numerically) on $`n`$. It definitely demonstrates that $`J_c\mathrm{\Delta }_n`$.
The transition in the level spacing statistics reflects a qualitative change in the structure of the eigenstates. While for $`JJ_c`$ the eigenstates are expected to be very close to the quantum register states $`|\psi _i>`$, for $`J>J_c`$ each eigenstate $`|\varphi _m>`$ becomes a superposition of an exponential number of states $`|\psi _i>`$. It is convenient to characterize the complexity of an eigenstate $`|\varphi _m>`$ by the quantum eigenstate entropy $`S_q=_iW_{im}\mathrm{log}_2W_{im}`$, where $`W_{im}`$ is the quantum probability to find the quantum register state $`|\psi _i>`$ in the eigenstate $`|\varphi _m>`$ of the Hamiltonian ($`W_{im}=|<\psi _i|\varphi _m>|^2`$). In this way $`S_q=0`$ if $`|\varphi _m>`$ is one quantum register state ($`J=0`$), $`S_q=1`$ if $`|\varphi _m>`$ is equally composed of two $`|\psi _i>`$, and the maximal value is $`S_q=n`$ if all $`2^n`$ states contribute equally to $`|\varphi _m>`$. We average $`S_q`$ over the states in the center of the energy band and $`N_D`$ realizations of $`\mathrm{\Gamma }_i`$ and $`J_{ij}`$.
The variation of this average $`S_q`$ as a function of $`J`$ for different values of $`n`$ is shown on Figs.6,7. It shows that indeed the entropy $`S_q`$ grows with $`J`$ until it saturates to a large value corresponding to an exponential number of mixed states. These data show that the critical coupling $`J_{cs}`$ at which $`S_q=1`$ (two states mixed) is proportional to $`J_c`$. Indeed, Fig.7 shows a small dispersion near $`S_q=1`$ when $`n`$ changes from 6 to 16, while $`\mathrm{\Delta }_n`$ varies by three orders of magnitude. This is confirmed by the data on Fig.5, which give $`J_{cs}0.13J_c0.4\delta /n`$. This result is in agreement with the results obtained by direct diagonalization of the SGQC model (1) at $`\delta \mathrm{\Delta }_0`$ (lower insert in Fig.2 of ).
The quantum eigenstate entropy $`S_q`$ characterizes the global properties of the eigenstates, while a more detailed information about them can be obtained from the local density of states $`\rho _W`$ introduced by Wigner :
$$\rho _W(EE_i)=\underset{m}{}W_{im}\delta (EE_m)$$
(3)
The function $`\rho _W`$ characterizes the average probability distribution of $`W_{im}`$ (see a numerical example in Fig.3 of ). For moderate coupling strength, $`\rho _W`$ is well described by the well-known Breit-Wigner distribution $`\rho _W=\rho _{BW}`$:
$$\rho _{BW}(EE_i)=\frac{\mathrm{\Gamma }}{2\pi ((EE_i)^2+\mathrm{\Gamma }^2/4)}$$
(4)
where $`\mathrm{\Gamma }`$ is the width of the distribution. This expression is valid when $`\mathrm{\Gamma }`$ is smaller than the bandwidth ($`\mathrm{\Gamma }<\sqrt{n}\delta `$) and many levels are contained inside this width. In this regime, the Breit-Wigner width $`\mathrm{\Gamma }`$ is given by the Fermi golden rule: $`\mathrm{\Gamma }=2\pi U_s^2\rho _c`$, where $`U_s`$ is the root mean square of the transition matrix element and $`\rho _c`$ is the density of directly coupled states. The validity of this formula was well checked in many-body systems with quantum chaos . In our case $`U_sJ`$ and $`\rho _cn/\delta `$, so that:
$$\mathrm{\Gamma }\frac{J^2n}{\delta }.$$
(5)
This dependence is confirmed by the data on Fig.8. However, for large $`J`$, when $`\mathrm{\Gamma }>\sqrt{n}\delta `$, the shape of $`\rho _W`$ becomes non-Lorentzian and is well fitted by a Gaussian distribution. The width of this modified distribution grows like $`\mathrm{\Gamma }J`$. This scaling naturally appears in the limit $`\delta =0,J\mathrm{\Delta }_0`$, since the noninteracting part of the Hamiltonian is simply a constant commuting with the perturbation. The change from one dependence to the other takes place for $`J>\delta /n^{1/4}`$. Above this limit $`\mathrm{\Gamma }`$ is still weakly dependent on the number of qubits $`n`$. We expect that for $`J\delta `$ the energy width of one band is $`\mathrm{\Gamma }J\sqrt{n}`$ (effective frequency of n Rabi frequencies with random signs), and have checked numerically this law for $`\delta =0`$ (data not shown).
According to the results obtained from many-body systems , the number of quantum register states mixed inside the width $`\mathrm{\Gamma }`$ is of the order of $`\mathrm{\Gamma }\rho _n`$, and is exponentially large. This however assumes that $`J>J_c`$ and the system is already in the quantum chaos regime. In this case the quantum eigenstate entropy $`S_q`$ is large ($`S_q\mathrm{log}_2(\mathrm{\Gamma }\rho _n)n`$) and the operability of the computer is quickly destroyed, since many quantum register states become mixed. The pictorial view of the quantum computer melting is shown on Fig.9. This image is qualitatively similar to the one in (Fig.5 there), which was obtained for the SGQC model at $`\delta =\mathrm{\Delta }_0`$. In Fig.9 the melting goes in a smoother way since all the states belong to the same central band (quantum computer core).
The effect of quantum chaos melting in the quantum register representation is shown on Fig.10 for $`J>J_c`$. The ideal register structure is manifestly washed out. On the contrary, below the chaos border ($`J<J_c`$), only few quantum register states are mixed. For comparison, Fig.11 shows the same part of the register in the regime $`JJ_{cs}`$ (no mixing of states) and Fig.12 in the regime $`JJ_{cs}`$ (few states are mixed).
## IV Time evolution in the SGQC model
In the previous section we determined the properties of eigenstates of the quantum computer in the presence of residual inter-qubit coupling. In the presence of this coupling the quantum register states $`|\psi _i>`$ are not any more stationary states, and therefore it is natural to analyze how they evolve in time. Indeed, if at time $`t=0`$ an initial state is $`|\chi (t=0)>=|\psi _{i_0}>`$ corresponding to the quantum register state $`i_0`$, then with time the probability will spread over the register and at a time $`t`$ the projection probability on the register state $`|\psi _i>`$ will be:
$$\begin{array}{c}F_{ii_0}(t)=|<\psi _i|\chi (t)>|^2\hfill \\ =_{m,m^{}}A_{im}A_{i_0m}^{}A_{im^{}}^{}A_{i_0m^{}}\mathrm{exp}(i(E_m^{}E_m)t),\hfill \end{array}$$
(6)
where $`A_{im}=<\psi _i|\varphi _m>`$ and $`E_m`$ is the energy of the stationary state $`|\varphi _m>`$ and we chose $`\mathrm{}=1`$. For $`JJ_c`$, the probability $`F_{i_0i_0}(t)`$ is very close to one for all times since the states are not mixed by the interaction. This means that all quantum register states $`|\psi _i>`$ remain well defined, and the computer can operate properly. For $`JJ_{cs}`$, only few states $`|\psi _i>`$ are mixed by the interaction, and $`F_{i_0i_0}(t)`$ oscillates in time regularly around an average value of order $`1/2`$. These oscillations are similar to the Rabi oscillations between two levels with frequency $`\mathrm{\Omega }J`$. An example is presented in Fig.13. In this regime, we expect that error-correcting codes may efficiently correct the spreading over few quantum register states.
For $`J>J_{cs}`$, quantum chaos sets in, and with time the probability spreads over more and more quantum register states until a quasi-stationary regime is reached where an exponentially large number of states are mixed. The probability $`F_{i_0i_0}(t)`$ drops approximately to zero, as shown on Fig.14. The chaotic time scale for this decay $`\tau _\chi `$ can be estimated as $`\tau _\chi 1/\mathrm{\Gamma }`$ where $`\mathrm{\Gamma }`$ is the width determined in the previous section. This estimate is very natural in the Fermi golden rule regime, with Breit-Wigner local density of states (4), since $`F_{i_0i_0}(t)`$ is essentially the Fourier transform of the local density of states $`\rho _W`$, and therefore decreases as $`\mathrm{exp}(\mathrm{\Gamma }t)`$. We note that the decay in this regime was recently discussed in . According to our data, when $`\mathrm{\Gamma }`$ becomes comparable to the energy bandwidth $`\sqrt{n}\delta `$, $`\rho _W`$ is close to a Gaussian distribution of width $`\mathrm{\Gamma }`$, and its Fourier transform $`F_{i_0i_0}(t)`$ is also a Gaussian of width $`1/\mathrm{\Gamma }`$. Therefore in both regimes we expect the time scale $`\tau _\chi `$ for the decay of $`F_{i_0i_0}(t)`$ to be $`\tau _\chi 1/\mathrm{\Gamma }`$. The data shown on Fig.14 correspond to the saturation regime for large values of $`n`$, and the insert shows that $`\tau _\chi 1/\mathrm{\Gamma }`$ is still valid. In fact the curve for $`n=16`$ in Fig.14 is already close to the limiting decay curve at $`\delta =0`$ (data not shown).
At the same time scale $`\tau _\chi `$ the quantum entropy $`S(t)`$ is large but is still growing. It reaches its maximal value on a larger time scale which seems independent of $`n`$. At this stage, an initial quantum register state is now spread over most of the register (Here $`S(t)=_iF_{ii_0}(t)\mathrm{log}_2F_{ii_0}(t)`$). This process is shown on Fig.15. This maximal value of $`S(t)`$ is approximately given by $`S_q`$ (see Fig.6) and accordingly decreases with decreasing $`J`$ as is illustrated in Fig.16.
Fig.17 illustrates this mixing process in the quantum register representation, evolving in time. The quantum computer hardware becomes quickly destroyed due to the inter-qubit coupling. It is necessary to decrease the coupling strength below the quantum chaos border to get well-defined quantum register states for $`t>0`$, as is illustrated in Fig.18.
The obtained data clearly show that exponentially many quantum register states become mixed after the finite chaotic time scale $`\tau _\chi 1/\mathrm{\Gamma }`$.
## V conclusions
The results presented in this paper show that residual inter-qubit coupling can lead to quantum chaos and ergodic very complicated eigenstates of the quantum computer. We have shown that in this regime the quantum register states disintegrate quickly in time over an exponentially large number of states and the computer operability is destroyed. We determined the dependence of the chaotic time scale $`\tau _\chi `$ of this process on coupling strength $`J`$, detuning fluctuations $`\delta `$ of one-qubit energy spacing, and number of qubits $`n`$. After this time $`\tau _\chi `$ the quantum computer hardware is melt. To prevent this melting one needs to introduce an efficient error-correcting code which operates on a time scale much shorter than $`\tau _\chi `$ and suppresses the development of quantum chaos. To avoid the quantum chaos regime dangerous for quantum computing, one should engineer the quantum computer in the integrable regime below the quantum chaos border $`J_c3\delta /n`$. It is important to note that this border decreases with the detuning $`\delta `$, showing that imperfections do not all conspire against the operability of the computer. We stress again that the transition to quantum chaos is an internal process which happens in a perfectly isolated system with no coupling to external world. Nevertheless, since a decoherence can be viewed as a result of internal interactions in a larger system, the results presented here may also apply to this problem.
Our main conclusion is that although in the quantum chaos regime a quantum computer cannot operate for long, fortunately the border for this process happens to be exponentially larger than the spacing between adjacent computer eigenstates, and therefore a broad parameter region remains available for realization of a quantum computer. Another possibility is to operate the quantum computer in the regime of quantum chaos. However, here one should keep in mind that after the chaotic time scale $`\tau _\chi `$ the computer hardware will melt due to inter-qubit coupling and quantum chaos. Therefor, the computer operability in this regime is possible only if many gate operations can be realized during the finite time $`\tau _\chi `$ (in a sense it becomes similar to the decoherence time). It is clear that the most preferable regime corresponds to quantum computer operation below the quantum chaos border.
We thank O.P. Sushkov for stimulating discussions, and the IDRIS in Orsay and the CICT in Toulouse for access to their supercomputers. One of us (DLS) thanks the Gordon Godfrey foundation at the University of New South Wales at Sydney for the hospitality at the final stage of this work. This research is partially done in the frame of EC program RTN1-1999-00400.
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# Radial velocities of pulsating subdwarf B stars: KPD 2109+4401 and PB 8783 Based on observations obtained with the William Herschel Telescope, ING, La Palma, Spain
## 1 Introduction
The discovery of small-amplitude non-radial pulsations in a number of subdwarf B stars (sdBVs) has introduced a powerful new tool for the study of stellar remnants (Kilkenny et al. 1997=SDBV I and subsequent papers). Long time-series photometric campaigns have detected rich spectra of oscillations in over a dozen targets, with frequencies and amplitudes generally indicative of low-order ($`\mathrm{}=02`$) and low-degree modes. The discovery has revolutionized the study of subdwarf B stars for the simple reason that pulsations in these stars were not expected. From a theoretical point of view, it appears that the pulsation mechanism is only effective when diffusion processes modify the outer layers of the star. Metal-enrichment in a specific layer must conspire with the local temperature to drive pulsations through the opacity mechanism (Charpinet et al. 1997).
In terms of their effective temperature and surface gravity, sdBVs are indistinguishable from their non-variable counterparts – there is no instability strip in which all sdBs pulsate. The origin of all sdBs remains a puzzle, although the increasing detection of sdB binarity, including several sdBVs, may point to a previous phase of common-envelope evolution in many, if not all. The structure of sdBs is partially hidden as a consequence of atmospheric diffusion, which disguises the true composition and the mass of the hydrogen-rich envelope, both of which are key diagnostics of previous evolution. By enabling an exploration of the composition and mass of these outer layers, asteroseismology may be of a pivotal importance.
Immediately after the discovery of sdBV pulsations, we recognized the potential for spectroscopy to provide additional diagnostics. Mode identification from photometry alone is challenging, and could be assisted by the identification of line-profile variations, whilst the comparison of radial and light amplitudes could be used to determine stellar radii directly. Spectroscopy might also demonstrate the presence of higher-order modes not detected photometrically. The limitations are that the periods are short (100-500s) compared with conventional CCD readout times, the stars are faint (12–15 mag.), and the photometric amplitudes are typically only a few tenths of one per cent. High-resolution high-S/N multi-line studies such as those obtained for non-radial oscillations in rapidly rotating bright O and B stars (e.g. Reid et al. 1993, Telting, Aerts & Mathias 1997) would not appear to be feasible. However, the development of new techniques offered the possibility to acquire high-speed spectroscopy of sdBVs and here we report our first successful observations obtained in 1998 October. Subsequently, O’Toole et al. (2000) announced preliminary results of independent radial velocity observations. They report a 9 $`\text{km}\text{s}^1`$ amplitude periodic variation in the large amplitude sdBV PG1605+072 at the principal frequency of 2.10 mHz found photometrically by Koen et al. (1998, SDBV VII).
## 2 Observations
The targets selected for our initial study were PB 8783 (Koen et al. 1997 = SDBV II) and KPD 2109+4401 (Koen et al. 1998 = SDBV XI, Billères et al. 1998). Although not ideal, our observations were scheduled when possibly more exciting targets such as PG1336–018 (Kilkenny et al. 1998 = SDBV VIII) and PG1605+072 (SDBV VII) were not accessible. On the other hand, both targets have been well-studied photometrically (O’Donoghue et al 1998 = SDBV V, SDBV XI, Billères et al. 1998) and show multiple periods in the ranges 7–10 and 5–6 mHz respectively, with amplitudes of a few millimagnitudes.
Observations were obtained with the 4.2m William Herschel Telescope on 1998 October 3 and 4 using the blue arm of the intermediate dispersion spectrograph ISIS. The R1200B grating and the TEK1 CCD with $`1124^2`$ 24$`\mu `$m pixels yielded an instrumental resolution (2 pixel) $`R=\mathrm{5\hspace{0.17em}000}`$ in the wavelength interval 4020–4420 Å. This wavelength region was chosen because it contains two strong Balmer lines and, potentially, a number of neutral helium and minor species lines that are normally observed in early-B stars. It also maximises the photon collection rate.
The CCD was read out in low-smear drift mode (Rutten et al. 1997), in which only a small number $`j`$ of CCD rows (parallel to the dispersion direction) are read out at one time. A dekker is used to limit the slit-length, thus only a fraction of the CCD window, $`j`$ rows, is exposed at one time. After exposing for an interval $`t`$, the CCD contents are stepped down by $`j`$ rows. This charge transfer imposes a small dead time $`d`$ which is roughly proportional to the CCD window size. Each set of $`j`$ rows is accumulated into a data cube containing $`n`$ individual 2D spectra, stacked adjacent to one another. In practice the first few spectra in the frame are null, corresponding to the unexposed part of the CCD being read out before exposed rows reach the CCD edge. Comparison lamps and stellar spectra can be obtained in exactly the same way. Maximum frame sizes limit $`n`$ to a few hundred; we normally adopted $`n=200`$ before beginning a new frame.
In our experiment, $`j`$ was set just wide enough to admit a small region of sky either side of the stellar image and to allow for edge effects between adjacent spectra. On-chip binning by a factor of 3 further reduced the read-out time and data-volume. $`t`$ was chosen so that $`t+d`$ was approximately 1/10 of the typical observed pulsation period. Values for $`t,j`$ and $`\mathrm{\Sigma }n`$ are given for each target in Table 1, together with the average number of counts detected per wavelength resolution element $`c`$, and the maximum S/N ratio $`s`$ anticipated in each exposure from nominal performance figures (SIGNAL, Benn 1997). In this configuration, $`d=1.57`$s and the read-out noise was 5.24 electrons per pixel.
Due to occasional errors in the CCD controller, not all exposure times $`t`$ were equal, roughly one in 100 exposures would drop by an unpredictable amount. These were identified by integrating the total counts in each individual spectrum and adjusting individual exposure times to reflect the photometry.
The data cubes were reduced using scripted standard IRAF routines. After bias subtraction and flatfielding, the 1-dimensional spectra were optimally extracted and sky subtracted. The copper-argon arc spectra obtained at the beginning and end of each data cube were time weighted for each stellar observation and the corrected calibration derived and applied. Over the duty cycle of one data cube the arc shifts were never greater than a few microns. Further reductions included normalization of the spectra and cosmic ray removal.
Cross-correlation of the arc spectra showed that the maximum shift during a half hour period was around 7 microns (about a quarter of a pixel). Arc exposures showed that this was varying smoothly during the night and as calibration data were obtained before and after each stellar data cube, then simple spectral interpolation of the extracted one dimensional spectra (weighted by the time interval from the bracketing arc spectra) removed any instrumental shift to probably within 1 micron ($`<0.9\text{km}\text{s}^1`$).
The final data product is a set of files each containing $`n`$ 1-dimensional wavelength-calibrated flux-normalized spectra. Each spectrum is time-tagged. The mean spectrum for each of our targets is shown in Fig. 1
## 3 Analysis and Results
The primary objective of our observations was to detect and measure radial motion due to the dominant oscillatory modes detected from photometry. There was the additional possibility that line-profile variations or multiple modes might be detected, as well as any mutual motion within any binary system.
The most commonly adopted method for measuring radial-velocity shifts in stellar spectra is the cross-correlation function (ccf). We constructed a mean spectrum for each target. This template was cleaned to remove stationary features such as bad CCD columns. Both template and individual spectra were continuum-subtracted; they were already linearized to the same wavelength scales. The ccf for each spectrum relative to the template was computed and stored as a 2d function $`\chi (t,v)`$, where $`t`$ is now time, $`\delta t`$ is the mean sampling interval and $`v`$ is radial velocity relative to the template. The velocity of each spectrum was found by a) fitting a gaussian to, b) fitting a parabola to and c) computing the centroid of $`\chi (v)`$ in the range $`\delta v`$ to $`+\delta v`$, to give functions $`v_g(t),v_p(t),`$ and $`v_c(t)`$, respectively. Moments $`M_j(t)=(vv_{\mathrm{ref}})^j\chi (v,t)𝑑v;j=1,2,3`$ were also computed. With $`v_{\mathrm{ref}}=0`$, $`M_1`$ corresponds to the centroid. $`v(t)`$ was sensitive both to the method chosen and the value of $`\delta v`$.
The functions $`v(t)`$ were investigated for periodic content by means of the discrete Fourier transform $``$. Either a low-degree polynomial or a long-period sine function $`b(t)`$ was first fit and subtracted from $`v(t)`$. $`\{v(t)b(t)\}`$ was then computed.
The velocity data $`v_p(t)`$ and $``$ are shown in Figs. 2 and 3. These figures also show a representation of the principal frequencies and amplitudes discovered photometrically $`𝒫`$ (SDBV V,IX) <sup>1</sup><sup>1</sup>1Throughout this paper, the term amplitude applied to periodic signals refers to the semi-amplitude $`a`$ of the sine function $`a\mathrm{sin}(2\pi \nu t+\varphi )`$.
The information we require from the amplitude spectrum is not in the first instance the significance of the peaks but rather the amplitude associated with expected frequencies. Over the entire range $`0.12/\mathrm{\Sigma }n\delta t<\nu /\mathrm{mHz}<1/\delta t90`$, $``$ shows many peaks of comparable amplitude which on their own have little statistical significance. However the coincidence in Figs. 2 and 3 of the highest peaks in $``$ with the highest peaks in $`𝒫`$ is more remarkable and is regarded as a positive signature of oscillatory behaviour in the spectroscopic data. The resolution $`2/n\delta t0.1`$mHz in $``$ is insufficient to resolve the fine structure observed in much longer photometric time series.
### 3.1 KPD 2109+4401
Because its spectrum is not contaminated by that of a late-type companion, the results for KPD 2109+4401 are the less complicated to interpret. A background drift of approximately 9 $`\text{km}\text{s}^1`$ over 5 hours is approximately linear and defines $`b`$. We adopted $`\delta v=150\text{km}\text{s}^1`$ to establish $`v_p(t)`$ and to construct Fig. 2. Principal peaks in the amplitude spectrum are evident at $`\nu _1=5.09`$ mHz and $`\nu _2=5.48`$ mHz.
The same frequency structure in $``$ was obtained for $`\delta v=150,300`$ and 450 $`\text{km}\text{s}^1`$ when either a gaussian, parabolic fit or centroid is used to measure $`v(t)`$. In the interval 4.5–6.4 mHz, minor peaks occured at $`\nu _{37}=4.77`$, 5.24, 5.33, 5.90 and 6.03 mHz in most of these analyses. Beating may be responsible for some peaks; for example $`\nu _3=\nu _1(\nu _2\nu _1),\nu _6=\nu _2+(\nu _2\nu _1)`$. Thus only $`\nu _1`$ and $`\nu _2`$ could be identified as genuine from the present data.
The probability of finding a set of peaks at prespecified frequencies and exceeding a given amplitude threshold by chance may be computed as follows. For a given frequency range $`\mathrm{\Delta }v`$, count the number of peaks $`n`$ which exceed a given threshold $`v_{\mathrm{thresh}}`$. Then count the number of peaks $`n_{\mathrm{id}}`$ which lie within a resolution width $`\pm \delta v`$ of the predicted frequencies. The probability of this number occurring by chance is
$$p=\left(n\frac{2\delta v}{\mathrm{\Delta }v}\right)^{n_{\mathrm{id}}}.$$
For KPD 2109+4401, we chose $`\mathrm{\Delta }v=6`$ mHz, $`v_{\mathrm{thresh}}=1.5\text{km}\text{s}^1`$ and $`\delta v=0.05`$ mHz which gave $`n=3`$, $`n_{\mathrm{id}}=2`$ and hence $`p=0.0025`$.
In order to interpret the amplitudes of the identified peaks, we carried out the following numerical experiment (table 3). We first formed the sum of a series of sine functions each having a frequency $`\nu _\mathrm{i}`$, amplitude $`a_\mathrm{i}`$ and a random phase shift $`\varphi _\mathrm{i}`$. The $`\nu _\mathrm{i}`$ and $`a_\mathrm{i}`$ were based on the photometric results for KPD 2109+4401 (SDBV IX). This sum was sampled at the same times $`t_\mathrm{j}`$ as the observations,
$$s(t_\mathrm{j})=\mathrm{\Sigma }_\mathrm{i}a_\mathrm{i}\mathrm{sin}(2\pi \nu _\mathrm{i}t_\mathrm{j}+\varphi _\mathrm{i}).$$
To this we added normally-distributed noise with a standard deviation $`\sigma =15\text{km}\text{s}^1`$ similar to the observational scatter. $``$ was derived and usually found to contain peaks at the principal input frequencies, but not always with the same amplitude ratios. Minor peaks were less frequently identified. The amplitudes of the highest peaks within $`\pm 0.1`$mHz of the test frequencies were measured. This procedure was repeated $`n`$ times and the average amplitudes for each detected frequency $`a_\mathrm{i}^{}`$ were computed. In the case of KPD 2109+4401, the spacing of some input frequencies is less than the resolution provided by the limited time series. Consequently, the measured amplitudes represent some sum from several independent sine functions. Interference, or beating, means that the apparent amplitude of such a signal will vary according to the relative phases of the components. This is apparent in table 3 (six-frequency simulation), where the relative standard deviation of the amplitude of the composite signal at 5.09 mHz is much higher than that of the single-frequency signal at 5.21 mHz. The experiment was repeated with only two well-resolved input frequencies and yielded the ratio $`a^{}/a1.0\pm 0.1`$. Reducing $`\sigma `$ reduces the error on $`a^{}/a`$ due to experimental noise, but not that due to mode beating or to data sampling. It was also repeated for $`\nu `$ and $`a`$ appropriate to the second target PB 8783 with a similar result.
Therefore, we conclude that the peak amplitudes of 2.49 and 2.68 $`\text{km}\text{s}^1`$ at frequencies 5.09 and 5.48 mHz in KPD 2109+4401 reflect to within $`20\%`$ (standard deviation) the combined real amplitudes of oscillations within $`0.05`$mHz of these frequencies (Table 2). The principal causes of uncertainty in the amplitudes are beating between unresolved modes and experimental noise.
The 9 $`\text{km}\text{s}^1`$ drift in $`b`$ may be real. Instrumental drift is limited to $`<0.9\text{km}\text{s}^1`$ (section 2). The heliocentric correction, which was not applied to the data, changes by less than 0.4 $`\text{km}\text{s}^1`$. In the simulations described above, the r.m.s. amplitude of $`b`$ was 1.7 $`\text{km}\text{s}^1`$. Although KPD 2109+4401 is not known to be a binary, it cannot be ruled out. PG 1336–018 was only recognised as a binary from its eclipses; the secondary is too faint to be detected spectroscopically (SDBV VIII).
### 3.2 Line profile variations?
In addition to radial-velocity information, $`\chi `$ contains information about line-profile variations, since the ccf represents an average line profile within the region of spectrum analyzed. In the present case, the average line profile is dominated by two Balmer lines, H$`\gamma `$ and H$`\delta `$. The average line width and asymmetry are indicated either by the moments $`M_{2,3}`$, respectively, or by the coefficients $`g_{2,4}`$ in the gaussian + linear fit $`g(v)=g_0\mathrm{exp}((vg_1)^2/g_2)+g_3+g_4v`$.
Neither the Fourier transform $`\{M_2(t)\}`$ or $`\{g_2(t)\}`$ showed peaks near $`\nu _1`$ or $`\nu _2`$. With suitably chosen $`\delta v`$, both $`\{M_3(t)\}`$ and $`\{g_4(t)\}`$ showed frequency structure similar to $`\{v(t)\}`$. However, the amplitude of peaks in $`\{M_1(t)\}`$ only approached that of $`\{v_p(t)\}`$ if $`\delta v`$ was sufficiently large, in which case the periodic content of $`M_3(t)`$ was essentially lost. This implied that $`M_3`$ contained velocity information not properly incorporated into $`M_1`$.
$`\{g_4(t)\}`$ showed peaks at 5.09 and 5.55 mHz with $`\delta v400600\text{km}\text{s}^1`$, but not otherwise. In the case of intrinsically sharp lines, such asymmetry may reflect differential radial motion across the rotationally broadened line profile. In the present study, any rotational broadening is convolved with Stark broadening in the Balmer lines and is difficult to interpret. Consequently, whilst we suspect that there may be evidence for line-profile variations with similar frequencies to the radial motions, we can draw no firm conclusions from the present data.
### 3.3 PB 8783
The binary PB 8783 is altogether more difficult to analyze because of the superposition of sdB and F star spectra. The results shown in Fig. 3 are based on fitting parabola to $`\chi (v)`$ with $`\delta v=420\text{km}\text{s}^1`$. In deriving the amplitude spectrum in Fig. 3, the background function $`b(t)`$ was found to be best given by a sinusoid with an amplitude of 10 $`\text{km}\text{s}^1`$, representing a systemic drift substantially larger than that in KPD 2109+4401.
On inspection, the typical ccf for PB 8783 consists of a strong broad component (FWHM $`1400\text{km}\text{s}^1`$) superimposed by a small narrow component (FWHM $`200\text{km}\text{s}^1`$, Fig. 4). The narrow feature disappears when the template is replaced by a theoretical sdB star spectrum containing H and He lines only, and is deduced to be due principally to metal-lines in the F-type secondary. This is confirmed when $`\chi `$ is calculated using only wavelengths between H$`\delta `$ and H$`\gamma `$; the broad component disappears leaving a small narrow peak.
The detection of two components in $`\chi `$ offers the possibility to resolve the radial motions of the binary components. The mutual motion of the F-star is clearly apparent when $`\chi `$ excludes the Balmer lines, and appears as a change of $`10\text{km}\text{s}^1`$ during the five-hour observing run. However, the sign of the change remains the same when the theoretical sdB spectrum is used as a template, indicating that the Balmer lines in the F-star contribute significantly to $`\chi `$.
The problem was resolved by fitting $`\chi (\delta v:+\delta v),\delta v=280\text{km}\text{s}^1`$ with the function $`g(v)=g_0\mathrm{exp}((vg_1)^2/g_2)+g_3+g_4v+g_5v^2`$, where the quadratic terms fit the broad component and the gaussian terms represent the narrow component (Fig. 4). The functions $`g_1(v)`$ and $`g_4(v)/2g_5(v)`$ were examined to determine the mutual motion of each star and to look for high-frequency content. The results are shown in Figs. 5 and 6 where it will be seen that the slowly varying components are of opposite sign. It is suggested that, in the absence of viable alternatives for these systematic drifts, these slowly varying components are most probably due to the mutual motion of the two stars within the binary.
If this is the case, a linear regression on the first 4 hours of both sets of velocity data gives the relative acceleration of the two components. The result, $`\mathrm{\Delta }v_\mathrm{F}/\mathrm{\Delta }v_{\mathrm{sdB}}=0.81\pm 0.10`$, provides an estimate for the mass ratio of the two stars $`M_{\mathrm{sdB}}/M_\mathrm{F}0.8`$. This is a little high for an F-star on the main-sequence ($`1.11.8\text{ }\mathrm{M}_{}`$) and a canonical sdB star ($`0.5\text{ }\mathrm{M}_{}`$), and suggests that the relative motions have still not been fully resolved. Further information concerning the velocities of the two components relative to one another is convolved into the composite Balmer profiles in the ccf template and will require observations of the complete binary orbit to disentangle. Whilst the binary period is clearly longer than the data sequence, it is constrained by the relation
$$K_\mathrm{F}=210\frac{1}{1+q}\left(\frac{M}{P}\right)^{1/3}\mathrm{sin}i\text{km}\text{s}^1.$$
(1)
$`M`$ is the total mass of the system ($`\mathrm{M}_{}`$), $`q=M_\mathrm{F}/M_{\mathrm{sdB}}`$ is the mass ratio, $`P`$ is the orbital period (d) and $`K_\mathrm{F}`$ is the velocity semi-amplitude of the F star ( $`\text{km}\text{s}^1`$). For example, with $`M_{\mathrm{sdB}}=0.55\text{ }\mathrm{M}_{}`$, $`M_\mathrm{F}=1.2\text{ }\mathrm{M}_{}`$, $`i=30^{}`$ and $`P=2`$d, $`K_\mathrm{F}=32\text{km}\text{s}^1`$. The maximum change in velocity over the timescale of our observations (0.15d) would be $`\mathrm{\Delta }v_\mathrm{F}15\text{km}\text{s}^1`$, compared with 13 $`\text{km}\text{s}^1`$ observed. Inverting the procedure, we can determine the maximum $`P`$ as a function of $`i`$ required to obtain $`\mathrm{\Delta }v_\mathrm{F}=13\text{km}\text{s}^1`$ in 0.15d. For $`M_\mathrm{F}=1.21.8\text{ }\mathrm{M}_{}`$, the maximum $`P`$ increases from $`1.00.8`$d at $`i=10^{}`$ to $`P=3.73.2`$d at $`i=80^{}`$. Eclipses are not expected at inclinations $`<80^{}`$.
It is satisfying to detect the photometric frequencies already reported for PB8783 (SDB V) in the deconvolved radial velocities of the sdB star (Fig. 5) and to find no such signal in the radial velocities of the F star (Fig. 6). This provides confirmation that the photometric and radial velocity variations are due to pulsations in the sdB star and not in the cool companion. A similar result was obtained for PG1336-018 (SDBV VIII), except in that case the persistence of oscillations during an eclipse of the unseen companion provided the necessary proof.
From this point, the interpretation of the sdB star velocities is the same as that for KPD2209+4101. Peaks in the velocity amplitude spectrum were identified at five frequencies also present in the photometry. The resulting pulsation velocity amplitudes were thus found to be in the range $`1.12.2\text{km}\text{s}^1`$ and are shown in Table 4. Regarding the probability that these frequencies were identified by chance and following the prescription already given for KPD 2209+4101, we found with $`v_{\mathrm{thresh}}=1.5\text{km}\text{s}^1`$ that $`n=9,n_{\mathrm{id}}=3`$ and $`p=0.0034`$ (Fig. 5). Relaxing $`v_{\mathrm{thresh}}=1.0\text{km}\text{s}^1`$ yielded $`n=27,n_{\mathrm{id}}=6`$ and $`p=0.0083`$.
### 3.4 Relating luminosity and velocity amplitudes
Kjeldsen & Bedding (1995) investigated the relation between velocity and luminosity amplitudes for several classes of main-sequence, giant and supergiant pulsators and derived the calibrated relation
$$(\delta L/L)_\lambda =\frac{v_{\mathrm{osc}}/\mathrm{m}\mathrm{s}^1}{(\lambda /550\mathrm{n}\mathrm{m})(\text{ }\text{T}\text{eff}/5777K)^2}20.1\mathrm{ppm},$$
(2)
where ppm denotes parts-per-million.
We tested whether these predictions could be extended to the sdBVs since we have both light ($`\delta L/L`$) and velocity ($`v_{\mathrm{osc}}`$) amplitudes for the same oscillation frequencies. We found that the observed velocity amplitudes are $`>2`$ times those expected from the photometry and Eqn. (2). However Eqn. (2) has been derived from cool stars pulsating in radial or non-radial p-modes and it may not be applicable to pulsations in the hot sdBVs. We note that the amplitudes of individual modes in sdBVs have been reported to vary from one season and possibly one night to another (SDBV/,VIII). Beating between unresolved modes will certainly affect the velocity amplitude measurements. Simultaneous high-speed photometry and spectrocopy will be required to properly relate the velocity and luminosity amplitudes and to make any futher validation of Eqn. 2.
### 3.5 For future investigation
We have left two issues for future investigation.
1) Compared with its apparent surface gravity, the observed pulsation frequencies of PB 8783 are too high (SDBV XII). The expectation is that its surface gravity ($`\mathrm{log}g=5.55`$, O’Donoghue et al. 1997, SDBV IV) should be much higher. It is recognised that this could be due to inadequate subtraction of the F-star spectrum, a conjecture which should be checked using higher quality data, including those presented here.
2) If the mode and amplitude of a spherical harmonic responsible for each pulsation frequency is specified, then it is relatively straightforward to integrate the projected surface velocities to provide a theoretical radial velocity curve. It is slightly more complicated to construct the theoretical light curve. From the combination it should be possible to refine the mode and amplitude measurements already established.
## 4 Conclusions
High-speed high-resolution spectroscopy of two non-radially pulsating subdwarf B stars KPD 2109+4401 and PB 8783 has been acquired with the William Herschel telescope. These data show radial velocity variations at both high- and low-frequencies. High-frequency peaks in the velocity amplitude spectrum correspond to frequencies identified in the light curves of both stars, and have allowed the velocity amplitudes associated with these stellar oscillations to be estimated. Typically $`2\text{km}\text{s}^1`$, these translate into radial variations of some $`100250`$km within 60–90s. The prospect for measuring the velocity amplitudes of non-radial pulsations in sdBVs with much higher amplitudes is promising. Line profile variations were not detected in the data although there are good reasons to suppose they should be present. The possibility that frequencies not present in the photometry may contribute to the velocity amplitude spectrum should also be pursued.
In the case of the binary sdB PB 8783, low-frequency velocity variations corresponding to the mutual motion of both components have been resolved. These provide an upper limit to the orbital period of between 0.8 and 3d, depending principally on the orbital inclination. Establishing the orbital period for this star should therefore be a priority for future observations. The deconvolution of component star velocities also demonstrated that high-frequency velocity variations occur only in the sdB star and not in the F star.
## Acknowledgments
This research has been supported by the Department of Education in Northern Ireland through a grant to the Armagh Observatory.
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# The halo of the exotic nucleus 11Li: a single Cooper pair
## Abstract
If neutrons are progressively added to a normal nucleus, the Pauli principle forces them into states of higher momentum. When the core becomes neutron-saturated, the nucleus expels most of the wavefunction of the last neutrons outside to form a halo, which because of its large size can have lower momentum. It is an open question how nature stabilizes such a fragile system and provides the glue needed to bind the halo neutrons to the core. Here we show that this problem is similar to that of the instability of the normal state of an electron system at zero temperature solved by Cooper, solution which is at the basis of BCS theory of superconductivity. By mimicking this approach using, aside from the bare nucleon-nucleon interaction, the long wavelength vibrations of the nucleus <sup>11</sup>Li, the paradigm of halo nuclei, as tailored glues of the least bound neutrons, we are able to obtain a unified and quantitative picture of the observed properties of <sup>11</sup>Li.
Aside from “dark matter” , atomic nuclei, little droplets made out of protons and neutrons 10 femtometer across , make up a sizable fraction of the present mass of the universe. Research into the structure of atomic nuclei concentrates largely on the limits of the nuclear stability , where new physics is expected to be, in particular at the limits of neutron and proton number defining the so called drip lines in the chart of nuclides. The most exotic nuclei , first produced in the laboratory only few years ago, are those that lie just within the drip lines, on the edges of nuclear stability. Of these, the atomic nucleus $`{}_{3}{}^{}{}_{}{}^{11}`$Li<sub>8</sub>, containing 3 protons and 8 neutrons, is rightly the most famous and better studied and provides the cleanest example of halo nuclei to date .
In halo nuclei, some of the constituent neutrons or protons venture beyond the drop’s surface and form a misty cloud or halo. Not surprisingly, these extended nuclei behave very differently from ordinary (“normal”) nuclei lying along the so-called stability valley in the chart of nuclides. In particular, they are larger than normal nuclei of the same mass number, and they interact with them with larger cross sections as well. In the case of <sup>11</sup>Li, the last two neutrons are very weakly bound. Consequently, these neutrons need very little energy to move away from the nucleus. There they can remain in their “stratospheric” orbits, spreading out and forming a tenuous halo. If one neutron is taken away from <sup>11</sup>Li, a second neutron will come out immediately, leaving behind the core of the system, the ordinary nucleus <sup>9</sup>Li. This result testifies to the fact that pairing, the attraction correlating pairs between the least bound particles in a system, plays a central role in the stability of <sup>11</sup>Li.
It is well known that pairing can radically affect the properties of a many-body system. In metals, pairing between the electrons gives rise to superconductivity . In a dilute neutron gas, pairing can influence the properties of neutron stars, cold remnants of fierce supernova explosions . Acting on a liquid made out of atoms of the lighter isotope of helium ($`{}_{2}{}^{}{}_{}{}^{3}`$He<sub>1</sub>) it leads to superfluidity . In nuclei, it controls almost every aspect of nuclear structure close to the ground state and determines, to a large extent, which nuclei are stable and which are not .
The basic experimental facts which characterize <sup>11</sup>Li and which are also of particular relevance in connection with pairing in this system are: a) $`{}_{3}{}^{}{}_{}{}^{9}`$Li<sub>6</sub> and $`{}_{3}{}^{}{}_{}{}^{11}`$Li<sub>8</sub> are stable, $`{}_{3}{}^{}{}_{}{}^{10}`$Li<sub>7</sub> is not, b) the two-neutron separation energy in <sup>11</sup>Li is only $`S_{2n}`$=0.294$`\pm `$0.03 MeV as compared with values of 10 to 30 MeV in stable nuclei, c) <sup>10</sup>Li displays s- and p-wave resonances at low energy, their centroids lying within the energy range 0.1-0.25 MeV and 0.5-0.6 MeV respectively , d) the mean square radius of <sup>11</sup>Li, $`r^2^{1/2}`$=3.55$`\pm `$ 0.10 fm , is very large as compared to the value 2.32$`\pm `$0.02 fm of the <sup>9</sup>Li core, and testifies to the fact that the neutron halo must have a large radius ($``$6-7 fm), e) the momentum distribution of the halo neutrons is found to be exceedingly narrow, its FWHM being equal to $`\sigma _{}=48\pm 10`$ MeV/c for the (perpendicular) distribution observed in the case of the break up of <sup>11</sup>Li on <sup>12</sup>C, a value which is of the order of one fifth of that measured during the break up of normal nuclei , f) the ground state of <sup>11</sup>Li is a mixture of configurations where the two halo nucleons move around the <sup>9</sup>Li core in $`s^2`$ and $`p^2`$configurations with almost equal weight .
Before discussing the sources of pairing correlations in <sup>11</sup>Li, we shall study the single-particle resonant spectrum of <sup>10</sup>Li. The basis of (bare) single-particle states used was determined by calculating the eigenvalues and eigenfunctions of a nucleon moving in a Saxon-Woods potential with spin-orbit and symmetry term (cf. , Vol. I, Eqs. (2.281) and (2.282)). The continuum states of this potential were calculated by solving the problem in a box of radius equal to 40 fm, chosen so as to make the results associated with <sup>10</sup>Li and <sup>11</sup>Li discussed below, stable. While mean field theory predicts the orbital $`p_{1/2}`$ to be lower than the $`s_{1/2}`$ orbital (cf. Fig. 1, I(a)), experimentally the situation is reversed. Similar parity inversions have been observed in other isotones of $`{}_{}{}^{10}{}_{3}{}^{}`$Li<sub>7</sub>, like e.g. $`{}_{4}{}^{}{}_{}{}^{11}`$Be<sub>7</sub>. Shell model calculations testify to the fact that the effect of core excitation, in particular of quadrupole type, play a central role in this inversion ,(cf. also ). In keeping with this result, we have studied the effect the coupling of the $`p_{1/2}`$ and $`s_{1/2}`$ orbitals of <sup>10</sup>Li to monopole-, dipole- and quadrupole-vibrations of the <sup>9</sup>Li core has on the properties of the 1/2<sup>+</sup> and 1/2<sup>-</sup> states of this system. In this study we have used vibrational states calculated by diagonalizing, in the random phase approximation (RPA), a multipole-multipole separable interaction taking into account the contributions arising from the excitation of particles into the continuum states. The coupling strengths were adjusted so as to reproduce the position and transition probabilities of vibrational states in the normal nucleus <sup>10</sup>Be . Unperturbed particle-hole excitations up to 70 MeV have been included in the calculations and phonon states up to 50 MeV have been considered. Within this space there are of the order of $`10^2`$ dipole states and about the same amount of quadrupole states, exhausting the associated energy-weighted sum rules. The resulting vibrational states were coupled to the particle states making use of the associated transition densities (formfactors) and particle-vibration coupling strengths. A Skyrme-type effective interaction (SLy4) was used to calculate the monopole linear response. The corresponding solutions were obtained in coordinate space making use of a mesh extending up to a radius of 80 fm. The monopole response exhausts 94% of the EWSR considering the summed contributions up to 40 MeV of excitation energy. This response function was discretized in bins of 300 keV and the resulting states coupled to the single-particle states making use of the corresponding transition densities and particle-vibration coupling strengths. Similar calculations were performed to determine the properties of the $`L`$=0,1 and 2-vibrational states of <sup>11</sup>Li needed in connection with the discussion of the pairing induced interaction carried out below. In this case the strength of the separable dipole interactions has been slightly changed from the value used in the previous calculation so as to provide an overall account of the experimental dipole response in <sup>11</sup>Li . For the quadrupole strength the same value used before was employed. Because the calculations have been carried out on the physical (correlated ) <sup>11</sup>Li ground state, the transitions associated with the vibrational states involving the $`p_{1/2}`$ and the $`s_{1/2}`$ states have been renormalized with the corresponding occupation numbers resulting from the full diagonalization.
In the study of the single-particle resonances of <sup>10</sup>Li we have considered not only the particle-vibration coupling vertices associated with effective-mass-like diagrams (upper-part graph of Fig. 1, I(b)) and leading to attractive (negative) contributions to the single-particle energies, but also those couplings leading to Pauli principle (repulsive) correction processes associated with diagrams containing two-particles, one-hole and a vibration in the intermediate states (lower-part diagram of Fig. 1, I(b)). Because of such Pauli correction processes, the $`p_{1/2}`$ state experiences an upward shift in energy, arising from the coupling of this orbital to the $`p_{3/2}`$ hole-state through quadrupole vibrational states, in keeping with the fact that the ($`p_{1/2}p_{3/2}^1`$) particle-hole excitation constitutes an important component of the quadrupole vibration wavefunction. As a consequence, the $`p_{1/2}`$ state becomes unbound, turning into a low-lying resonance with centroid $`E_{res}`$ 0.5 MeV. Due to the coupling to the vibrations the $`s`$states are instead shifted downwards. In fact, in this case there are essentially no (repulsive) contributions arising from the Pauli correction processes. On the other hand, (attractive) effective-mass-like processes with intermediate states consisting of one particle plus a vibrational state of the type ($`s_{1/2}\times 0^+`$), ($`p_{1/2}\times 1^{}`$) and ($`d_{5/2}\times 2^+`$) lead to a virtual state with $`E_{virt}=0.2`$ MeV (cf. Fig. 1, I(b)), the coupling to quadrupole vibrations providing by far the largest contribution to the total energy shift. The important difference between the distribution of the single-particle strength associated with the resonant state $`p_{1/2}`$ and the virtual state $`s_{1/2}`$ can be observed in Fig. 1, I(c), where the partial cross section $`\sigma _l`$ for neutron elastic scattering off <sup>9</sup>Li is shown. While $`\sigma _p`$ displays a clear peak at 0.5 MeV, $`\sigma _s`$ is a smoothly decreasing function of the energy. On the other hand, a small increase in the depth of the potential felt by the $`s`$neutron will lead to a (slightly) bound state, hence the name of virtual .
At the basis of the variety of pairing phenomena observed in systems so apparently different as atomic nuclei and metals is the formation of Cooper pairs . In the case of metals, Cooper pairs arise from the combined effect of Pauli principle and of the exchange of lattice vibrations (phonons) between pairs of electrons (fermions) moving in time reversal states lying close to the Fermi energy. In the superconducting or BCS state of metals , Cooper pairs are strongly overlapping. In fact, there are on average 10<sup>6</sup> pairs which have their centers of mass falling within the extent of a given pair wavefunction. In spite of the modest number of Cooper pairs present in the ground state of atomic nuclei ($``$ 10), BCS theory gives a quite accurate description of pairing in nuclei , the analogy between the pairing gap typical of metallic superconductors and of atomic nuclei being very much to the point . On the other hand, the finiteness of the nucleus introduces in the BCS treatment of pairing important modifications (quantal size effects (QSE), ,-). In particular, while in the infinite system the existence of a bound state of the (Cooper) pair happens for an arbitrarily weak interaction , in the nuclear case this phenomenon takes place only if the strength of the nucleon-nucleon potential is larger than a critical value connected with the discreteness of the nuclear spectrum. In fact, calculations carried out making use of a particularly successful parametrization of the (bare) potential (Argonne potential ), testify to the fact that the nuclear forces are able to bind Cooper pairs in open shell nuclei (leading to sizable pairing gaps (1-2 MeV) ), but not in closed shell nuclei, the most important contributions to the nucleon-nucleon (pairing) interaction arising from high multipole components of the force .
The situation is however quite different for the ”open shell” nucleus <sup>11</sup>Li. In fact, allowing the two neutrons to interact through the Argonne potential produces almost no mixing between $`s`$ and $`p`$waves, but essentially it only shifts the energy of the unperturbed (resonant) states $`s_{1/2}^2(0)`$ and $`p_{1/2}^2(0)`$ by about 80 keV without giving rise to a bound system. Similarly, making use of the matrix elements of the same nucleon-nucleon potential in connection with the BCS equations does not lead to a solution but the trivial one of zero pairing gap ($`\mathrm{\Delta }_\nu =0,U_\nu V_\nu `$ = 0) . At the basis of this negative result is the fact that the most important single-particle states allowed to the halo neutrons of <sup>11</sup>Li to correlate are the $`s_{1/2}`$, $`p_{1/2}`$ and $`d_{5/2}`$ orbitals. Consequently the two neutrons are not able, in this low-angular momentum phase space, to profit fully from the strong force-pairing interaction, as only the components of multipolarity $`L=0`$, 1 and 2 of this force are effective in <sup>11</sup>Li because of angular momentum and parity conservation .
In keeping with this result and with those of ref. , and making use of the fact that <sup>11</sup>Li displays low-lying collective vibrations , one can posit that the exchange of these vibrations between the two neutrons is the main source of pairing available to them to correlate. In fact, allowing the two neutrons to both exchange phonons (induced interaction, Fig. 1, II(a)), as well as to emit and later reabsorb them ( self-energy correction, Fig. 1, I(b)), leads to a bound (Cooper) pair, the lowest eigenstate of the associated secular matrix being $`E_{gs}=`$ -0.270 MeV. Adding to the induced interaction the nucleon-nucleon Argonne potential one obtains $`E_{gs}`$ = -0.330 MeV, and thus a two-neutron separation energy quite close to the experimental value. Measured from the unperturbed energy of a pair of neutrons in the lowest state calculated for <sup>10</sup>Li, namely the $`s`$-resonance ($`E_{unp}=2E_{s_{1/2}}=`$ 400 keV, cf. Fig. 1,I(b)), it leads to a pairing correlation energy $`E_o=E_{unp}E_{g.s.}=`$ 0.730 MeV (cf. Fig. 1, II(b)). From the associated two-particle ground state wavefunction $`\mathrm{\Psi }_0(\stackrel{}{r}_1,\stackrel{}{r}_2)(\stackrel{}{r}_1,\stackrel{}{r}_2|0^+)`$, one obtains a momentum distribution (whose FWHM is $`\sigma _{}=56`$ MeV/c, for <sup>11</sup>Li on <sup>12</sup>C) and ground state occupation probabilities of the two-particle states $`s_{1/2}^2`$(0), $`p_{1/2}^2`$(0) and $`d_{5/2}^2(0)`$ (0.40, 0.58 and 0.02 respectively, cf. Fig. 1,II(b)) which provide an overall account of the experimental findings. The radius of the associated single-particle distribution is 7.1 fm. Adding to this density that of the core nucleons one obtains the total density of <sup>11</sup>Li. The associated mean square radius (3.9 fm) is slightly larger than the experimental value.
The spatial structure of the Cooper pair described by the wavefunction $`\mathrm{\Psi }_0(\stackrel{}{r}_1,\stackrel{}{r}_2)`$ is displayed in Fig. 2. The mean square radius of the center of mass of the two neutrons is $`r_{cm}^2^{1/2}`$ = 5.4 fm. This result testifies to the importance the correlations have in collecting the small (enhanced) amplitudes of the uncorrelated two-particle configuration $`s_{1/2}^2(0)`$ in the region between 4 to 5 fm, region in which the $`p_{1/2}^2(0)`$, helped by the centrifugal barrier, displays a somewhat larger concentration (Fig. 3). From the above results, it emerges that the exchange of vibrations between the least bound neutrons leads to a (density-dependent) pairing interaction acting essentially only outside the core (cf. also ref. ). To be noted that the long wavelength behaviour of these vibrations is connected with the excitation of the neutron halo, the large size of which not only makes the system easily polarizable but provides also the elastic medium through which the loosely bound neutrons exchange vibrations with each other .
The average mean square distance between the halo neutrons is $`<r_{12}^2>^{1/2}`$ 9.2 fm, in keeping with the fact that the coherence length associated with Cooper pairs in nuclei is larger than the nuclear dimensions thus preventing the possibility of a nuclear supercurrent (cf. e.g. Vol. II, p. 398). On the other hand, this value of $`<r_{12}^2>`$ does not prevent the two correlated neutrons to be close together. The corresponding (small) probability (cf. Fig. 2) being much larger than that associated with the uncorrelated neutrons (cf. Fig. 3).
Similar results as those reported above are obtained solving the BCS equation for the two-neutron system making use of the matrix elements used in the diagonalization, sum of those of the nucleon-nucleon Argonne potential and those of the induced interaction. In this case, the correlation energy is $`E_o`$= 0.7 MeV, the separation energy of the two neutrons becoming $`S_{2n}=`$ 0.3 MeV. The radial structure of the projected BCS wavefunctions $`_{\nu >0}(V_\nu /U_\nu )\phi _\nu (\stackrel{}{r}_1)\phi _\nu (\stackrel{}{r}_2)`$ displays a spatial structure quite similar to $`\mathrm{\Psi }_0(\stackrel{}{r}_1,\stackrel{}{r}_2)`$, the admixture of s- , p- and d- two particle configurations being now 46$`\%`$ and 51$`\%`$ and $`3\%`$ respectively. The coherence length $`\xi `$, that is the mean square distance between the two neutrons forming the Cooper pair, is in this case, $`r_{12}^2^{1/2}`$ = 7.8 fm.
Arguably, the understanding of halo nuclei is the single, most important issue of nuclear structure research still awaiting a satisfactory explanation. We have shown that a substantial advance towards this goal is made by properly characterizing the role the surface of the system plays in renormalizing the bare nucleon-nucleon potential and the single-particle motion of the nucleons. In fact, we find compelling evidence which testifies to the fact that the mechanism which is at the basis of the presence of a low density halo in <sup>11</sup>Li, is the coupling of weakly bound nucleons to long wavelength vibrations of the system leading to pairing instability and thus to a bound system. This result suggests a general strategy for designing nuclei with a large excess of neutrons or protons which may prove valuable in the current exploration of the drip line of the chart of nuclides, and thus in the design of new nucleon species: select those systems which display, under an increase of the excess of one type of nucleons, a marked softening of the long wavelength linear response. It is likely that such systems could venture far inside the region delimited by the drip lines, in keeping with the fact that the associated vibrational modes are expected to be an important source of (induced) pairing interaction, and thus to contribute significantly to the stability of the system.
We wish to thank G. Gori for technical help.
Correspondence and request for materials should be addressed to RAB (e-mail: broglia@mi.infn.it)
Caption to the figures
Fig. 1
(I) Single-particle neutron resonances in <sup>10</sup>Li. In (a) the position of the levels $`s_{1/2}`$ and $`p_{1/2}`$ calculated making use of mean field theory is shown (hatched area and thin horizontal line respectively). The coupling of a single-neutron (upward pointing arrowed line) to a vibration (wavy line) calculated making use of the Feynman diagrams displayed in (b) (schematically depicted also in terms of either solid dots (neutron) or open circles (neutron hole) moving in a single-particle level around or in the <sup>9</sup>Li core (hatched area)), leads to conspicuous shifts in the energy centroid of the $`s_{1/2}`$ and $`p_{1/2}`$ resonances (shown by thick horizontal lines) and eventually to an inversion in their sequence. In (c) we show the calculated partial cross section $`\sigma _l`$ for neutron elastic scattering off <sup>9</sup>Li.
(II) The two-neutron system <sup>11</sup>Li. We show in (a) the mean-field picture of <sup>11</sup>Li, where two neutrons (solid dots) move in time-reversal states around the core <sup>9</sup>Li (hatched area) in the $`s_{1/2}`$ resonance leading to an unbound $`s_{1/2}^2(0)`$ state where the two neutrons are coupled to zero angular momentum. The exchange of vibrations between the two neutrons shown in the upper part of the figure leads to a density dependent interaction which, added to the nucleon-nucleon interaction, correlates the two-neutron system leading to a bound state $`|0^+`$, where the two neutrons move with probability 0.40, 0.58 and 0.02 in the two-particle configurations $`s_{1/2}^2(0)`$, $`p_{1/2}^2(0)`$ and $`d_{5/2}^2(0)`$ respectively.
Fig. 2
Spatial structure of two-neutron Cooper pair. The modulus squared wavefunction $`|\mathrm{\Psi }_0(\stackrel{}{r}_1,\stackrel{}{r}_2)|^2=|\stackrel{}{r}_1,\stackrel{}{r}_2|0^+|^2`$ (cf. Fig. 1, II (b)) describing the motion of the two halo neutrons around the <sup>9</sup>Li core (normalized to unity and multiplied by 16$`\pi ^2r_1^2r_2^2)`$ is displayed as a function of the cartesian coordinates $`x_2=r_2\mathrm{cos}(\theta _{12})`$ and $`y_2=r_2\mathrm{sin}(\theta _{12})`$ of particle 2, for fixed value of the position of particle 1 ($`r_1`$=2.5, 5, 7.5 fm) represented in the right panels by a solid dot, while the core <sup>9</sup>Li is shown as a red circle. The numbers appearing on the $`z`$-axis of the three-dimensional plots displayed on the left side of the figure are in units of fm<sup>-2</sup>.
Fig. 3 Spatial distribution of the pure two-particle configurations $`s_{\mathrm{𝟏}\mathbf{/}\mathrm{𝟐}}^\mathrm{𝟐}`$ (0) and $`p_{\mathrm{𝟏}\mathbf{/}\mathrm{𝟐}}^\mathrm{𝟐}`$(0) as a function of the $`x`$\- and $`y`$-coordinates of particle 2, for a fixed value of the coordinate of particle 1 ($`r_1`$=5 fm). For more details cf. caption to Fig. 2.
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# 1 Introduction
## 1 Introduction
Two-dimensional field theories have been widely explored in the last years and various phenomena such as dynamical mass generation, asymptotic freedom, and quark confinement, relevant in more realistic models, have been tested. It is well known the relevance of localized classical solutions of non-linear relativistic field equations to the corresponding quantum theories . In particular, solitons can be associated with quantum extended-particle states. This picture will be valid for the ‘gauge’ field $`\phi ,`$ of our model. The relevance of the solitons to non-perturbative phenomena comes, in general, from the fact that their interactions are inversely proportional to the coupling constants governing the dynamics of the fundamental fields appearing in the Lagrangian. Therefore, the solitons are weakly coupled in the strong regime of the theory, and that is the basic fact underlying several duality ideas. The reason is that one can describe the theory in the strong coupling regime, by replacing the fundamental Lagrangian by another one where the excitations of its fields correspond to the solitonic states.
We will be interested in the conformal affine Toda system coupled to the matter fields (CATM) . In the authors have studied the classical solitonic solutions, and one of the remarkable features of the model, the equivalence of the Noether and topological currents, and its consequences for the confinement mechanism were outlined. The first thing to be established is if such equivalence is not spoiled by quantum anomalies the currents may present. Fortunately, that issue can be established exactly by using, instead of perturbative approaches, bosonization techniques following the lines of . By bosonizing the spinor field, it is shown that the reduced theory made of the spinor and “gauge” particles, the so-called affine Toda model coupled to the matter (ATM), is equivalent to a theory of a free massless scalar and a sine-Gordon field. In addition, the condition for the equivalence of the Noether and topological currents is simply the condition that the free massless scalar modes must be forbidden in the physical spectrum. Therefore, by performing a quantum reduction where the excitations of the free scalar are eliminated, we obtain a submodel where the equivalence of those currents holds true exactly at the quantum level, proving that there are no anomalies. Consequently, we show that in such reduced theory the confinement of the spinor particles does take place. An important property of the model is that it possesses another type of spinor particle. That is obtained by fermionizing the sine-Gordon field using the well known equivalence between the sine-Gordon and massive Thirring models . In that scenario the solitons of the sine-Gordon model are interpreted as the spinor particles of the Thirring theory. We are then lead to an interesting analogy with what one expects to happen in QCD. The original spinor particles of our model that get confined inside the solitons play the role of the quarks, and the second spinor particle (Thirring), which are the solitons, play the role of the hadrons. The $`U(1)`$ Noether charge is also confined and is analogous to color in QCD (however, see a companion paper ). In this sense, our model constitute an excellent laboratory to test ideas about confinement, the role of solitons in quantum field theory, and dualities interchanging the role of solitons and fundamental particles.
The construction of the conformal affine Toda system coupled to the matter fields (CATM), in the particular case of the affine Lie algebra $`\widehat{sl}(2),`$ using the parlance of the original reference, will be summarised in the following subsection. In section 2 we will obtain, using the dressing procedure formalism, the classical solitonic solutions of the model . In this paper we will provide the explicit forms of the two-soliton solutions and compare with the ones of the sine-Gordon theory. In section 3 we will consider some quantum aspects of the model. We use the perturbative Lagrangian point of view to have an insight into the quantum structure of the theory . Moreover, by bosonizing the spinor field, and subsequently performing a quantum reduction, we outline a confinement mechanism present in a special class of models .
### 1.1 The model
We discuss the example associated with the principal gradation of the untwisted affine Kac-Moody algebra $`\widehat{sl}(2)`$. This belongs to a special class of models introduced in possessing a $`U(1)`$ Noether current depending only on the matter fields. It is then possible, under some circunstances, to choose one solution in each orbit of the conformal group, such that for these solutions, that $`U(1)`$ current is proportional to a topological current depending only on the (gauge) zero grade fields. The zero curvature condition in light-cone coordinates $`x_\pm =t\pm x`$ takes the form (for $`\widehat{sl}(2)`$ affine Lie algebra theory, notations and conventions used here, see, for example, an appendix of Ref )
$`_+A_{}_{}A_++[A_+,A_{}]=0.`$ (1.1)
The connections are of the form
$$A_+=BF^+B^1,A_+=_{}BB^1+F^{},$$
(1.2)
where the mapping $`B`$ is parametrized as
$$B=be^{\nu C}e^{\eta Q_𝐬}=e^{\phi H^0}e^{\widehat{\nu }C}e^{\eta Q_𝐬},\text{with}b=e^{\phi \stackrel{~}{H}^0},$$
(1.3)
and so $`\widehat{\nu }=\nu \frac{1}{2}\phi .`$ Here $`Q_𝐬`$ ($`𝐬=(1,1)`$) is the grading operator for the principal gradation .
The special class of models in which the $`U(1)`$ Noether current is proportional to a topological current, occur for those models where the first constant terms of the nonvanishing grade potentials $`F^\pm `$ are equal to $`\pm N_𝐬,(N_𝐬=𝐬_1+𝐬_2=2)`$ respectively. So, the potentials $`F^\pm `$ are of the form
$$F^+E_{N_𝐬}+F_1^+,F^{}E_{N_𝐬}+F_1^{},$$
(1.4)
where
$`F_1^+=2\sqrt{im}(\psi _RE_+^0+\stackrel{~}{\psi }_RE_{}^1),F_1^{}=2\sqrt{im}(\psi _LE_+^1\stackrel{~}{\psi }_LE_{}^0),`$ (1.5)
and $`E_{\pm N_𝐬}=mH^{\pm 1}`$ ($`m`$ =constant).
Introducing the Dirac fields $`\psi ^T=(\psi _R,\psi _L)`$ and substituting the explicit form of the connections $`b`$ and $`F_1^\pm `$into (1.1), we get the following equations of motion
$`^2\phi =4m_\psi \overline{\psi }\gamma _5e^{\eta +2\phi \gamma _5}\psi ,`$
$`^2\stackrel{~}{\nu }=2m_\psi \overline{\psi }(1\gamma _5)e^{\eta +2\phi \gamma _5}\psi {\displaystyle \frac{1}{2}}m_\psi ^2e^{2\eta },`$
$`^2\eta =\mathrm{\hspace{0.17em}0},i\gamma ^\mu _\mu \psi =m_\psi e^{\eta +2\phi \gamma _5}\psi ,m_\psi 4m.`$ (1.6)
The corresponding Lagrangian has the form
$`{\displaystyle \frac{1}{k}}={\displaystyle \frac{1}{4}}_\mu \phi ^\mu \phi +{\displaystyle \frac{1}{2}}_\mu \nu ^\mu \eta {\displaystyle \frac{1}{8}}m_\psi ^2e^{2\eta }+i\overline{\psi }\gamma ^\mu _\mu \psi `$
$`m_\psi \overline{\psi }e^{\eta +2\phi \gamma _5}\psi .`$ (1.7)
It is real (for $`\eta `$ $`=`$real constant) if $`\stackrel{~}{\psi }`$ is proportional to the complex conjugate of $`\psi ,`$ and if $`\phi `$ is pure imaginary. This will be true for the particular solutions of (1.6) such as: the 1-soliton (1-antisoliton), soliton-soliton (antisoliton-antisoliton).
The equations (1.6) are invariant under the conformal transformation
$`x_+f(x_+),x_{}f(x_{}),`$ (1.8)
with $`f`$ and $`g`$ being analytic functions; and with the fields transforming conveniently.
Therefore we can construct two chiral currents
$`𝒥(x_+)=m_+\phi +m_+\eta 2im\stackrel{~}{\psi }_R\psi _R\text{and}\overline{𝒥}(x_{})=m_{}\phi +m_{}\eta +2im\stackrel{~}{\psi }_L\psi _L,`$ (1.9)
satisfying
$`_{}𝒥=0,_+\overline{𝒥}=0.`$ (1.10)
Defining the current
$`\stackrel{~}{J}_+=2im\stackrel{~}{\psi }_R\psi _R,\stackrel{~}{J}_{}=2im\stackrel{~}{\psi }_L\psi _L,`$ (1.11)
we may check that it is conserved
$`_\mu \stackrel{~}{J}^\mu =0.`$ (1.12)
It is the $`U(1)`$ Noether current. In addition, this model possesses the current
$`\stackrel{~}{j}_+=m_+\phi ,\stackrel{~}{j}_{}=m_{}\phi ,`$ (1.13)
which is a topological current, i.e., it is conserved independently of the equations of motion
$`_\mu \stackrel{~}{j}^\mu =0.`$ (1.14)
Under the conformal transformations the chiral currents transform as
$`𝒥(x_+)(\mathrm{ln}f^{}\left(x_+\right))^1(𝒥(x_+)m(\mathrm{ln}f^{}\left(x_+\right))^{}),`$
$`\overline{𝒥}(x_{})(\mathrm{ln}g^{}\left(x_{}\right))^1(\overline{𝒥}(x_{})m(\mathrm{ln}g^{}\left(x_{}\right))^{}).`$ (1.15)
One concludes that, given a solution of the model, one can always map it, under a conformal transformation, into a solution where
$`𝒥(x_+)=\overline{𝒥}(x_{})=0.`$ (1.16)
Such a procedure amounts to a gauge fixing of the conformal symmetry. We are choosing one solution in each orbit of the conformal group. The degree of fredom eliminated corresponds to the field $`\eta .`$ So, in the model defined by (1.7) one observes that the Noether current $`\stackrel{~}{J}^\mu `$ and the topological current $`\stackrel{~}{j}^\mu `$ are equal $`\stackrel{~}{J}_\mu =\stackrel{~}{j}_\mu ,`$ as a result of the field equations.
From now on, unless otherwise stated, we consider the Lagrangian (1.7) with its conformal symmetry gauge fixed by choosing $`\eta =\eta _0=`$const. Then we define the off-critical model
$`_{ATM}={\displaystyle \frac{1}{4}}_\mu \phi ^\mu \phi +i\overline{\psi }\gamma ^\mu _\mu \psi m_\psi e^{\eta _0}\overline{\psi }e^{2\phi \gamma _5}\psi {\displaystyle \frac{1}{8}}m_\psi ^2e^{2\eta _0},`$ (1.17)
which we call affine Toda model coupled to the matter (ATM).
Now, let us study the symmetries of this Lagrangian. The conservation law (1.12) corresponds to the global U(1) symmetry
$`\psi e^{i\theta }\psi ,\phi \phi ,\stackrel{~}{\nu }\stackrel{~}{\nu },\text{with the Noether current}J^\mu =\overline{\psi }\gamma ^\mu \psi .`$ (1.18)
The fields $`\psi `$ and $`\stackrel{~}{\psi }`$ have charges $`1`$ and $`1`$, respectively; and $`\phi `$ has charge zero.
The global chiral U(1) symmetry
$`\psi e^{i\gamma _5\alpha }\psi ,\phi \phi i\alpha ,\stackrel{~}{\nu }\stackrel{~}{\nu },`$ (1.19)
also leaves the action invariant. The relevant conservation law reads
$`_\mu [\overline{\psi }\gamma ^\mu \gamma _5\psi +{\displaystyle \frac{1}{2}}^\mu \phi ]=0.`$ (1.20)
Let us next see the relationship between the Noether and topological currents. The topological current and charge are
$$j^\mu =\frac{1}{2\pi i}ϵ^{\mu \nu }_\nu \phi ,Q_{topol}𝑑xj^0.$$
Indeed, the Lagrangian (1.7) has infinitely many degenerate vacua due to the invariance under $`\phi \phi +in\pi .`$
Combining the chiral current and the vector conservation laws, and applying the arguments presented above we can set $`𝒥`$ $`=\overline{𝒥}=0.`$ This gives, altogether,
$`{\displaystyle \frac{1}{2\pi i}}ϵ^{\mu \nu }_\nu \phi ={\displaystyle \frac{1}{\pi }}\overline{\psi }\gamma ^\mu \psi ,`$ (1.21)
so that the topological and Noether currents are proportional.
The equation (1.21) means that the Noether density is non zero only where $`\phi 0,`$ that is inside the solitons. Thus the $`\psi `$ field is confined inside the solitons. It will hold at the quantum level after a suitable redefinition, we shall come back to this point below.
## 2 Dressing procedure and classical solitons
At this point the field $`\psi `$ seems to be a c-number field. It is well known, that in two dimensions the statistics of the fields depends upon the coupling constant; so, we will postpone the study of their statistics to section 3. Field equations with the replacement of the Fermi fields by a c-number Dirac wave function have often been introduced, ab initio, in semiclassical calculations in the literature, particularly in constructing models for hadrons .
Therefore, we will examine the classical soliton type solutions to get insight into the quantum expectrum of the model in much the same way as in the remarkable sine-Gordon model. We shall argue that, at the classical level, the solutions for the $`\phi `$ and $`\psi `$ fields share some features of the sine-Gordon and the massive Thirring theories, respectively.
The dressing transformation procedure provides a powerful way of solving a wide class of nonlinear equations presenting soliton solutions . For example, an application of the method to the vector NLS hierarchies can be found in . We use this method to obtain the 1-soliton and 2-soliton solutions of the model. We want to study the behavior of the solitonic solutions taking care of the reality conditions of the fields ($`\stackrel{~}{\psi }_R=\psi _R^{},`$ $`\stackrel{~}{\psi }_L=\psi _L^{}`$ and $`\phi =`$ pure imaginary and $`\eta =`$ constant) which make the Lagrangian real.
Let us define $`\epsilon _\pm =mH^\pm `$ and let $`\widehat{\lambda }_0>`$ and $`\widehat{\lambda }_1>`$ be the highest weight states of two fundamental representations of the affine Kac-Moody algebra $`\widehat{sl}(2),`$ respectively the scalar and spinor ones. Then the solutions on the orbit of the vaccum are the same as in the equations (10.31) of
$`\psi _R=\sqrt{{\displaystyle \frac{m}{i}}}{\displaystyle \frac{\tau _{01}}{\tau _0}},\stackrel{~}{\psi }_R=\sqrt{{\displaystyle \frac{m}{i}}}{\displaystyle \frac{\tau _{12}}{\tau _1}},`$
$`\psi _L=\sqrt{{\displaystyle \frac{m}{i}}}{\displaystyle \frac{\tau _{11}}{\tau _1}},\stackrel{~}{\psi }_L=\sqrt{{\displaystyle \frac{m}{i}}}{\displaystyle \frac{\tau _{02}}{\tau _0}},`$ (2.1)
$`e^\phi ={\displaystyle \frac{\tau _1}{\tau _0}},e^{(\stackrel{~}{\nu }\nu _o)}=\tau _0,`$
where the so-called tau functions are given by
$`\tau _{0,1}=<\widehat{\lambda }_{0,1}`$ $`G`$ $`\widehat{\lambda }_{0,1}>,`$ (2.2)
$`\tau _{01}=<\widehat{\lambda }_0E_{}^1G\widehat{\lambda }_0>`$ , $`\tau _{02}=<\widehat{\lambda }_0GE_+^1\widehat{\lambda }_0>,`$ (2.3)
$`\tau _{12}=<\widehat{\lambda }_1E_+^0G\widehat{\lambda }_1>`$ , $`\tau _{11}=<\widehat{\lambda }_1GE_{}^0\widehat{\lambda }_1>,`$ (2.4)
with G=e$`{}_{}{}^{x_+\epsilon _+}e_{}^{x_{}\epsilon _{}}\rho e^{x_{}\epsilon _{}}e^{x_+\epsilon _+}`$; $`\rho `$ being a constant group element of the $`\widehat{SL}(2)`$ Kac-Moody group.
### 2.1 The 1-soliton solutions
The soliton solutions of the system are constructed as follows. Let us choose
$`\rho =e^{\sqrt{i}a_+V_+(z)}e^{\sqrt{i}a_{}V_\mathrm{\_}(z)},`$ (2.5)
with \[$`\epsilon _\pm `$,V<sub>±</sub>\]=$`\varpi _\pm ^\pm `$V<sub>±</sub>. The particular factor $`\sqrt{i}`$ is chosen such that the reality condition will be obeyed with $`a_+=a_{}^{}`$ .
Computing the matrix elements one gets the explicit form of the solutions generated by $`\rho `$
$`\phi =2\mathrm{arctan}\left(\mathrm{exp}\left(2m_\psi \left(xx_0vt\right)/\sqrt{1v^2}\right)\right),`$
$`\psi =e^{i\theta }\sqrt{m_\psi }e^{m_\psi \left(xx_0vt\right)/\sqrt{1v^2}}`$
$`\left(\begin{array}{c}\left(\frac{1v}{1+v}\right)^{1/4}\frac{1}{1+ie^{2m_\psi \left(xx_0vt\right)/\sqrt{1v^2}}}\\ \left(\frac{1+v}{1v}\right)^{1/4}\frac{1}{1ie^{2m_\psi \left(xx_0vt\right)/\sqrt{1v^2}}}\end{array}\right),`$ (2.8)
$`\nu ={\displaystyle \frac{1}{2}}\mathrm{log}\left(1+\mathrm{exp}\left(4m_\psi \left(xx_0vt\right)/\sqrt{1v^2}\right)\right)`$
$`+{\displaystyle \frac{1}{8}}m_\psi ^2x_+x_{},`$
$`\eta =0,`$ (2.9)
and $`\stackrel{~}{\psi }`$ is the complex conjugate of $`\psi .`$ The $`\psi `$ field solution is similar to the 1-soliton solution of the massive Thirring model . Notice that the $`\phi `$ solution is of the sine-Gordon type soliton/antisoliton .
It can be directly verified that these solutions satisfy the relation (1.21). It is clear from the explicit expressions for the fermion field and the scalar, that $`\psi `$ vanishes exponentially when $`xx_0\pm \mathrm{},`$ so that the Dirac field is confined inside the soliton.
### 2.2 The 2-soliton solutions
Now let us move to the study of the 2-soliton solutions of the model. As usual in this case consider the following steps:
Let us choose
$`G=e^{x_+\epsilon _+}e^{x_{}\epsilon _{}}\rho e^{x_{}\epsilon _{}}e^{x_{}\epsilon _+},`$ (2.10)
and as the constant group element
$`\rho =e^{b_1V_+(\rho _1)}e^{a_1V_{}(\nu _1)}e^{b_2V_+(\rho _2)}e^{a_2V_{}(\nu _2)}.`$ (2.11)
Defining $`A_i=a_ie^{\mathrm{\Gamma }(\nu _i)},`$ $`B_i=b_ie^{\mathrm{\Gamma }(\rho _i)}`$ and $`\mathrm{\Gamma }(z)=2m(zx_+\frac{1}{z}x_{}),`$ and using the properties of level 1 vertex operators to determine the matrix elements in the tau functions , one can obtain the soliton-soliton and the soliton-antisoliton solutions.
By direct substitution, one can verify that these solutions, indeed, satisfy the equation (1.21) without any further restriction. This allows us to conclude that, also, for 2-soliton type solutions the Noether density is non zero only where $`\phi 0,`$ that is inside the solitons.
Let us choose the parameters conveniently in order to express in the form of a soliton-soliton type solution.
Choosing the following relationships between the parameters $`\nu _{1+}\nu _2+\nu _1^1+\nu _2^1=0,`$ $`\nu _1\nu _2=1,`$ and defining $`u=(\nu _1^21)/(\nu _1^2+1),`$ one can write
$`\phi _{SS}/2i=\mathrm{arctan}\left(8u^3\sqrt{1u^2}{\displaystyle \frac{\mathrm{sinh}(8m\gamma x\mathrm{log}\overline{a})}{\mathrm{cosh}(8m\gamma ut)2u\sqrt{1u^2}}}\right),`$ (2.12)
where $`\gamma =1/\sqrt{1u^2},`$ $`u<1.`$
Its asymptotic behaviour in time can be written as
$`{\displaystyle \frac{\phi _{SS}}{2i}}(t\mathrm{})\mathrm{tan}^1(e^{8m\gamma (xu(t+\frac{\delta _{SS}}{2}))})\mathrm{tan}^1(e^{8m\gamma (x+u(t+\frac{\delta _{SS}}{2}))}),`$ (2.13)
$`{\displaystyle \frac{\phi _{SS}}{2i}}(t+\mathrm{})\mathrm{tan}^1(e^{8m\gamma (x+u(t\frac{\delta _{SS}}{2}))})\mathrm{tan}^1(e^{8m\gamma (xu(t\frac{\delta _{SS}}{2}))}),`$ (2.14)
$`\delta _{SS}={\displaystyle \frac{1}{4mu\gamma }}\mathrm{log}(8u^2\sqrt{1u^2}).`$
Therefore, this solution correponds to two solitons approaching one another with relative velocity $`u`$ in the distant past. Notice that the time delay $`\delta _{SS},`$ of course, differs from the corresponding soliton-soliton interaction in the sine-Gordon model . We will come back to this point below, when we compare our soliton solutions with the corresponding solitons of the sine-Gordon theory.
In order to obtain soliton-antisoliton type solution we convert the real parameters into an imaginary ones by making $`\nu _1i\nu _1,`$ $`\nu _2i\nu _2`$ , $`mim.`$
Setting the relationships $`\nu _1+\nu _2\nu _1^1+\nu _2^1=0,`$ $`\nu _1\nu _2=1,`$ we may obtain
$`\phi _{S\overline{S}}/2i=\mathrm{arctan}\left({\displaystyle \frac{16}{u^2}}{\displaystyle \frac{\mathrm{sinh}(8m\gamma ut+\mathrm{log}\overline{a})}{\mathrm{cosh}(8m\gamma x\mathrm{log}(4/\gamma u^2))2/\gamma u^2}}\right),`$ (2.15)
with $`\gamma =1/\sqrt{u^21},`$ $`u>1`$.
Even though the parameter complexifications do not modify the purely imaginary character of $`\phi ,`$ the reality condition $`\stackrel{~}{\psi }=\psi ^{}`$ of the $`\psi `$ field is lost, converting the Lagrangian into a complex one. Moreover, with this choice of parameters the velocity $`u`$ becomes greater than the velocity of light, and the mass of the soliton becomes imaginary. Nevertheless, from the mathematical point of view, still we can write the asymptotic behaviour, extracting a ‘soliton’ and an ‘antisoliton’ approaching one another with relative velocity $`u`$ in the distant past and ‘time’ delay $`\delta =\frac{1}{2m\gamma }\mathrm{log}(u/4).`$
Therefore, we have not obtained a physical bound soliton-antisoliton pair (breather) solution of the model. This fact is quite remarkable if one intends to study the quantum spectrum of the model. A discussion of this point is presented in a companion work .
In order to have a better insight into the behaviour of the two-solitons of our model, let us consider some known facts about the sine-Gordon field $`\varphi `$ solitons. The $`\varphi `$ field is related to its corresponding tau functions as $`i\varphi =\mathrm{ln}(\frac{\tau _+}{\tau _{}}).`$ The tau functions of the N-soliton solutions of the sine-Gordon equation are given by
$$\tau _\pm ^{(N)}(z_+,z_{})=det(1\pm V),$$
(2.16)
with $`V`$ a $`N`$ x $`N`$ matrix with elements
$`V_{ij}=2{\displaystyle \frac{\sqrt{\mu _i\mu _j}}{\mu _i+\mu _j}}\sqrt{X_iX_j},\text{where}X_i=a_ie^{2(\mu _iz_++\frac{1}{\mu _i}z_{})},z_\pm =2mx_\pm .`$ (2.17)
For the 2-soliton solution the determinant takes a simple form. With a convenient choice of the parameters, one can write for the soliton-soliton and soliton-antisoliton
$$\varphi _{SS}/2=\mathrm{tan}^1\left(u\frac{\mathrm{sinh}(8m\gamma x)}{\mathrm{cosh}(8m\gamma ut)}\right),\varphi _{S\overline{S}}/2=\mathrm{tan}^1\left(\frac{1}{u}\frac{\mathrm{sinh}(8m\gamma ut)}{\mathrm{cosh}(8m\gamma x)}\right),$$
(2.18)
with $`\gamma =1/\sqrt{1u^2},u<1.`$
If one makes the change $`uiv`$ in $`\varphi _{S\overline{S}}`$ one gets
$$\varphi _v/2=\mathrm{arctan}\left(\frac{1}{v}\frac{\mathrm{sinh}(16m\gamma vt)}{\mathrm{cosh}(16m\gamma x)}\right),$$
(2.19)
which still gives a solution of the sine-Gordon equation. It is the doublet or breather solutions, which represents a soliton and antisoliton oscillating about a common center.
Then we can conclude that our solution (2.12) resembles the sine-Gordon soliton-soliton, at least asymptotically; the differences manifest in their time delays, which can be interpreted as the effect of the spinor field $`\psi `$ on the Toda field $`\phi `$, since our model couples those fields. Also, as in the sine-Gordon case, the equation (2.15) can be put asymptotically in the form of the mathematical ‘soliton’ and ‘antisoliton’, even though, each of them with imaginary mass. Notice that in the soliton-soliton sector the time delay $`\delta _{SS}`$ can eventually become negative, unlike the corresponding $`\delta `$ of the sine-Gordon model, which is always positive, indicating in the latter case a repulsive force .
To conclude this section, let us point out that the model possesses at the classical level some solitonic solutions which resemble the sine-Gordon and the massive Thirring models, for the scalar and the Dirac fields, respectively. We have verified, up to the 2-soliton solution, the equivalence between the Noether and topological currents. The classical breather solution was not found, and it is expected that it will not appear at the quantum level.
## 3 Quantum aspects of the model
The aim of the present section is to study some quantum aspects of the model defined by the Lagrangian (1.7).
### 3.1 Perturbative Lagrangian viewpoint and the CATM $``$ ATM reduction
In this subsection we will be interested in verifying how the classical reduction of the CATM to the ATM by setting $`\eta =\eta _0=`$constant is recovered at the quantum level. This will give us some information about the vacuum of the theory, in particular about the trivial vacuum (non-topological) of the ATM model. We use the perturbative Lagrangian point of view to have an insight into the quantum structure of the theory. Let us begin expanding the action around a simple solution of the classical equations of motion (1.6)
$`\phi =\eta =\psi =\stackrel{~}{\psi }=0,\nu ={\displaystyle \frac{1}{8}}m_\psi ^2x_+x_{}.`$ (3.1)
Denote collectively all the fields as
$`\mathrm{\Phi }=\mathrm{\Phi }_0+\mathrm{\Phi }^{},`$ (3.2)
where $`\mathrm{\Phi }_0`$ represent the relevant classical solutions (3.1). Let us expand every field around (3.1). Then relabelling the $`\mathrm{\Phi }^{}`$ with the same symbol as the original ones, the action corresponding to (1.7) turns into
$`S`$ $`=`$ $`{\displaystyle }d^2x\{{\displaystyle \frac{1}{4}}\phi \mathrm{}\phi {\displaystyle \frac{1}{2}}\nu \mathrm{}\eta {\displaystyle \frac{1}{8}}m_\psi ^2{\displaystyle \frac{1}{4}}m_\psi ^2+i\overline{\psi }\gamma _\mu ^\mu \psi m_\psi \overline{\psi }\psi 2im_\psi \overline{\psi }\gamma _5\psi \phi `$ (3.3)
$`m_\psi \eta \overline{\psi }\psi +\mathrm{}\},`$
where the ellipsis denotes quartic or higher interaction terms in which the $`\nu `$ field never appears, and $`\mathrm{}=_t^2_x^2`$. Next let us perform a Wick rotation and compute the Fourier transformed propagators. They become
$`<\phi ,\phi >`$ $`=`$ $`{\displaystyle \frac{1}{\frac{1}{2}k^2}},<\eta ,\nu >={\displaystyle \frac{1}{\frac{1}{2}k^2}},`$ (3.4)
$`<\nu ,\nu >`$ $`=`$ $`{\displaystyle \frac{\frac{1}{2}m_\psi ^2}{\frac{1}{4}k^4}},<\psi ,\overline{\psi }>={\displaystyle \frac{\gamma _\mu p^\mu +im_\psi }{p^2+m_\psi ^2}},`$ (3.5)
where $`k^2=k_0^2+k_x^2`$. All the other propagators vanish. Let us observe that the pole at $`p^2=m_\psi ^2`$ in the spinor field propagator, corresponds to a massive particle of the ATM model. Let us inspect now a scattering process whose external legs consist only of the $`\psi (\overline{\psi })`$ and $`\phi `$ particles. From the structure of the propagators, and of the interaction terms it is easy to see that no connected graph relevant for this process will have propagation of other modes than $`\psi (\overline{\psi })`$ and $`\phi `$ themselves. That is to say, the other modes decouple. This fact allows us to set the fields $`\nu `$ and $`\eta `$ to zero in the action. Therefore, the action becomes exactly the action of the ATM model. In other words, for scattering processes involving massless $`\phi `$ modes and massive $`\psi `$ modes, the theory may be described by the perturbative ATM model.
Then we have derived an effective Lagrangian with $`\nu =0`$, $`\eta =\eta _0=`$const., this is exactly the off-critical affine Toda model coupled to the matter (ATM) defined above, Eq. (1.17). Then we can regard the ATM model as a spontaneously broken and reduced version of the CATM.
The symplectic structure of the reduced model (1.17) has recently been studied . It was performed in the context of Faddeev-Jackiw and (constrained) symplectic methods (Wotzasek, Montani and Barcelos-Neto); by imposing the equivalence between the Noether and topological currents, Eq. (1.21), as a constraint, the authors have been able to obtain either, the sine-Gordon model or the massive Thirring model, through a Hamiltonian reduction and gauge fixing the symmetries of the model in two different ways.
### 3.2 Bosonization procedure and the ATM model
In two-dimensional quantum field theories it is possible to transform Fermi fields into Bose fields, and vice versa (for a complete review of the most important references in the field see, e.g., the second Ref. of ). The existence of such a transformation, called bosonization, provided a powerful tool to obtain nonperturbative information in two-dimensional field theories.
Let us consider the ATM Lagrangian (1.17) in a slightly modified form. This modification is undertaken by imposing a convenient reality conditions (e.g. $`\stackrel{~}{\psi }=\psi ^{}`$, $`\phi `$ pure imaginary) on the fields and dropping a overall minus sign. See Refs. for discussions concerning the positive-definite caracter of the kinetic terms of the real Lagrangian submodel, and for the implications of these conditions on the solitonic solutions of the original model. Then let us write the real action in a form which is more convenient for the bosonization of the fermion bilinears
$`{\displaystyle \frac{1}{k}}S={\displaystyle }d^2x\{{\displaystyle \frac{1}{4}}_\mu \phi ^\mu \phi +i\overline{\psi }\gamma ^\mu _\mu \psi `$
$`m_\psi e^{\eta _0}[\overline{\psi }{\displaystyle \frac{(1+\gamma _5)}{2}}\psi e^{2i\phi }+\overline{\psi }{\displaystyle \frac{(1\gamma _5)}{2}}\psi e^{2i\phi }]{\displaystyle \frac{1}{8}}m_\psi ^2e^{2\eta _0}\},`$ (3.6)
where we have considered the pure imaginary character of the scalar field by making the change $`\phi i\phi `$ in (1.17).
The model (3.6) was considered in Refs. . Some of the points that folllow were discussed in those papers. The model defined by (3.6) possesses a chiral symmetry $`\psi e^{i\beta \gamma _5}\psi ,`$ $`\phi \phi \beta `$. Apparently, if unbroken, the symmetry would prevent the fermions from having a mass. As we will see, however, the symmetry is not broken, but the fermion has a mass.
Following one can use the boson representation of fermions as
$$i\overline{\psi }\gamma ^\mu _\mu \psi =\frac{1}{2}(_\mu c)^2,$$
(3.7)
$$\overline{\psi }(1\pm \gamma _5)\psi =\mu \mathrm{exp}(\pm i\sqrt{4\pi }c),$$
(3.8)
$$\overline{\psi }\gamma ^\mu \psi =\frac{1}{\sqrt{\pi }}ϵ^{\mu \nu }_\nu c.$$
(3.9)
Then the action now becomes
$`S_B`$ $`=`$ $`{\displaystyle }d^2x\{{\displaystyle \frac{1}{2}}_\mu \phi ^\mu \phi +{\displaystyle \frac{1}{2}}(_\mu \varphi )^2`$ (3.10)
$`{\displaystyle \frac{1}{2}}\mu m_\psi [\mathrm{exp}i[\sqrt{8/k}\phi +\sqrt{4\pi }\varphi ]+\mathrm{exp}i[\sqrt{8/k}\phi +\sqrt{4\pi }\varphi ]]\}.`$
Introducing new fields
$$\stackrel{~}{c}\frac{\phi /a+\sqrt{4\pi }\varphi }{\sqrt{4\pi +1/a^2}},\stackrel{~}{\sigma }\frac{\varphi /a+\phi \sqrt{4\pi }}{\sqrt{4\pi +1/a^2}},$$
(3.11)
the action takes the form
$$S_{sG+\stackrel{~}{\sigma }}=d^2x[\frac{1}{2}(_\mu \stackrel{~}{c})^2+\frac{1}{2}(_\mu \stackrel{~}{\sigma })^2+\mu m_\psi \mathrm{cos}(\sqrt{4\pi +1/a^2}\stackrel{~}{c}),$$
(3.12)
where $`a^2k/8`$.
Therefore we obtain a sine-Gordon model for the field $`\stackrel{~}{c}`$ and a free, massless scalar $`\stackrel{~}{\sigma }`$ field.
Why do the fermion and antifermion acquire masses despite chiral symmetry? To answer this issue, we should ask what form the original chiral current takes in terms of $`\stackrel{~}{c}`$ and $`\stackrel{~}{\sigma }`$. The chiral current defined in section 1 is
$$A_\mu i\overline{\psi }\gamma _\mu \gamma _5\psi +\frac{i}{2}_\mu \phi .$$
(3.13)
In terms of the new variables, one finds
$$A_\mu =\frac{i}{2}_\mu \stackrel{~}{\sigma }.$$
(3.14)
Thus the chiral current involves only $`\stackrel{~}{\sigma }`$ and not $`\stackrel{~}{c}.`$ It is conserved at the quantum level. This means that the field $`\stackrel{~}{c}`$, and therefore also the physical fermion and antifermion associated with this field, are neutral under chirality. Thus, even though the elementary fermion field $`\psi `$ has nonzero chirality, the physical fermion particle has zero chirality.
The presence of the physical fermions (zero chirality particles) can be clarified introducing a new fermion field that has the quantum numbers of the physical particles. We simply introduce a new fermion field $`\widehat{\psi }`$ with
$`{\displaystyle \frac{1}{2\pi \alpha }}ϵ^{\mu \nu }_\nu \stackrel{~}{c}=\overline{\widehat{\psi }}\gamma ^\mu \widehat{\psi }.`$ (3.15)
According to the standard rules which identify the sine-Gordon theory to the charge-zero sector of the massive Thirring model , the action(3.12) can now be written as
$`S={\displaystyle d^2x[i\overline{\widehat{\psi }}\gamma ^\mu _\mu \widehat{\psi }m_F\overline{\widehat{\psi }}\widehat{\psi }\frac{1}{2}g(\overline{\widehat{\psi }}\gamma ^\mu \widehat{\psi })^2+\frac{1}{2}(_\mu \stackrel{~}{\sigma })^2]}.`$ (3.16)
where
$`{\displaystyle \frac{4\pi +1/a^2}{4\pi }}={\displaystyle \frac{1}{1+g/\pi }}.`$ (3.17)
Then, we observe that our theory consists of a massive fermion with self-interaction, and a free, massless scalar.
We now discuss the quantum version of the reduction (1.16). The chiral currents become
$`𝒥=i/2_+\stackrel{~}{\sigma },\overline{𝒥}=i/2_{}\stackrel{~}{\sigma }.`$ (3.18)
Therefore, one can argue that semiclassically, the constraints (1.16) are equivalent to $`\stackrel{~}{\sigma }=\mathrm{const}.`$, and consequently the degree of freedom eliminated by them corresponds to the $`\stackrel{~}{\sigma }`$ field.
The reduction at the quantum level can be realized as follows. Since the scalar field $`\stackrel{~}{\sigma }`$ is decoupled from all others fields in the theory (3.12), we can denote the space of states as $`=_{\stackrel{~}{\sigma }}_0`$, where $`_{\stackrel{~}{\sigma }}`$ is the Fock space of the free massless scalar $`\stackrel{~}{\sigma }`$, and $`_0`$ carries the states of the rest of the theory. We shall denote $`\stackrel{~}{\sigma }=\stackrel{~}{\sigma }^{(+)}+\stackrel{~}{\sigma }^{()}`$, where $`\stackrel{~}{\sigma }^{(+)}`$ ($`\stackrel{~}{\sigma }^{()}`$) corresponds to the part containing annihilation (creation) operators in its expansion on plane waves.
The reduction corresponds to the restriction of the theory to those states satisfying
$$_\pm \stackrel{~}{\sigma }^{(+)}\mathrm{\Psi }=0.$$
(3.19)
By taking the complex conjugate one gets
$$\mathrm{\Psi }_\pm \stackrel{~}{\sigma }^{()}=0,$$
(3.20)
and consequently the expectation value of $`_\pm \stackrel{~}{\sigma }`$ vanishes on such states. Indeed, if $`\mathrm{\Psi }`$ and $`\mathrm{\Psi }^{}`$ are two states satisfying (3.19), then
$$\mathrm{\Psi }^{}_\pm \stackrel{~}{\sigma }\mathrm{\Psi }=\mathrm{\Psi }^{}[_\pm \stackrel{~}{\sigma }^{()}+_\pm \stackrel{~}{\sigma }^{(+)}]\mathrm{\Psi }=0.$$
(3.21)
That provides the correspondence with the classical constraints (1.16). Therefore, the Hilbert space of the reduced theory is $`_c=\mathrm{\Psi }_0`$.
For the theory described by $`_c`$, the equivalence between Noether and topological currents, given by (1.21), holds true at the quantum level, since as we have shown before, (1.21) is equivalent to the vanishing of the currents $`𝒥`$ and $`\overline{𝒥}`$. Notice that such quantum equivalence is exact, since we have not used perturbative or semiclassical methods.
One of the consequences of that quantum equivalence is that in the states of $`_c`$, like the one-soliton (2.9), where the space derivative of $`\phi `$ is localized, one has that the spinor $`\psi `$ is confined, since from (1.21), $`_x\phi \psi ^{}\psi `$. Therefore, we have shown that the confinement of $`\psi `$ does take place in the quantum theory.
The properties of the theory (3.6) at the quantum level are quite remarkable. In the weak coupling regime, i.e. small $`k`$, the excitations around the vacuum correspond to the spinor $`\psi `$ and the “gauge”particle $`\phi `$. The $`U(1)`$ symmetry (1.19) is not broken and the charged states correspond to the $`\psi `$ particles. Consider now those states satisfying (3.19) and look for the fluctuations around the state corresponding, for instance, to the one-soliton solution (2.9). The $`\psi `$ particles disappear from the spectrum since they are confined inside the soliton. The $`\psi `$ particles can live outside the soliton only in bound states with vanishing $`U(1)`$ charge. The theory, however, presents another spinor particle corresponding to the excitations of the Thirring field $`\chi `$, which have zero $`U(1)`$ charge. However, according to Coleman’s interpretation of the sine-Gordon/Thirring equivalence, such excitations correspond to the solitons themselves. Therefore, we can make an analogy with what is expected to happen in QCD. The $`\psi `$ and $`\chi `$ particles are like the quarks and hadrons respectively. The $`U(1)`$ charge is analogous to color in QCD, since it is also confined. For the rigorous treatment of the dynamical flavor and color symmetries of the ATM model, as well as their relevant realizations, and the role played in a topological confinement mechanism, see a companion paper . Then, our model can be considered as a one dimensional bag model for QCD .
It could be interesting to compute the quantum corrections to the mass of the solitons associated to the $`\phi `$ field of the ATM model, for example, in the lines of Refs. . In our case, the quantum fluctuations of the spinor fields, of course, must be taken into account.
## Appendix A Appendix: Notations and Conventions
We use the following conventions in two dimensions. The metric tensor is $`g_{\mu \nu }=\text{diag}(1,1)`$ and the antisymmetric tensor $`ϵ_{\mu \nu }`$ is defined so that $`ϵ_{01}=ϵ^{01}=1`$. $`_\pm `$ are derivatives w.r.t. to the light cone variables $`x_\pm =t\pm x`$. The gamma matrices are in the following representation:
$`\gamma _0=i\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\gamma _1=i\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\gamma _5=\gamma _0\gamma _1=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$ (A.7)
satisfying anticommutation relations
$`\{\gamma _\mu ,\gamma _\nu \}=2g_{\mu \nu }\mathrm{𝟏},`$ (A.8)
so the spinors $`\psi `$ and $`\overline{\psi }`$ are of the form
$`\psi =\left(\begin{array}{c}\psi _R\\ \psi _L\end{array}\right),\stackrel{~}{\psi }=\left(\begin{array}{c}\stackrel{~}{\psi }_R\\ \stackrel{~}{\psi }_L\end{array}\right),\overline{\psi }=\left(\begin{array}{cc}\stackrel{~}{\psi }_R\stackrel{~}{\psi }_L& \end{array}\right)\gamma _0.`$ (A.14)
Acknowledgements
I would like to thank Professor L.A. Ferreira for his colaboration on parts of this work, and Professors G.M. Sotkov and A.H. Zimerman for valuable discussions. R. Bentín and C. Tello are also acknowledged for interesting conversations. I am grateful to Professor H.G. Valqui for introducing me to the study of non-linear phenomena and solitons. Thanks are also due to the organizers of the VII Hadron Physics 2000 for a very enjoyable Workshop. This work has been supported by FAPESP.
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# Generalized Master Function Approach to Quasi-Exactly Solvable Models
## 1 INTRODUCTION
During the last decade a remarkable new class of quasi-exactly solvable spectral problems was introduced . These occupy an intermediate position between exactly solvable and unsolvable models in the sense that exact solution in an algebraized form exists only for a part of the spectrum.
The usual approach to the analysis of quasi-exactly solvable systems is an algebraic one in which the operator is expressed as a non-linear combination of generators of a Lie algebra. Another recent developement is the work of Bender-Dunne where they have shown that the eigen-functions of a quasi-exactly solvable schrodinger equation is the generating function for a set of orthogonal polynomials $`P_m(E)`$ in energy variable. It was further shown that, these polynomials satisfy the three-term recursion relation. Also, all polynomials beyond a critical polynomial $`P_m(E)`$ factorize into the product of polynomial $`P_{n+1}(E)`$ and another arbitrary polynomial.
In this paper we suggest a generalization of Bender-Dunne approach to all possible one-dimentional quasi-exactly second order differential equations.
For this purpose, the succesful master function approach of references to exactly solvable models, is generalized to a master function of order up to four which gives all possible one-dimensional quasi-exactly solvable models, where Bender-Dunne model and Heun differential equation are among them.
The paper is organized as follows: In section II we show that we can generalize the usual quadratic master function to a master function of at most four order polynomials, then the most general quasi-exactly solvable differential operators related to generalized master function of degree $`k=3`$ and $`k=4`$ are given, respectively.
In section III, expanding their solutions in powers of $`x`$, we get 3-term and 4-term recursion relations among their coefficeints, where Bender-Dunne factorization follows through imposing the quasi-exactly solvability conditions of section II. At the end of this section we list all possible related quasi-exactly solvable differential equations for $`k=3`$ and $`k=4`$ in Tables I and II, respectively.
Finally at section IV, we derive all possible one-dimensional quasi-exactly solvable quantum Hamiltonian from the differential operators of section III, via prescription of references , where we have listed them at the end of section III, except for those which can given in terms of elliptic functions. Paper ends with a brief conclusion.
## 2 QUASI-EXACTLY SOLVABLE DIFFERENTIAL EQUATIONS ASSOCIATED WITH GENERALIZED MASTER FUNCTIONS
By generalizing master function of order up to two to polynomial of order up to k, together with the non-negative weight function $`W(x)`$, defined at interval (a,b) such that $`\frac{1}{W(x)}\frac{d}{dx}\left(A(x)W(x)\right)`$ to be a polynomial of degree at most $`(k1)`$, we can define the operator
$$L=\frac{1}{W(x)}\frac{d}{dx}\left(A(x)W(x)\frac{d}{dx}\right)+B(x),$$
(2-1)
where $`B(x)`$ is a polynomial of order up to $`(k2)`$. The interval $`(a,b)`$ is chosen so that, we have $`A(a)W(a)=A(b)W(b)=0`$.
It is straightforward to show that the above defined operator $`L`$ is a self adjoint linear operator which at most maps a given polynomial of order $`m`$ to another polynomial of order $`(m+k2)`$. Now, by an appropriate choice of $`B(x)`$ and weight function $`W(x)`$, the operator $`L`$ can have an invariant subspace of polynomials of order up to $`n`$. Then by choosing the set of orthogonal polynomials $`\{\varphi _0(x),\varphi _1(x),\mathrm{},\varphi _n(x)\}`$ defined in the interval $`(a,b)`$ with respect to the weight function $`W(x)`$:
$$_a^b\varphi _m(x)\varphi _n(x)W(x)𝑑x=0,\text{for}m=n$$
(2-2)
as the base, the matrix elements of the operator $`L`$ on this base will have the following block diagonal form:
$$L_{ij}=0,if\{inandjn+1\}or\{in+1andjn\}.$$
(2-3)
Since, according to the well known theorem of orthogonal polynomials, $`\varphi _n(x)`$ is orthogonal to any polynomial of order up to $`n1`$, therefore, for matrix $`L`$ we get
$$L=\left[\begin{array}{cc}M& 0\\ \multicolumn{2}{c}{}\\ 0& N\end{array}\right],$$
(2-4)
where $`M`$ is an $`(n+1)\times (n+1)`$ matrix with matrix elements
$$M_{ij}=_a^b𝑑xW(x)\varphi _i(x)L(x)\varphi _j(x),i,j=0,1,2,\mathrm{},n,$$
(2-5)
and $`N`$ is an infinite matrix element defined as above with $`i,jn+1`$.
The block diagonal form of the operator $`L`$ indicates that by diagonalizing the $`(n+1)\times (n+1)`$ matrix M, we can find $`(n+1)`$ eigen-values of the operator $`L`$ together with the related eigen-functions as linear functions of orthogonal polynomials $`\{\varphi _0(x),\varphi _1(x),\mathrm{},\varphi _n(x)\}`$ .
In order to determine the appropriate $`B(x)`$ and $`W(x)`$ for a given generalized master function $`A(x)`$, we Taylor expand those functions:
$$A(x)=\underset{i=0}{\overset{k}{}}\frac{A^{(i)}(0)}{i!}x^i,\text{where}A^{(i)}(0)=\frac{d^iA(x)}{dx^i}_{x=0}$$
(2-6)
$$\frac{\left(A(x)W(x)\right)^{}}{W(x)}=\underset{i=0}{\overset{k1}{}}\frac{\left(\frac{(AW)^{}}{W}\right)^{(i)}(0)}{i!}x^i,\text{ where}\left(\frac{(AW)^{}}{W}\right)^{(i)}(0)=\frac{d^i\left(\frac{(A(x)W(x))^{}}{W(x)}\right)}{dx^i}_{x=0}$$
(2-7)
$$B(x)=\underset{i=0}{\overset{k2}{}}\frac{B^{(i)}(0)}{i!}x^i,\text{where}B^{(i)}(0)=\frac{d^iB(x)}{dx^i}_{x=0}.$$
(2-8)
Then, the existence of invariant subspace of the polynomials of order $`n`$ of the operator $`L`$ leads to the following linear equation between the coefficients of above Taylor expansions:
$`{\displaystyle \frac{A^{(i+2)}}{(i+2)!}}l(l1){\displaystyle \frac{\left(\frac{(AW)^{}}{W}\right)^{(i+1)}}{(i+1)!}}l+{\displaystyle \frac{B^{(i)}}{i!}}=0,`$ (2-9)
where
$`\{\begin{array}{cccccc}l=n,& and& i=1,& 2,& \mathrm{},& k2\\ l=n1,& and& i=2,& 3,& \mathrm{},& k2\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ l=nk+4,& and& i=k3,& k2& & \\ l=nk+3,& and& i=k2& & & \end{array}.`$ (2-15)
The number of above equations, for a given value of $`k`$, is $`\frac{(k1)(k2)}{2}`$. If we are to determine only the unknown function $`B(x)`$ without having any further constraint on the weight function $`W(x)`$, then the above $`\frac{(k1)(k2)}{2}`$ equations should be satisfied with $`(k2)`$ coefficients of Taylor expansion of $`B(x)`$ as the only unknowns, since $`B^{(0)}`$ can be absorbed in the eigen-spectrum operator $`L`$. Therefore, we are left with $`(k2)`$ unkowns to be determined, where the compatibility of equations (2-9) require $`k=3`$ at most. On the other hand, if we add the coefficients of Taylor expansions of $`A(x)`$ and $`\frac{\left(A(x)W(x)\right)^{}}{W(x)}`$ to our list of unknowns, ( to be determined by solving equations (2-9) ), then their compatibility conditions require that:
$$3(k1)\frac{(k1)(k2)}{2},$$
(2-16)
or $`k8`$, where further investigations show that we can have at most $`k=4`$, since for $`k5`$ the coefficients $`A^{(k)}(0)`$ and $`\left(\frac{(AW)^{}}{W}\right)^{(k1)}(0)`$ will vanish. Below we summarize the above-mentioned discussion for $`k=3`$ and $`k=4`$, separately.
### 2.1 $`k=3`$
In this case, $`B(x)`$ is a second order polynomial where $`B^{(1)}`$ can be determined by solving equation (2-9):
$$B^{(1)}=\frac{n}{2}\left(\frac{A^{(3)}}{3}(n1)+\left(\frac{(AW)^{}}{W}\right)^{(2)}\right),$$
(2-17)
which is the only unknown in this case.
### 2.2 $`k=4`$
Again, the solving of equations (2-9) leads to:
$$B^{(1)}=\frac{n}{2}\left(\frac{A^{(3)}}{3}(n1)+\left(\frac{(AW)^{}}{W}\right)^{(2)}\right),$$
(2-18)
$$B^{(2)}=\frac{A^{(4)}}{12}n(n1),$$
(2-19)
and
$$\left(\frac{(AW)^{}}{W}\right)^{(3)}=\frac{A^{(4)}}{2}(n1).$$
(2-20)
Here, besides having constraint over second order polynomial $`B(x)`$, we have to put further constraints on the weight function $`W(x)`$ given in (2-14).
Definetly we can determine $`n+1`$ eigen-spectrum of the operator $`L`$, simply by diagonalizing the $`(n+1)\times (n+1)`$ matrix $`M`$, since it is a self-adjoint operator in Hilbert space of polynomials and it has a block diagonal form given in (2-4).
As we are going to see in the next section, we can determine its eigen-spectrum analytically, using some recursion relations.
## 3 RECURSION RELATIONS
In this section we show that the eigen-functions of the operator $`L`$ are a generating function for a new set of polynomials $`P_m(E)`$ where the eigen-function equation of the operator $`L`$ leads to the recursion relation between these polynomials. Quasi-exact solvable constraints (2-9) will lead to their factorization, that is, $`P_{n+N+1}(E)=P_{n+1}(E)Q_N(E)`$ for $`N0`$, where roots of polynomials $`P_{n+1}(E)`$ turn out to be the eigen-values of the operator $`L`$.
To achieve these results, first we expand $`\psi (x)`$, the eigen-function of $`L`$, as:
$$\psi (x)=\underset{m=0}{\overset{\mathrm{}}{}}P_m(E)x^m,$$
(3-21)
where eigen-function equation:
$$L\psi (x)=E\psi (x)$$
(3-22)
can be expressed as:
$$A(x)\underset{m=2}{\overset{\mathrm{}}{}}m(m1)P_m(E)x^{m2}\frac{\left(A(x)W(x)\right)^{}}{W(x)}\underset{m=1}{\overset{\mathrm{}}{}}mP_m(E)x^{m1}$$
$$+B(x)\underset{m=0}{\overset{\mathrm{}}{}}P_m(E)x^m=E\underset{m=0}{\overset{\mathrm{}}{}}P_m(E)x^m,$$
(3-23)
which leads to the following recursion relations for the coefficients $`P_m(E)`$:
$$\left(A^{(1)}(m+1)(m+2)+\left(\frac{(AW)^{}}{W}\right)^{(0)}(m+2)\right)P_{m+2}(E)$$
$$+\left(\frac{A^{(2)}}{2!}m(m+1)+\left(\frac{(AW)^{}}{W}\right)^{(1)}(m+1)+E\right)P_{m+1}(E)$$
$$+\left(\frac{A^{(3)}}{3!}m(m1)+\frac{\left(\frac{(AW)^{}}{W}\right)^{(2)}}{2!}mB^{(1)}\right)P_m(E)$$
$$+\left(\frac{A^{(4)}}{4!}(m1)(m2)+\frac{\left(\frac{(AW)^{}}{W}\right)^{(3)}}{3!}m\frac{B^{(2)}}{2!}\right)P_{m1}(E)=0.$$
(3-24)
Below we investigate recursion relations thus obtained for $`k=3`$ and $`k=4`$, separately.
### 3.1 $`k=3`$
In this case the 4-term general recursion relation reduse to the following 3-term recursion relation:
$$\left(A^{(1)}(m+1)(m+2)+\left(\frac{(AW)^{}}{W}\right)^{(0)}(m+2)\right)P_{m+2}(E)$$
$$+\left(\frac{A^{(2)}}{2!}m(m+1)+\left(\frac{(AW)^{}}{W}\right)^{(1)}(m+1)+E\right)P_{m+1}(E)$$
$$+\left(\frac{A^{(3)}}{3!}m(m1)+\frac{\left(\frac{(AW)^{}}{W}\right)^{(2)}}{2!}mB^{(1)}\right)P_m(E)=0.$$
(3-25)
In order to have finite eigen-spectrum, that is, quasi-integrable differential equation, the above recursion relation should be truncated for some value of $`m=n`$, which is obviously possible by an appropriate choice of:
$$B^{(1)}=\frac{n}{2}\left(\frac{A^{(3)}}{3}(n1)+\left(\frac{(AW)^{}}{W}\right)^{(2)}\right),$$
(3-26)
which is in agreement with the result of previous section given in (2-11).
Using the recursion relation (3-19), with $`B^{(1)}`$ given in (3-20), we get a factorization of polynomial $`P_{n+N+1}(E)`$ for $`N0`$ in terms of $`P_{n+1}(E)`$ as follows:
$$P_{n+N+1}(E)=P_{n+1}(E)Q_N(E)N0$$
(3-27)
where, by choosing the eigen-values $`E`$ as roots of polynomial $`P_{n+1}(E)`$, all polynomials of order higher than $`n`$ will vanish.
In order to determine corresponding eigen-functions, it is sufficient to evaluate $`P_m(E_i)`$ for $`m=0,\mathrm{\hspace{0.33em}1},\mathrm{\hspace{0.33em}2}\mathrm{}n`$ with $`E_i`$ as roots of $`P_{n+1}(E)`$, then eigen-function $`\psi _i(x)`$ corresponding to eigen-value $`E_i`$ can be written as:
$$\psi _i(x)=\underset{m=0}{\overset{n}{}}P_m(E_i)x^m,i=0,1\mathrm{},n.$$
(3-28)
The above eigen-functions are polynomials of order $`n`$, hence they can have at most $`n`$ roots in the interval $`(a,b)`$, where, according to the well known oscillation and comparison theorem of second-order linear differential equation , these numbers order the eigen-values according to the number of roots of corresponding eigen-functions. Therefore, we can say that the eigen-values thus obtained are the first $`n+1`$ eigen-values of the operator $`L`$.
Using the recursion relations (3-19), we can evaluate the polynomials $`P_m(E)`$ in terms of $`P_0(E)`$, where we have chosen $`P_0(E)=1`$. We have evaluated the first five polynomials which appear in the Appendix I.
As an illustration we give the results for $`A(x)=x`$ and $`n=3`$ with $`\alpha =1`$, $`\beta =0`$, $`\gamma =1`$ which is equivalent to the Bender-Dunne model:
$$P_1(E)=\frac{1}{2}E,$$
$$P_2(E)=1+\frac{1}{12}E^2,$$
$$P_3(E)=\frac{1}{4}E\frac{1}{144}E^3,$$
$$P_4(E)=\frac{1}{10}\frac{1}{48}E^2+\frac{1}{2880}E^4.$$
Obviously $`P_m(E)`$ have the parity of $`m`$.
By finding the 4-roots of $`P_4(E)`$ we determine the corresponding four eigen-values:
$$E_0=7.398556194,$$
$$E_1=2.293766823,$$
$$E_2=2.293766823,$$
$$E_3=7.398556194.$$
Finally for the coefficient $`P_m(E_i)`$ we get:
$$P_0(E_0)=1,P_1(E_0)=3.699278097,P_2(E_0)=3.561552813,P_3(E_0)=.962769686,$$
$$P_0(E_1)=1,P_1(E_1)=1.146883412,P_2(E_1)=.5615528135,P_3(E_1)=.4896337383,$$
$$P_0(E_2)=1,P_1(E_2)=1.146883412,P_2(E_2)=.5615528135,P_3(E_2)=.4896337383,$$
$$P_0(E_3)=1,P_1(E_3)=3.699278097,P_2(E_3)=3.561552813,P_3(E_3)=.962769686.$$
Using the above coefficients we can determine the corresponding eigen-functions through formula (3-22).
In Table I we give all quasi-exactly solvable operators which can be obtained by choosing different generalized master function of order 3. This Table contains all possible models corresponding to different choice of $`A(x)`$ up to translation and rescaling of variable $`x`$. Also by choosing $`A(x)`$ as a polynomial of up to second order with $`\gamma =0`$ we lead to the exactly solvable models of references .
### 3.2 $`k=4`$
Again in order to truncate the recursion relation (3-18) and to factorize polynomials $`P_{n+N+1}(E)`$ in terms of $`P_{n+1}(E)`$, we should have:
$$B^{(1)}=\frac{n}{2}\left(\frac{A^{(3)}}{3}(n1)+\left(\frac{(AW)^{}}{W}\right)^{(2)}\right),$$
(3-29)
$$\frac{B^{(2)}}{2!}=\frac{A^{(4)}}{4!}(n1)(n2)+\frac{\left(\frac{(AW)^{}}{W}\right)^{(3)}}{3!}n,$$
(3-30)
and
$$\frac{B^{(2)}}{2!}=\frac{A^{(4)}}{4!}n(n1)+\frac{\left(\frac{(AW)^{}}{W}\right)^{(3)}}{3!}(n+1).$$
(3-31)
Solving the above equations we get:
$$B^{(2)}=\frac{A^{(4)}}{12}n(n1),$$
(3-32)
and
$$\left(\frac{(AW)^{}}{W}\right)^{(3)}=\frac{A^{(4)}}{2}(n1).$$
(3-33)
The equations (3-23) , (3-26) and (3-27) are the same equations which are required in the reduction of the operator $`L`$ to its block diagonal form.
Again roots of polynomials $`P_{n+1}(E)`$ will correspond to $`n+1`$ eigen-values of the differential operator $`L`$ with eigen-functions which can be expressed in terms of $`P_m(E_i)`$ for $`mn`$, where polynomials $`P_m(E)`$ can be obtained from recursion relation by choosing $`P_0(E)=1`$ and $`P_1(E)=0`$, where we have given the first four polynomials in Appendix II.
In Table II we list all quasi-exactly differential operators which can be obtained from the generalized master function of order up to four.
Tables I and II contain all quasi-exactly second order differential equations which can be obtained from Lie algebraic methods . For example we get the Bender-Dunne model for the choice of $`A(x)=x`$ and $`\beta =0`$, and similarly, we get the Heun differential operator ( Fuxian equation with four regular singular point ) , for the choice of $`A(x)=x(x1)(xa)`$.
## 4 QUASI-EXACTLY POTENTIAL ASSOCIATED WITH GENERALIZED MASTER FUNCTION
As in references , writting:
$$\psi (t)=A^{1/4}(x)W^{1/2}(x)\varphi (x),$$
(4-34)
with a change of variable $`\frac{dx}{dt}=\sqrt{A(x)}`$, the eigen-value equation of the operator $`L`$ reduces to the schrodinger equation:
$$H(t)\psi (t)=E\psi (t),$$
(4-35)
with the same eigen-value E and $`\psi (t)`$ given in (4-29), in terms of the eigen function of $`L`$, where $`H(t)=\frac{d^2}{dt^2}+V(t)`$ is the similarity transformation of $`L(x)`$ defined as:
$$H(t)=A^{1/4}(x)W^{1/2}(x)L(x)A^{1/4}(x)W^{1/2}(x)$$
(4-36)
with:
$$V(t)=\frac{3}{16}\frac{\dot{A}^2(t)}{A^2(t)}\frac{1}{4}\frac{\dot{W}^2(t)}{W^2(t)}+\frac{1}{4}\frac{\dot{A}(t)\dot{W}(t)}{A(t)W(t)}+\frac{1}{4}\frac{\ddot{A}(t)}{A(t)}+\frac{1}{2}\frac{\ddot{W}(t)}{W(t)}+B(t).$$
(4-37)
It is also straightforward to show that:
$$𝑑t\varphi (t)H(t)\psi (t)=_a^b𝑑xW(x)\psi (x)L(x)\psi (x).$$
(4-38)
Hence block diagonalization of $`L`$ leads to block-diagonalization of $`H`$.
As an illustration we give below an example with $`A(x)=x^3`$, weight function $`W(x)=x^\alpha e^{\beta /x^2\gamma /x}`$ and interval $`[0,\mathrm{})`$, where $`\alpha <3`$ and $`\beta ,\gamma >0`$.
From a change of variable $`dx/dt=\sqrt{A(x)}`$, we get $`x(t)=4/t^2`$, hence using equation (4-31) we have for potential $`V(t)`$:
$$V_n(t)=\frac{\gamma }{2}\left(\alpha +1\right)+\left(\frac{15}{4}+\alpha ^2+4n\alpha +4\alpha +4n^2+8n\right)\frac{1}{t^2}+\frac{1}{4}\left(\alpha \beta +\frac{1}{4}\gamma ^2\right)t^2$$
$$+\frac{\beta \gamma }{16}t^4+\frac{\beta ^2}{64}t^6$$
(4-39)
Below we give a list of quasi-exact solvable potentials, except for those potentials which can be expressed in terms of elliptic functions, since in this case we get rather long expressions for them:
$$A(x)=x,x(t)=\frac{t^2}{4}$$
$$V_n(t)=\frac{\beta }{2}\left(\alpha +1\right)+\left(\alpha ^2\frac{1}{4}\right)\frac{1}{t^2}+\frac{1}{2}\left(\frac{\beta ^2}{8}+\gamma \left(n+1+\frac{\alpha }{2}\right)\right)t^2+\frac{\beta \gamma }{16}t^4+\frac{\gamma ^2}{64}t^6,$$
$$A(x)=x^2,x(t)=e^t$$
$$V_n(t)=\frac{1}{4}\left(1+\alpha ^22\beta \gamma +2\alpha 2\alpha \beta e^t+\beta ^2e^{2t}+\gamma ^2e^{2t}+2\left(2\gamma +2n\gamma +\alpha \gamma \right)e^t\right),$$
$$A(x)=x(1x),x(t)=\frac{1+sin(t)}{2}$$
$$V_n(t)=\frac{1}{2}\left(+n\gamma \alpha \beta \beta \alpha +\frac{1}{2}\left(\beta \gamma \alpha ^2\beta ^2\alpha \gamma 1\right)+\left(\frac{\alpha \gamma }{2}+\gamma +\frac{\beta \gamma }{2}+n\gamma \right)sin(t)\right)$$
$$\frac{1}{2}\left(\alpha ^2+\beta ^2\frac{1}{2}+\left(\beta ^2\alpha ^2\right)sin(t)\right)\frac{1}{cos^2(t)}+\frac{\gamma ^2}{16}cos^2(t),$$
$$A(x)=x^3,x(t)=\frac{4}{t^2}$$
$$V_n(t)=\frac{\gamma }{2}\left(\alpha +1\right)+\left(\frac{15}{4}+\alpha ^2+4n\alpha +4\alpha +4n^2+8n\right)\frac{1}{t^2}+\frac{1}{4}\left(\alpha \beta +\frac{\gamma ^2}{4}\right)t^2$$
$$+\frac{\beta \gamma }{16}t^4+\frac{\beta ^2}{64}t^6,$$
$$A(x)=x^2(1x),x(t)=1tanh^2(\frac{t}{2})$$
$$V_n(t)=\left(\frac{1}{cosh^2(t)1}\right)((2n^2+2+2n\alpha +\frac{\alpha ^2}{2}+4n+\alpha \beta +2\alpha +\frac{\alpha \gamma }{4}+2n\beta +2\beta )cosh(t)$$
$$+\frac{1}{2}\left(\frac{\gamma ^2}{4}\frac{\alpha \gamma }{2}+\frac{1}{2}\gamma +\alpha \beta \gamma +\frac{\alpha ^2}{2}\right)cosh^2+\frac{\alpha \gamma }{4}cosh^3+\frac{\gamma ^2}{16}cosh^4$$
$$+(\frac{\gamma }{2}+4n+\alpha \beta +\frac{3\alpha }{2}+\beta ^2+2\beta +\frac{3}{2}+\frac{\alpha ^2}{4}+\frac{\alpha \gamma }{4}+\frac{\beta \gamma }{2}+\frac{\gamma ^2}{16}+2n\beta +2n^2+2n\alpha )),$$
$$A(x)=x^4,x(t)=\frac{1}{t}$$
$$V_n(t)=\left(\frac{\delta ^2}{4}+\gamma +2n\gamma \right)+\left(\gamma \delta +3n\beta +3\beta \right)t+\left(\frac{3\beta \delta }{2}+\gamma ^2\right)t^2+3\beta \gamma t^3+\frac{9\beta ^2}{4}t^4,$$
$$A(x)=x^3(1x),x(t)=\frac{4}{4+t^2}$$
$$V_n(t)=\left(\frac{\gamma }{2}+\delta +\beta \delta +\frac{\beta \gamma }{2}+n\gamma \right)+\left(\beta ^2\frac{1}{4}\right)\frac{1}{t^2}$$
$$+\frac{1}{2}\left(n\delta +\frac{\delta ^2}{2}\frac{\beta \delta }{2}\delta +\frac{\gamma ^2}{8}+\frac{\gamma \delta }{2}\right)t^2+\frac{\delta }{8}\left(\frac{\gamma }{2}+\delta \right)t^4+\frac{\delta ^2}{64}t^6,$$
$$A(x)=x^2(1+x^2),x(t)=\frac{1}{sinh(t)}$$
$$V_n(t)=\left(n+n^2\frac{\gamma \delta }{2}+\frac{1}{4}+2n\beta +\beta +\beta ^2\left(n\gamma +\gamma +\beta \gamma \right)sinh(t)\right)$$
$$+\left(\frac{\delta ^2}{4}+\beta \delta sinh(t)\beta ^2+\frac{1}{4}\right)\frac{1}{cosh^2(t)}+\frac{\gamma ^2}{4}cosh^2(t),$$
$$A(x)=x^2(1x^2),x(t)=\frac{1}{cosh(t)}$$
$$V_n(t)=\left(\frac{1}{cosh^2(t)1}\right)((\frac{\gamma }{2}n\frac{\beta }{2}+\frac{\gamma ^2}{4}+\frac{\delta ^2}{4}\frac{\beta \gamma }{2}\frac{\beta \delta }{2}+\frac{\gamma \delta }{2}\frac{1}{2}n\beta +\frac{\beta ^2}{4}n\gamma n^2)$$
$$+\left(\frac{\beta \delta }{2}+n\beta \frac{\gamma \delta }{2}+n\gamma +\frac{\gamma }{2}+\frac{1}{4}+\frac{\beta }{2}+\frac{\beta \gamma }{2}\frac{\delta ^2}{2}+n^2+\frac{\beta ^2}{4}+n+\frac{\gamma ^2}{4}\right)cosh(t)^2$$
$$+(\frac{\gamma \delta }{2}n\delta \delta \frac{\gamma ^2}{2}\frac{\beta \delta }{2}+\frac{\beta ^2}{2})cosh(t)+(\delta +n\delta +\frac{\beta \delta }{2}+\frac{\gamma \delta }{2})cosh^3(t)+\frac{\delta ^2}{4}cosh^4(t)),$$
$$A(x)=x^2(1x)^2,x(t)=\frac{e^t}{1+e^t}$$
$$V_n(t)=\left(\frac{1}{\left(e^t+1\right)^4S_6}\right)(\frac{\delta ^2}{4}S_4+\delta (\delta \frac{\beta }{2})S_5$$
$$+\left(\frac{1}{4}\frac{\gamma \delta }{2}+n^2+\frac{\beta ^2}{4}n\gamma +\frac{\beta }{2}+n2\beta \delta +\frac{3\delta ^2}{2}+n\beta \right)S_6$$
$$\left(3\beta \delta +\frac{\beta \gamma }{2}3n\gamma +\gamma +4n+\delta ^22\gamma \delta +\beta ^2+1+4n^2+2\beta +4n\beta \right)S_7$$
$$+\left(3\beta +6n+\frac{3\beta ^2}{2}3\gamma \delta +\frac{3}{2}+\frac{\delta ^2}{4}+\frac{\gamma ^2}{4}+6n^22\beta \delta +4\gamma +6n\beta +2\beta \gamma 2n\gamma \right)S_8$$
$$+\left(4n^2\frac{\beta \delta }{2}2\gamma \delta +3\beta \gamma +1+2n\gamma +4n+\beta ^2+2\beta +\gamma ^2+6\gamma +4n\beta \right)S_9$$
$$+\left(\frac{\gamma \delta }{2}+\frac{3\gamma ^2}{2}+4\gamma +n+2\beta \gamma +3n\gamma +\frac{1}{4}+\frac{\beta ^2}{4}+\frac{\beta }{2}+n\beta +n^2\right)S_{10}$$
$$+(n\gamma +\gamma ^2+\frac{\beta \gamma }{2}+\gamma )S_{11}+\frac{\gamma ^2}{4}S_{12}),$$
where $`S_k`$ is defined as:
$$S_k=exp\left(2\left(\frac{kt}{2}+\left(\gamma e^t+\gamma +\delta \left(e^t+1\right)\right)\right)\right),k=4,5,\mathrm{},12.$$
## 5 CONCLUSION
As we saw by introducing of master function $`A(x)`$ as a polynomial of order at most four, we could obtain all quasi-exactly second order differential equations. It is shown that the eigen-equation relation $`L\mathrm{\Psi }(x)=E\mathrm{\Psi }(x)`$ generates a set of polynomials $`P_m(E)`$, where these polynomials satisfy 3-term and 4-term recursion relations for master function of at most three and four, respectively. Finally the quasi-exactly solvablity leads to factorization of polynomials $`P_{n+N+1}(E)`$ for $`N0`$ in terms of $`P_{n+1}(E)`$, where by determining the roots of $`P_{n+1}(E)`$ we can determine first $`n+1`$ eigen-values of these quasi-exactly solvable differential equations.
APPENDIX I
The First Five Polynomials $`P_n(E)`$, For $`k=3`$
To abbreviate, we set $`F^{(i)}=(\frac{AW^{}}{W})^{(i)}`$,
$$P_1(E)=\frac{E}{F^{(\mathit{0})}},$$
$$P_2(E)=\frac{1}{2}\frac{EF^{(\mathit{1})}+E^2+B^{(\mathit{1})}F^{(\mathit{0})}}{F^{(\mathit{0})}\left(A^{(\mathit{1})}+F^{(\mathit{0})}\right)},$$
$$P_3(E)=(A^{(\mathit{2})}EF^{(\mathit{1})}A^{(\mathit{2})}E^2A^{(\mathit{2})}B^{(\mathit{1})}F^{(\mathit{0})}2EF_{}^{(\mathit{1})}{}_{}{}^{2}3F^{(\mathit{1})}E^22F^{(\mathit{1})}B^{(\mathit{1})}F^{(\mathit{0})}E^3$$
$$3EB^{(\mathit{1})}F^{(\mathit{0})}+EF^{(\mathit{2})}A^{(\mathit{1})}+EF^{(\mathit{2})}F^{(\mathit{0})}2EB^{(\mathit{1})}A^{(\mathit{1})})$$
$$/\left(6F^{(0)}\left(A^{(1)}+F^{(0)}\right)\left(2A^{(1)}+F^{(0)}\right)\right),$$
$$P_4(E)=(3A_{}^{(2)}{}_{}{}^{2}E^23B_{}^{(1)}{}_{}{}^{2}F_{}^{(0)}{}_{}{}^{2}6B_{}^{(1)}{}_{}{}^{2}F^{(0)}A^{(1)}+3F^{(2)}B^{(1)}F_{}^{(0)}{}_{}{}^{2}+A^{(3)}B^{(1)}F_{}^{(0)}{}_{}{}^{2}$$
$$+A^{(\mathit{3})}E^2F^{(\mathit{0})}+2A^{(\mathit{3})}E^2A^{(\mathit{1})}8E^2B^{(\mathit{1})}A^{(\mathit{1})}+4E^2F^{(\mathit{2})}F^{(\mathit{0})}+7E^2F^{(\mathit{2})}A^{(\mathit{1})}6E^2B^{(\mathit{1})}F^{(\mathit{0})}$$
$$6F_{}^{(\mathit{1})}{}_{}{}^{2}B^{(\mathit{1})}F^{(\mathit{0})}13A^{(\mathit{2})}F^{(\mathit{1})}E^29A^{(\mathit{2})}EF_{}^{(\mathit{1})}{}_{}{}^{2}3A_{}^{(\mathit{2})}{}_{}{}^{2}B^{(\mathit{1})}F^{(\mathit{0})}3A_{}^{(\mathit{2})}{}_{}{}^{2}EF^{(\mathit{1})}$$
$$4A^{(\mathit{2})}E^36EF_{}^{(\mathit{1})}{}_{}{}^{3}11F_{}^{(\mathit{1})}{}_{}{}^{2}E^26F^{(\mathit{1})}E^39A^{(\mathit{2})}F^{(\mathit{1})}B^{(\mathit{1})}F^{(\mathit{0})}6A^{(\mathit{2})}EB^{(\mathit{1})}A^{(\mathit{1})}$$
$$+3A^{(\mathit{2})}EF^{(\mathit{2})}F^{(\mathit{0})}+3A^{(\mathit{2})}EF^{(\mathit{2})}A^{(\mathit{1})}10A^{(\mathit{2})}EB^{(\mathit{1})}F^{(\mathit{0})}12F^{(\mathit{1})}EB^{(\mathit{1})}A^{(\mathit{1})}+6F^{(\mathit{1})}EF^{(\mathit{2})}F^{(\mathit{0})}$$
$$+9F^{(\mathit{1})}EF^{(\mathit{2})}A^{(\mathit{1})}14F^{(\mathit{1})}EB^{(\mathit{1})}F^{(\mathit{0})}+2A^{(\mathit{3})}B^{(\mathit{1})}F^{(\mathit{0})}A^{(\mathit{1})}+A^{(\mathit{3})}EF^{(\mathit{1})}F^{(\mathit{0})}$$
$$+2A^{(\mathit{3})}EF^{(\mathit{1})}A^{(\mathit{1})}+6F^{(\mathit{2})}B^{(\mathit{1})}F^{(\mathit{0})}A^{(\mathit{1})}E^4)$$
$$/\left(24F^{(\mathit{0})}\left(A^{(\mathit{1})}+F^{(\mathit{0})}\right)\left(2A^{(\mathit{1})}+F^{(\mathit{0})}\right)\left(3A^{(\mathit{1})}+F^{(\mathit{0})}\right)\right),$$
$$P_5(E)=(46EA^{(\mathit{3})}B^{(\mathit{1})}F^{(\mathit{0})}A^{(\mathit{1})}88EF^{(\mathit{2})}B^{(\mathit{1})}F^{(\mathit{0})}A^{(\mathit{1})}+16A^{(\mathit{3})}EF^{(\mathit{2})}A^{(\mathit{1})}F^{(\mathit{0})}$$
$$+24EF_{}^{(\mathit{2})}{}_{}{}^{2}A^{(\mathit{1})}F^{(\mathit{0})}48F^{(\mathit{2})}EB^{(\mathit{1})}A_{}^{(\mathit{1})}{}_{}{}^{2}+18EF_{}^{(\mathit{2})}{}_{}{}^{2}A_{}^{(\mathit{1})}{}_{}{}^{2}+6EF_{}^{(\mathit{2})}{}_{}{}^{2}F_{}^{(\mathit{0})}{}_{}{}^{2}+24EB_{}^{(\mathit{1})}{}_{}{}^{2}A_{}^{(\mathit{1})}{}_{}{}^{2}$$
$$+12A^{(\mathit{3})}EF^{(\mathit{2})}A_{}^{(\mathit{1})}{}_{}{}^{2}+4A^{(\mathit{3})}EF^{(\mathit{2})}F_{}^{(\mathit{0})}{}_{}{}^{2}24A^{(\mathit{3})}EB^{(\mathit{1})}A_{}^{(\mathit{1})}{}_{}{}^{2}32F^{(\mathit{1})}A^{(\mathit{3})}B^{(\mathit{1})}F^{(\mathit{0})}A^{(\mathit{1})}$$
$$60F^{(\mathit{1})}F^{(\mathit{2})}B^{(\mathit{1})}F^{(\mathit{0})}A^{(\mathit{1})}+50EB_{}^{(\mathit{1})}{}_{}{}^{2}F^{(\mathit{0})}A^{(\mathit{1})}25EF^{(\mathit{2})}B^{(\mathit{1})}F_{}^{(\mathit{0})}{}_{}{}^{2}13EA^{(\mathit{3})}B^{(\mathit{1})}F_{}^{(\mathit{0})}{}_{}{}^{2}$$
$$46F^{(\mathit{1})}A^{(\mathit{3})}E^2A^{(\mathit{1})}+80F^{(\mathit{1})}E^2B^{(\mathit{1})}A^{(\mathit{1})}40F^{(\mathit{1})}E^2F^{(\mathit{2})}F^{(\mathit{0})}91F^{(\mathit{1})}E^2F^{(\mathit{2})}A^{(\mathit{1})}$$
$$+50F^{(\mathit{1})}E^2B^{(\mathit{1})}F^{(\mathit{0})}54A^{(\mathit{2})}F^{(\mathit{1})}EF^{(\mathit{2})}F^{(\mathit{0})}84A^{(\mathit{2})}F^{(\mathit{1})}EF^{(\mathit{2})}A^{(\mathit{1})}+137A^{(\mathit{2})}F^{(\mathit{1})}EB^{(\mathit{1})}F^{(\mathit{0})}$$
$$24A^{(\mathit{2})}A^{(\mathit{3})}B^{(\mathit{1})}F^{(\mathit{0})}A^{(\mathit{1})}10A^{(\mathit{2})}A^{(\mathit{3})}EF^{(\mathit{1})}F^{(\mathit{0})}24A^{(\mathit{2})}A^{(\mathit{3})}EF^{(\mathit{1})}A^{(\mathit{1})}$$
$$54A^{(\mathit{2})}F^{(\mathit{2})}B^{(\mathit{1})}F^{(\mathit{0})}A^{(\mathit{1})}+20F^{(\mathit{1})}B_{}^{(\mathit{1})}{}_{}{}^{2}F_{}^{(\mathit{0})}{}_{}{}^{2}5A^{(\mathit{3})}E^3F^{(\mathit{0})}14A^{(\mathit{3})}E^3A^{(\mathit{1})}+20E^3B^{(\mathit{1})}A^{(\mathit{1})}$$
$$10E^3F^{(\mathit{2})}F^{(\mathit{0})}25E^3F^{(\mathit{2})}A^{(\mathit{1})}+10E^3B^{(\mathit{1})}F^{(\mathit{0})}+15EB_{}^{(\mathit{1})}{}_{}{}^{2}F_{}^{(\mathit{0})}{}_{}{}^{2}+66A^{(\mathit{2})}E^2B^{(\mathit{1})}A^{(\mathit{1})}$$
$$33A^{(\mathit{2})}E^2F^{(\mathit{2})}F^{(\mathit{0})}63A^{(\mathit{2})}E^2F^{(\mathit{2})}A^{(\mathit{1})}+50A^{(\mathit{2})}E^2B^{(\mathit{1})}F^{(\mathit{0})}+72A^{(\mathit{2})}F_{}^{(\mathit{1})}{}_{}{}^{2}B^{(\mathit{1})}F^{(\mathit{0})}$$
$$+72F_{}^{(\mathit{1})}{}_{}{}^{2}EB^{(\mathit{1})}A^{(\mathit{1})}36F_{}^{(\mathit{1})}{}_{}{}^{2}EF^{(\mathit{2})}F^{(\mathit{0})}72F_{}^{(\mathit{1})}{}_{}{}^{2}EF^{(\mathit{2})}A^{(\mathit{1})}+70F_{}^{(\mathit{1})}{}_{}{}^{2}EB^{(\mathit{1})}F^{(\mathit{0})}$$
$$12A^{(\mathit{3})}EF_{}^{(\mathit{1})}{}_{}{}^{2}F^{(\mathit{0})}32A^{(\mathit{3})}EF_{}^{(\mathit{1})}{}_{}{}^{2}A^{(\mathit{1})}+48F^{(\mathit{1})}B_{}^{(\mathit{1})}{}_{}{}^{2}F^{(\mathit{0})}A^{(\mathit{1})}24F^{(\mathit{1})}F^{(\mathit{2})}B^{(\mathit{1})}F_{}^{(\mathit{0})}{}_{}{}^{2}$$
$$12F^{(\mathit{1})}A^{(\mathit{3})}B^{(\mathit{1})}F_{}^{(\mathit{0})}{}_{}{}^{2}17F^{(\mathit{1})}A^{(\mathit{3})}E^2F^{(\mathit{0})}24A^{(\mathit{2})}F^{(\mathit{2})}B^{(\mathit{1})}F_{}^{(\mathit{0})}{}_{}{}^{2}10A^{(\mathit{2})}A^{(\mathit{3})}E^2F^{(\mathit{0})}$$
$$24A^{(\mathit{2})}A^{(\mathit{3})}E^2A^{(\mathit{1})}10A^{(\mathit{2})}A^{(\mathit{3})}B^{(\mathit{1})}F_{}^{(\mathit{0})}{}_{}{}^{2}+108A^{(\mathit{2})}F^{(\mathit{1})}EB^{(\mathit{1})}A^{(\mathit{1})}$$
$$+10F^{(\mathit{1})}E^4+18A_{}^{(\mathit{2})}{}_{}{}^{3}E^2+27A_{}^{(\mathit{2})}{}_{}{}^{2}E^3+10A^{(\mathit{2})}E^4+24EF_{}^{(\mathit{1})}{}_{}{}^{4}+50F_{}^{(\mathit{1})}{}_{}{}^{3}E^2$$
$$+35F_{}^{(\mathit{1})}{}_{}{}^{2}E^3+E^5+93A_{}^{(\mathit{2})}{}_{}{}^{2}F^{(\mathit{1})}E^2+66A_{}^{(\mathit{2})}{}_{}{}^{2}EF_{}^{(\mathit{1})}{}_{}{}^{2}+18A_{}^{(\mathit{2})}{}_{}{}^{3}B^{(\mathit{1})}F^{(\mathit{0})}+18A_{}^{(\mathit{2})}{}_{}{}^{3}EF^{(\mathit{1})}$$
$$+22A^{(\mathit{2})}B_{}^{(\mathit{1})}{}_{}{}^{2}F_{}^{(\mathit{0})}{}_{}{}^{2}+72A^{(\mathit{2})}EF_{}^{(\mathit{1})}{}_{}{}^{3}+127A^{(\mathit{2})}F_{}^{(\mathit{1})}{}_{}{}^{2}E^2+65A^{(\mathit{2})}F^{(\mathit{1})}E^3$$
$$+24F_{}^{(\mathit{1})}{}_{}{}^{3}B^{(\mathit{1})}F^{(\mathit{0})}+66A_{}^{(\mathit{2})}{}_{}{}^{2}F^{(\mathit{1})}B^{(\mathit{1})}F^{(\mathit{0})}+36A_{}^{(\mathit{2})}{}_{}{}^{2}EB^{(\mathit{1})}A^{(\mathit{1})}$$
$$18A_{}^{(\mathit{2})}{}_{}{}^{2}EF^{(\mathit{2})}F^{(\mathit{0})}18A_{}^{(\mathit{2})}{}_{}{}^{2}EF^{(\mathit{2})}A^{(\mathit{1})}+63A_{}^{(\mathit{2})}{}_{}{}^{2}EB^{(\mathit{1})}F^{(\mathit{0})}+48A^{(\mathit{2})}B_{}^{(\mathit{1})}{}_{}{}^{2}F^{(\mathit{0})}A^{(\mathit{1})})$$
$$/\left(120F^{(\mathit{0})}\left(A^{(\mathit{1})}+F^{(\mathit{0})}\right)\left(2A^{(\mathit{1})}+F^{(\mathit{0})}\right)\left(3A^{(\mathit{1})}+F^{(\mathit{0})}\right)\left(4A^{(\mathit{1})}+F^{(\mathit{0})}\right)\right).$$
APPENDIX II
The First Four Polynomials $`P_n(E)`$, For $`k=4`$
$$P_1(E)=\frac{E}{F^{(0)}},$$
$$P_2(E)=\frac{1}{2}\frac{EF^{(\mathit{1})}+E^2B^{(\mathit{1})}F^{(\mathit{0})}}{F^{(\mathit{0})}\left(A^{(\mathit{1})}+F^{(\mathit{0})}\right)},$$
$$P_3(E)=(A^{(\mathit{2})}EF^{(\mathit{1})}+A^{(\mathit{2})}E^2A^{(\mathit{2})}B^{(\mathit{1})}F^{(\mathit{0})}+2EF_{}^{(\mathit{1})}{}_{}{}^{2}+3F^{(\mathit{1})}E^22F^{(\mathit{1})}B^{(\mathit{1})}F^{(\mathit{0})}$$
$$+E^3+EB^{(\mathit{1})}F^{(\mathit{0})}EF^{(\mathit{2})}A^{(\mathit{1})}EF^{(\mathit{2})}F^{(\mathit{0})}+2EB^{(\mathit{1})}A^{(\mathit{1})}+B^{(\mathit{2})}F^{(\mathit{0})}A^{(\mathit{1})}+B^{(\mathit{2})}F_{}^{(\mathit{0})}{}_{}{}^{2})$$
$$/\left(6F^{(\mathit{0})}\left(A^{(\mathit{1})}+F^{(\mathit{0})}\right)\left(2A^{(\mathit{1})}+F^{(\mathit{0})}\right)\right),$$
$$P_4(E)=(3A^{(\mathit{2})}B^{(\mathit{2})}F_{}^{(\mathit{0})}{}_{}{}^{2}+A^{(\mathit{3})}B^{(\mathit{1})}F_{}^{(\mathit{0})}{}_{}{}^{2}+9A^{(\mathit{2})}EF_{}^{(\mathit{1})}{}_{}{}^{2}3A^{(\mathit{2})}EF^{(\mathit{2})}A^{(\mathit{1})}$$
$$+2A^{(\mathit{2})}EB^{(\mathit{1})}F^{(\mathit{0})}9A^{(\mathit{2})}F^{(\mathit{1})}B^{(\mathit{1})}F^{(\mathit{0})}2EB^{(\mathit{2})}F_{}^{(\mathit{0})}{}_{}{}^{2}+3A^{(\mathit{2})}B^{(\mathit{2})}F^{(\mathit{0})}A^{(\mathit{1})}+6A^{(\mathit{2})}EB^{(\mathit{1})}A^{(\mathit{1})}$$
$$3A^{(\mathit{2})}EF^{(\mathit{2})}F^{(\mathit{0})}9F^{(\mathit{1})}EF^{(\mathit{2})}A^{(\mathit{1})}+4F^{(\mathit{1})}EB^{(\mathit{1})}F^{(\mathit{0})}3B_{}^{(\mathit{1})}{}_{}{}^{2}F_{}^{(\mathit{0})}{}_{}{}^{2}8EB^{(\mathit{2})}F^{(\mathit{0})}A^{(\mathit{1})}$$
$$2A^{(\mathit{3})}EF^{(\mathit{1})}A^{(\mathit{1})}A^{(\mathit{3})}EF^{(\mathit{1})}F^{(\mathit{0})}+2A^{(\mathit{3})}B^{(\mathit{1})}F^{(\mathit{0})}A^{(\mathit{1})}+6F^{(\mathit{2})}B^{(\mathit{1})}F^{(\mathit{0})}A^{(\mathit{1})}+3EF^{(\mathit{3})}A^{(\mathit{1})}F^{(\mathit{0})}$$
$$6F^{(\mathit{1})}EF^{(\mathit{2})}F^{(\mathit{0})}+3A_{}^{(\mathit{2})}{}_{}{}^{2}E^2+4A^{(\mathit{2})}E^3+6EF_{}^{(\mathit{1})}{}_{}{}^{3}+11F_{}^{(\mathit{1})}{}_{}{}^{2}E^2+6F^{(\mathit{1})}E^3+3A_{}^{(\mathit{2})}{}_{}{}^{2}EF^{(\mathit{1})}$$
$$3A_{}^{(\mathit{2})}{}_{}{}^{2}B^{(\mathit{1})}F^{(\mathit{0})}+13A^{(\mathit{2})}F^{(\mathit{1})}E^26F_{}^{(\mathit{1})}{}_{}{}^{2}B^{(\mathit{1})}F^{(\mathit{0})}+3F^{(\mathit{1})}B^{(\mathit{2})}F_{}^{(\mathit{0})}{}_{}{}^{2}+4E^2B^{(\mathit{1})}F^{(\mathit{0})}$$
$$7E^2F^{(\mathit{2})}A^{(\mathit{1})}4E^2F^{(\mathit{2})}F^{(\mathit{0})}+8E^2B^{(\mathit{1})}A^{(\mathit{1})}2A^{(\mathit{3})}E^2A^{(\mathit{1})}A^{(\mathit{3})}E^2F^{(\mathit{0})}+3F^{(\mathit{2})}B^{(\mathit{1})}F_{}^{(\mathit{0})}{}_{}{}^{2}$$
$$6B_{}^{(\mathit{1})}{}_{}{}^{2}F^{(\mathit{0})}A^{(\mathit{1})}+2EF^{(\mathit{3})}A_{}^{(\mathit{1})}{}_{}{}^{2}+EF^{(\mathit{3})}F_{}^{(\mathit{0})}{}_{}{}^{2}6EB^{(\mathit{2})}A_{}^{(\mathit{1})}{}_{}{}^{2}+E^4+3F^{(\mathit{1})}B^{(\mathit{2})}F^{(\mathit{0})}A^{(\mathit{1})}$$
$$+12F^{(\mathit{1})}EB^{(\mathit{1})}A^{(\mathit{1})})/(24F^{(\mathit{0})}(A^{(\mathit{1})}+F^{(\mathit{0})})(2A^{(\mathit{1})}+F^{(\mathit{0})})(3A^{(\mathit{1})}+F^{(\mathit{0})}))$$
ACKNOWLEDGEMENT
We wish to thank Dr. S. K. A. Seyed Yagoobi for his careful reading the article and for his constructive comments.
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# Space-Time Invariant Measures, Entropy, and Dimension for Stochastic Ginzburg–Landau Equations
## 1 Introduction
The use of dynamical system techniques and ideas in the study of extended partial differential equations has proved extremely fruitful in the past, see for example P. Collet’s talk at ICM’98, \[C2\] (where he also emphasises the limitations of such an approach). However, until now only results using topological or geometric properties of the dynamics have been used (like invariant manifolds, bifurcation theory, topological entropy, Hausdorff dimension). That is to say, extended dynamical systems are usually regarded as topological dynamical systems. In contrast, most of the very deep results in finite dimensional dynamical systems use measure-theoretic ideas, namely ergodic theory (as advocated for instance in the review by L.-S. Young at the ICMP in 1997, \[Y2\]).
One of the favourite models of infinite dimensional dynamical systems studied recently is the Ginzburg–Landau equation. It appears as a generic normal form describing the amplitude of periodic bifurcated solutions (see \[C2\]) and it is also believed to be a good example of spatio-temporally chaotic dynamics \[LO\]. It is known that its attractor is infinite-dimensional and has positive $`\epsilon `$–entropy (see \[CE1, CE2, CE3, Ro\]).
Here we propose to use random perturbations to obtain, by probabilistic techniques, the existence of invariant measures for the corresponding random dynamical system. The existence result is based on the observation by J. Ginibre and G. Velo that the Ginzburg–Landau equation has global solutions both in uniformly local Sobolev spaces and in local $`\mathrm{L}_{}^2`$ space. Their proofs go through to the stochastic case without much effort, if we assume the noise to be smooth in space. Since uniformly local Sobolev spaces of sufficiently high order are compactly embedded into local $`\mathrm{L}_{}^2`$ space, we get the Feller property of the semi-group and the tightness of the Cesàro means therefore existence of an invariant measure (by standard arguments for stochastic differential equations, see \[DZ1\]). These measures are also translation invariant, because the noise, the deterministic part of the equation and the spaces used are all translation invariant. We refer to the property of being invariant under the time evolution as “stationarity” and the invariance under space translations as “homogeneity” of the measure, following Vishik and Fursikov \[VF\].
In a second part of the paper we define the topological entropy and the measure-theoretic entropy, or rather their spatial densities, since both quantities are extensive (this has been discovered in this context by Collet and Eckmann in \[CE1\], see e.g. \[Ru\] for earlier similar ideas). Usual inequalities from ergodic theory can be proved in this case and the Collet-Eckmann bound on the topological entropy is also valid (see \[CE2\]).
The paper is organised as follows: In Section 2, we set the model and the functional analysis background needed for the remainder of the paper. The main results of the paper are summarised in Section 3. In Section 4 we obtain uniform bounds on the solutions in Sobolev spaces, these bounds being then used in Section 5 to prove the existence of invariant measures. Section 6 is devoted to the results on existence and properties of the (measure-theoretic and topological) entropies. Various technical proofs have been relegated to Sections 712.
We finish this introduction by commenting on the fact that many new results on invariant measures for nonlinear PDEs have recently appeared. We mention for instance \[BKL, FM, EH, KS, Ku, Ma, S\]. To the best of our knowledge the present work is the first where the model considered enjoys: infinite volume (hence continuous spectrum without gap), genuinely nonlinear interaction, homogeneous noise (hence infinite supply of energy at each time), non-trivial deterministic dynamics (e.g. the attractor of the deterministic Ginzburg–Landau is infinite dimensional). However, we are still unable to prove uniqueness of the invariant measure (i.e. an ergodicity result, as for example in \[KS, F1, FM, Ma, BKL, DZ2\]).
Acknowledgements. This work was supported by the Fonds National Suisse. I am grateful to Sergei Kuksin, Armen Shirikyan, and Martin Hairer for their comments and suggestions.
## 2 Model and Definitions
We consider equations of the form
$$\begin{array}{cc}& \mathrm{d}u=\left((1+\mathrm{i}\alpha )\mathrm{\Delta }u+u(1+\mathrm{i}\beta )|u|^{2q}u\right)\mathrm{d}t+\xi \mathrm{d}w(t),\hfill \\ & u(x,t)𝐂,x𝐑^d,t\mathrm{\hspace{0.17em}0},\alpha ,\beta 𝐑,\hfill \end{array}$$
(2.1)
where $`w(t)`$ is a Wiener process and Eq.(2.1) is understood as an Itô stochastic differential in $`t`$. For a while we simply assume $`\xi (,t)𝒞_\mathrm{b}^{\mathrm{}}(𝐑^d)`$ uniformly in $`t`$ and it is adapted to the Wiener process in $`t`$. A specific example will be considered in Section 5. We also assume $`u(,0)=u_0𝒞_\mathrm{b}^{\mathrm{}}(𝐑^d)`$. We make the following
###### Hypothesis 2.1
. We assume $`d2`$, $`q>\frac{1}{2}`$, and
$$(1+\alpha \beta )<|\alpha \beta |\frac{\sqrt{2q+1}}{q},|\beta |\frac{\sqrt{2q+1}}{q}.$$
(2.2)
Remark. The second inequality in (2.2) implies the first one. We wrote the first one because it appears in this form in the proof of Proposition 4.4, while the second condition appears in the proof of Lemma 10.3, see Ginibre and Velo \[GV1, GV2\] for the most general results available in this direction.
We next introduce the function spaces used in this paper. Let
$$\begin{array}{c}\hfill \phi _{\delta ,y}(x)=\mathrm{exp}\left(\sqrt{1+\delta ^2|xy|^2}\right).\end{array}$$
(2.3)
The main feature of this function is that it belongs to $`\mathrm{L}_{}^p`$ for all $`p`$ and
$$\begin{array}{c}\hfill \frac{^n\phi _{\delta ,y}}{\phi _{\delta ,y}}_{\mathrm{}}=A_n\delta ^n<\mathrm{}\end{array}$$
(2.4)
for all $`y𝐑^d`$ and $`n𝐍`$. This function is used as weight on Sobolev and Lebesgue spaces:
###### Definition 2.2
. The local Lebesgue space $`\mathrm{L}_{\delta ,y}^2`$ is defined as the completion of $`𝒞_\mathrm{b}^{\mathrm{}}`$ (bounded smooth functions) in the norm induced by the scalar product
$$(f,g)_{\delta ,y}=\phi _{\delta ,y}(x)\overline{f(x)}g(x)dx.$$
The local Sobolev spaces $`\mathrm{H}_{\delta ,y}^m`$ are defined as
$$\mathrm{H}_{\delta ,y}^m=\{f:^kf\mathrm{L}_{\delta ,y}^2,k=0,\mathrm{},m\}.$$
The uniformly local Sobolev spaces $`\mathrm{H}_{\mathrm{ul}}^m`$ are defined as the completion of $`𝒞_\mathrm{b}^{\mathrm{}}`$ in the norm
$$f_{\mathrm{H}_{\mathrm{ul}}^m}^2=\underset{k=0}{\overset{m}{}}\underset{y𝐑^d}{sup}(^kf,^kf)_{\delta ,y}.$$
Remark that $`\mathrm{H}_{\mathrm{ul}}^m`$ is actually independent of $`\delta >0`$, since the following inclusion holds:
$$\begin{array}{c}\hfill \mathrm{H}_{\mathrm{ul}}^m=\underset{y𝐑^d}{}\mathrm{H}_{\delta ,y}^m\underset{\delta >0}{}\mathrm{H}_{\delta ,y}^m.\end{array}$$
(2.5)
Usual Sobolev embeddings hold \[Ad\], for example if $`m>d/2`$, the inequality
$$\begin{array}{c}\hfill f_{\mathrm{}}^2C\delta f_{\mathrm{H}_{\mathrm{ul}}^m}^2\end{array}$$
(2.6)
implies the continuous embedding $`\mathrm{H}_{\mathrm{ul}}^m\mathrm{L}_{}^{\mathrm{}}`$. Moreover, by the Rellich–Kondrachov Theorem \[Ad\],
$$\begin{array}{c}\hfill \mathrm{H}_{\delta ,y}^{m+k}\mathrm{H}_{\delta ^{},y}^m\end{array}$$
(2.7)
is compact if $`k>d/2`$ and $`0<\delta <\delta ^{}`$ (see Section 11).
Notations. Throughout the paper, $`\overline{z}`$ denotes the complex conjugate of $`z`$, $`f_t(x)f(x,t)`$ hence $`f_t_𝒳`$ is the norm of $`f(x,t)`$ in the space $`𝒳(\mathrm{d}x)`$ (e.g. $`𝒳=\mathrm{L}_{}^2`$). Norms in Lebesgue spaces $`\mathrm{L}_{}^p`$ are denoted $`_p`$ ($``$ is usually the norm on $`\mathrm{L}_{\delta ,y}^2`$ for the current choice of $`\delta `$ and $`y`$). Expectations and probabilities with respect to the Wiener measure are denoted $`𝐄`$ and $`𝐏`$. We denote the integer part of the positive real $`x`$ by $`[x]\mathrm{max}\{n𝐍:nx\}`$. Symbols $`C,C_1,C_2,\mathrm{},c,c_1,c_2,\mathrm{}`$ usually denote generic numerical constant. The product $`fg`$ means the convolution of the functions $`f`$ and $`g`$. The cube of side $`L`$ and centre $`0`$ in $`𝐑^d`$ is $`Q_L=[\frac{1}{2}L,\frac{1}{2}L]^d`$.
## 3 Summary of Results
In this section, we describe in a rather informal way the main results of this paper. The first result (Section 4) is the following theorem of existence of smooth bounded solutions to Eq.(2.1): Theorem A. If Hypothesis 2.1 holds, then Eq.(2.1) with initial data $`u_0𝒞_\mathrm{b}^{\mathrm{}}`$ has a unique solution $`u(x,t)=u_t(x)`$. For all real $`p1`$ and integer $`m`$, there is a $`B_{p,m}<\mathrm{}`$ such that for all $`t>0`$:
$$𝐄u_t_{\mathrm{H}_{\mathrm{ul}}^m}^pB_{p,m}.$$
The proof relies on well-known estimates \[GV1, Mi, C1\] using the dissipative nature of the nonlinear term in Eq.(2.1) for the deterministic part and on Itô’s Lemma to treat the stochastic term. Actually, in the evolution equation for $`𝐄u_t_{\mathrm{H}_{\mathrm{ul}}^m}^p`$, Itô’s Lemma only generates terms which are dominated by the nonlinear dissipative term and this implies that the techniques which were developed for the deterministic equation are applicable.
By Lemma 10.3 we can extend the existence result to the space $`\mathrm{L}_{\delta ,y}^2`$, hence we can define the Markovian Feller semi-group $`𝒫_t`$ acting on $`𝒞_\mathrm{b}(\mathrm{L}_{\delta ,y}^2,𝐂)`$ for any specific choice of $`y`$, for example $`y=0`$:
$$\begin{array}{c}\hfill \left(𝒫_tf\right)(u)=_{\mathrm{L}_{\delta ,0}^2}f(\eta )𝐏\left(u_t\mathrm{d}\eta \right).\end{array}$$
An invariant measure for Eq.(2.1) is a fixed point of the dual semi-group $`𝒫_t^{}`$. We next assume that $`\xi `$ is an homogeneous process adapted to the Brownian motion. Since $`\mathrm{H}_{\mathrm{ul}}^m`$ is compactly embedded into $`\mathrm{L}_{\delta ,y}^2`$ for $`m>d/2`$ (see (2.7)) the following theorem (see Section 5) is an immediate consequence of Theorem A by the Prokhorov and Krylov–Bogolyubov Theorems (see \[Ar, VF, DZ2\]):
Theorem B. There exists at least one invariant measure $`\mu `$ for Eq.(2.1). This measure is homogeneous in $`x`$ and its support is contained in $`_{m0}\mathrm{H}_{\mathrm{ul}}^m`$.
Finally, in Section 6, we define the random attractor (see \[CDF\])
$$\begin{array}{cc}\hfill 𝒜_\omega =& \overline{\underset{R>0}{}𝒜(\omega ,R)}^{\mathrm{H}_{\mathrm{ul}}^m},\hfill \\ \hfill 𝒜(\omega ,R)=& \underset{T>0}{}\overline{\underset{t>T}{}\mathrm{\Phi }_{\theta ^t\omega }^t(B_R)}^{\mathrm{H}_{\mathrm{ul}}^m}.\hfill \end{array}$$
Here and below $`\mathrm{\Phi }_\omega ^t`$ is the semi-group generated by Eq.(2.1), $`\theta ^t`$ is the time-shift of the noise, and $`T_x`$ the group of spatial translations. Moreover $`B_R\mathrm{H}_{\mathrm{ul}}^m`$ is the ball of radius $`R`$ and centre $`0`$. We introduce the following dynamical observables (see \[KH, LQ\]):
$$\begin{array}{cc}\hfill h_{\mathrm{top}}=& \underset{\epsilon 0}{lim}\underset{L\mathrm{}}{lim}\frac{1}{L^d}\underset{n\mathrm{}}{lim}\frac{1}{n\tau }\mathrm{log}N_{\omega ,n,\tau ,Q_L,\epsilon }𝐏(\mathrm{d}\omega ),\hfill \\ \hfill h_\mu =& \underset{\epsilon 0}{lim}\underset{L\mathrm{}}{lim}\frac{1}{L^d}\underset{n\mathrm{}}{lim}\frac{1}{n\tau }H_\mu \left(\underset{x𝐙^dQ_L}{}\underset{k=0}{\overset{n1}{}}\mathrm{\Phi }_\omega ^{k\tau }T_x(\mathrm{\Sigma }_{\theta ^{k\tau }T_x\omega ,\epsilon })\right)𝐏(\mathrm{d}\omega ),\hfill \\ \hfill _\epsilon =& \underset{L\mathrm{}}{lim}\frac{\mathrm{log}M_{\epsilon ,Q_L,\omega }}{L^d}𝐏(\mathrm{d}\omega ),\hfill \\ \hfill d_{\mathrm{up}}=& \underset{\epsilon 0}{lim\; sup}\frac{_\epsilon }{\mathrm{log}\epsilon ^1},\hfill \end{array}$$
(3.1)
where $`N_{\omega ,n,\tau ,Q,\epsilon }`$ is the cardinality of a minimal $`(n,\epsilon )`$–cover of $`𝒜_\omega |_Q`$, $`\mathrm{\Sigma }_{\omega ,\epsilon }`$ is a sequence of partitions of $`𝒜_\omega `$ in sets of diameter at most $`\epsilon `$ in the metric of $`\mathrm{L}_{}^{\mathrm{}}(Q_1)`$, $`M_{\epsilon ,Q,\omega }`$ is the least cardinality of an $`\epsilon `$–cover of $`𝒜_\omega |_Q`$ and $`Q_L=[\frac{1}{2}L,\frac{1}{2}L]^d`$ (see Section 6 for detailed definitions).
The quantities in Eq.(3.1) are called respectively the topological entropy, the metric or measure-theoretic entropy \[KH\], the Kolmogorov–Tikhomirov $`\epsilon `$–entropy \[KT\] and the upper (box-counting) dimension. It is important to note that the above numbers are all spatial densities (limit as $`L\mathrm{}`$ of quantities divided by $`L^d`$) although the limits are not taken in the most natural order. They are thus spatially localised versions of the usual entropies and dimensions.
We then prove the following estimates:
Theorem C. There is a $`\gamma <\mathrm{}`$ such that $`h_\mu h_{\mathrm{top}}\gamma d_{\mathrm{up}}<\mathrm{}`$.
The proof that all the various limits in Eq.(3.1) exist relies on standard subadditive bounds \[KH\]. The upper bound on $`d_{\mathrm{up}}`$ follows from spatially localised estimates of the rate of divergence of nearby orbits (Lemma 6.7) as well as the smoothing action of the evolution (see Section 7, in particular Lemma 7.1). It is similar to the proof of the deterministic case \[CE2\].
## 4 Bounded Smooth Solutions
Our first result in this paper is the existence (and uniqueness by Lemma 10.1) of smooth bounded solutions to Eq.(2.1).
###### Theorem 4.1
. If Hypothesis 2.1 holds, then Eq.(2.1) with initial data $`u_0𝒞_\mathrm{b}^{\mathrm{}}`$ has a unique solution $`u(x,t)=u_t(x)`$. For all real $`p1`$ and integer $`m`$, there is a $`B_{p,m}<\mathrm{}`$ such that for all $`t>0`$:
$$𝐄u_t_{\mathrm{H}_{\mathrm{ul}}^m}^pB_{p,m}.$$
Remark. This proof is amply simplified by our assumptions on the regularity of $`\xi _t`$ in Eq.(2.1). A much more general theory of stochastic PDEs on $`𝐑^d`$ can be found, for example, in Krylov \[Kr\]. Funaki \[F1, F2\] has studied a similar equations with stronger assumptions on the nonlinearity and Eckmann–Hairer \[EH\] have recently proved a similar result for stochastic forcings with finite energy.
Proof. In the first part of the proof, we fix $`y𝐑^d`$ and $`\delta >0`$ such that $`A_1\delta +A_2\delta ^2<1`$ (see Eq.(2.4)). We write $``$ and $`(,)`$ for the norm and scalar product in the corresponding space $`\mathrm{L}_{\delta ,y}^2`$. All bounds will actually turn out to be uniform in $`y`$. We stress that scalar products denoted $`(,)`$ contain the weight $`\phi _{\delta ,y}`$ (see Definition 2.2) hence integration by parts produces commutators of the form $`\phi _{\delta ,y}/\phi _{\delta ,y}`$. From now on, we also write $`\phi `$ for $`\phi _{\delta ,y}`$.
Let $`=(1+\mathrm{i}\alpha )\mathrm{\Delta }+1`$. For $`fD_m(\mathrm{\Delta })\mathrm{H}_{\delta ,y}^m`$ (the domain of the closure in $`\mathrm{H}_{\delta ,y}^m`$ of $`\mathrm{\Delta }`$ with core $`𝒞_\mathrm{b}^{\mathrm{}}`$), the following holds by Eq.(2.4):
$$\begin{array}{cc}& \mathrm{Re}(^mf,^mf)\hfill \\ & =(^{m+1}f,^{m+1}f)+(^mf,^mf)+\mathrm{Re}(\phi \phi ^1^mf,(1+\mathrm{i}\alpha )^{m+1}f)\hfill \\ & \frac{1}{2}(^{m+1}f,^{m+1}f)+\left(1+\frac{1}{2}(1+\alpha ^2)\right)(^mf,^mf).\hfill \end{array}$$
Namely $`(1+(1+\alpha ^2)/2)`$ is a dissipative operator hence by the Lumer–Phillips Theorem \[Y1\], $``$ generates a strongly continuous quasi-bounded semi-group $`\mathrm{exp}(t)`$ on $`\mathrm{H}_{\delta ,y}^m`$, with
$$\begin{array}{c}\hfill e^t_{\mathrm{H}_{\delta ,y}^m\mathrm{H}_{\delta ,y}^m}e^{ct}\end{array}$$
(4.1)
for some $`c<\mathrm{}`$. Remark that we may have chosen $`\delta `$ such that $`c=1+\epsilon `$ for arbitrarily small $`\epsilon >0`$.
We define mild solutions to Eq.(2.1) in $`\mathrm{H}_{\delta ,y}^m`$ by the Duhamel formula (and the Itô integral):
$$\begin{array}{c}\hfill z_t=e^tz_0(1+\mathrm{i}\beta )_0^te^{(ts)}|z_s|^{2q}z_sds+_0^te^{(ts)}\xi _sdw_s.\end{array}$$
(4.2)
We let $`P_M:𝐑^+𝐑^+`$ be a smooth cutoff function satisfying $`P_M(x)=1`$ if $`x<M`$ and $`P_M(x)=0`$ if $`x>M+1`$. We introduce this cutoff into the nonlinear term above, effectively rendering the nonlinearity uniformly Lipschitz:
$$\begin{array}{c}\hfill \stackrel{~}{z}_t=e^t\stackrel{~}{z}_0(1+\mathrm{i}\beta )_0^te^{(ts)}P_M(|\stackrel{~}{z}_s|)|\stackrel{~}{z}_s|^{2q}\stackrel{~}{z}_sds+_0^te^{(ts)}\xi _sdw_s.\end{array}$$
(4.3)
We next define the random stopping time $`\tau (R)`$ by
$$\begin{array}{c}\hfill \tau (R)=\mathrm{min}\left\{t>\mathrm{\hspace{0.17em}0}:\stackrel{~}{z}_t_{\mathrm{}}R\right\}.\end{array}$$
(4.4)
We fix arbitrarily a positive number $`R<M`$, and if $`\chi _I`$ denotes the characteristic function of the set $`I`$, we consider the following integral equation for $`t<\tau (R)`$:
$$\begin{array}{c}\hfill u_t=e^tu_0(1+\mathrm{i}\beta )_0^te^{(ts)}P_M(|u_s|)|u_s|^{2q}u_s\chi _{\{s\tau (R)\}}ds+_0^te^{(ts)}\xi _s\chi _{\{s\tau (R)\}}dw_s.\end{array}$$
(4.5)
The following is a simple consequence of our construction:
###### Lemma 4.2
. There is almost surely a unique function $`u_t`$ satisfying Eq.(4.5), this function is independent of $`M>R`$ and it also satisfies Eq.(4.2) for $`t<\tau (R)`$.
Proof. See \[DZ1, Ku\] for the properties of the stochastic convolution and Section 10 for the contraction argument needed to prove uniqueness.
The remaining part of the proof of Theorem 4.1 follows very closely the paper by Mielke \[Mi\] which is based on \[BGO, C1, GV1, GV2\]. We first establish uniform bounds in $`\mathrm{L}_{\delta ,y}^2`$.
###### Lemma 4.3
. For all $`\delta >0`$ and $`p1`$, there are $`C_{0,p}(\delta )`$ such that the following bound holds for all $`t>0`$ and all $`y𝐑^d`$:
$$\begin{array}{c}\hfill 𝐄u_t_{\mathrm{L}_{\delta ,y}^2}^pC_{0,p}(\delta ).\end{array}$$
(4.6)
Proof. We first estimate the square of the norm in $`\mathrm{L}_{\delta ,y}^2=\mathrm{L}_{}^2`$. By Itô’s formula, we have
$$\begin{array}{cc}\hfill \mathrm{d}u_t^2=& 2u_t^2\mathrm{d}t2\mathrm{R}\mathrm{e}(\phi \phi ^1u_t,(1+\mathrm{i}\alpha )u_t)\mathrm{d}t+2u_t^2\mathrm{d}t\hfill \\ & 2\mathrm{R}\mathrm{e}(u_t,(1+\mathrm{i}\beta )|u_t|^{2q}u_t)\mathrm{d}t+\xi _t^2\mathrm{d}t+2\mathrm{R}\mathrm{e}(u_t,\xi _t)\mathrm{d}w_t\hfill \\ \hfill & u_t^2\mathrm{d}t+\left(2+(1+\alpha ^2)\right)u_t^2\mathrm{d}t2(u_t,|u_t|^{2q}u_t)\mathrm{d}t\hfill \\ & +\xi _t^2\mathrm{d}t+2\mathrm{R}\mathrm{e}(u_t,\xi _t)\mathrm{d}w_t\hfill \\ \hfill & u_t^2\mathrm{d}t+C(\alpha ,\xi _t,q)\mathrm{d}tu_t^2\mathrm{d}t+2\mathrm{R}\mathrm{e}(u_t,\xi _t)\mathrm{d}w_t.\hfill \end{array}$$
(4.7)
We integrate this last inequality over $`t`$ and take expectations. By standard arguments the expectation of the Itô integral vanishes (recall that we consider stopped solutions, Eq.(4.5), see \[DZ1\]) and we obtain
$$𝐄u_T^2𝐄u_0^2𝐄_0^Tu_t^2dt𝐄_0^T\left(u_t^2C\right)dt.$$
By Gronwall’s inequality, this is
$$𝐄u_T^2\mathrm{max}\{C_{0,2},(𝐄u_0^2C_{0,2})e^T+C_{0,2}\}.$$
For higher powers of the $`\mathrm{L}_{}^2`$ norm, we use Itô’s formula again:
$$\frac{1}{p}\mathrm{d}u_t_2^{2p}=u_t^{2p2}\mathrm{d}u_t^2+2(p1)u_t^{2p4}\left(\mathrm{Re}(u_t,\xi _t)\right)^2\mathrm{d}t,$$
hence (after substituting the estimate (4.7))
$$𝐄u_T_2^{2p}𝐄u_0_2^{2p}𝐄_0^T\left(u_t^{2p}C_{0,2p}\right)dt,$$
which by Gronwall’s inequality gives a uniform bound on $`u_t^p`$ for $`p>2`$. For $`p[1,2)`$, we use Jensen’s inequality:
$$𝐄u_t^p\left(𝐄u_t^2\right)^{p/2}C_{0,2}^{p/2}=C_{0,p}.$$
If $`u_0`$ is uniformly bounded and because $`\xi _t_{\mathrm{L}_{\delta ,y}^2}^p`$ is bounded uniformly in $`y`$ and $`t`$, we obtain the uniform bound in the spaces $`\mathrm{L}_{\delta ,y}^2`$ for all $`y`$
$$\begin{array}{c}\hfill \underset{t>0}{sup}\underset{y𝐑^d}{sup}𝐄u_t_{\mathrm{L}_{\delta ,y}^2}^pC_{0,p}(\delta ),\end{array}$$
(4.8)
which proves Lemma 4.3.
###### Proposition 4.4
. For all $`\delta >0`$ and $`p1`$, there are $`C_{1,p}(\delta )`$ such that the following bound holds for all $`t>0`$ and all $`y𝐑^d`$:
$$𝐄u_t_{\mathrm{H}_{\delta ,y}^1}^pC_{1,p}(\delta ).$$
Proof. We first consider the differential
$$\begin{array}{cc}\hfill \mathrm{d}u_t^2=& 2\mathrm{\Delta }u_t^2\mathrm{d}t2\mathrm{R}\mathrm{e}(\phi \phi ^1u_t,(1+\mathrm{i}\alpha )\mathrm{\Delta }u_t)\mathrm{d}t+2u_t^2\mathrm{d}t\hfill \\ & +2\mathrm{R}\mathrm{e}(\mathrm{\Delta }u_t,(1+\mathrm{i}\beta )|u_t|^{2q}u_t)\mathrm{d}t+2\mathrm{R}\mathrm{e}(\phi \phi ^1u_t,(1+\mathrm{i}\beta )|u_t|^{2q}u_t)\mathrm{d}t\hfill \\ & +\xi _t^2\mathrm{d}t+2\mathrm{R}\mathrm{e}(u_t,\xi _t)\mathrm{d}w_t\hfill \\ \hfill & 2\mathrm{\Delta }u_t^2\mathrm{d}t+2u_t^2\mathrm{d}t+2\mathrm{R}\mathrm{e}(\mathrm{\Delta }u_t,(1+\mathrm{i}\beta )|u_t|^{2q}u_t)\mathrm{d}t\hfill \\ & +2\left(\sqrt{1+\alpha ^2}\mathrm{\Delta }u_t+\sqrt{1+\beta ^2}|u_t|^{2q+1}\right)u_t\mathrm{d}t\hfill \\ & +\xi _t^2\mathrm{d}t2\mathrm{R}\mathrm{e}(u_t,\phi ^1(\phi \xi _t))\mathrm{d}w_t,\hfill \end{array}$$
(4.9)
and we also compute the following differential that will help us to cancel out some of the terms above:
$$\begin{array}{cc}\hfill \frac{1}{q+1}\mathrm{d}|u_t|^{q+1}^2=& \mathrm{\hspace{0.17em}2}\mathrm{Re}(|u_t|^{2q}u,(1+\mathrm{i}\alpha )\mathrm{\Delta }u_t)\mathrm{d}t+2|u_t|^{q+1}^2\mathrm{d}t\hfill \\ & 2|u_t|^{2q}u_t^2\mathrm{d}t+2q\left(\mathrm{Re}(|u_t|^{q1}u_t,\xi _t)\right)^2\mathrm{d}t+2\mathrm{R}\mathrm{e}(|u_t|^{2q}u_t,\xi _t)\mathrm{d}w_t\hfill \\ \hfill & \mathrm{\hspace{0.17em}2}\mathrm{Re}(|u_t|^{2q}u,(1+\mathrm{i}\alpha )\mathrm{\Delta }u_t)\mathrm{d}t+2|u_t|^{q+1}^2\mathrm{d}t+2q\xi _t^2|u_t|^q^2\mathrm{d}t\hfill \\ & 2|u_t|^{2q}u_t^2\mathrm{d}t+2\mathrm{R}\mathrm{e}(|u_t|^{2q}u_t,\xi _t)\mathrm{d}w_t.\hfill \end{array}$$
(4.10)
We take a convex combination of Inequalities (4.10) and (4.9) (here $`\lambda [0,1]`$):
$$\begin{array}{cc}& \frac{1}{2}\left(\lambda \mathrm{d}u_t^2+\frac{(1\lambda )}{q+1}\mathrm{d}|u_t|^{q+1}^2\right)\hfill \\ & \left(\lambda u_t^2+(1\lambda )|u_t|^{q+1}^2\right)\mathrm{d}t+\left((1\lambda )q\xi _t^2|u_t|^q^2+\frac{\lambda }{2}\xi _t^2\right)\mathrm{d}t\hfill \\ & +\lambda \left(\sqrt{1+\alpha ^2}\mathrm{\Delta }u_t+\sqrt{1+\beta ^2}|u_t|^{2q+1}\right)u_t\mathrm{d}t\hfill \\ & +\mathrm{d}t+\mathrm{Re}\left(\lambda (u_t,\phi ^1(\phi \xi _t))+(1\lambda )b(|u_t|^{2q}u_t,\xi _t)\right)\mathrm{d}w_t.\hfill \end{array}$$
The term denoted by $``$ is treated separately:
$$\begin{array}{cc}\hfill =& \left(\lambda \mathrm{\Delta }u_t^2+(1\lambda )|u_t|^{2q+1}^2\right)+\mathrm{Re}\left(1\mathrm{i}(\lambda \beta (1\lambda )\alpha )\right)(|u_t|^{2q}u_t,\mathrm{\Delta }u_t)\hfill \\ \hfill & \epsilon \left(\lambda \mathrm{\Delta }u_t^2+\frac{(1\lambda )}{q+1}|u_t|^{2q+1}^2\right)2(1\epsilon )\sqrt{\lambda (1\lambda )}\left|(|u_t|^{2q}u_t,\mathrm{\Delta }u_t)\right|\hfill \\ & +\mathrm{Re}\left(1\mathrm{i}(\lambda \beta (1\lambda )\alpha )\right)(|u_t|^{2q}u_t,\mathrm{\Delta }u_t)\hfill \\ \hfill & \epsilon \left(\lambda \mathrm{\Delta }u_t^2+\frac{(1\lambda )}{q+1}|u_t|^{2q+1}^2\right)+\stackrel{~}{}(\epsilon ).\hfill \end{array}$$
Under Hypothesis 2.1, there is an $`\epsilon >0`$ such that $`\stackrel{~}{}(\epsilon )`$ is negative (see e.g. \[GV1, Mi, BGO\]). The proof goes as follows: we first remark that integration by parts leads to
$$(|u_t|^{2q}u_t,\mathrm{\Delta }u_t)=(\frac{\phi }{\phi }|u_t|^{2q}u_t,u_t)(q+1)\phi |u_t|^{2q}|u_t|^2\left(1+\frac{q}{1+q}\frac{\overline{u_t}^2}{|u_t|^2}\frac{u_t^2}{|u_t|^2}\right).$$
The last bracket above is of the form $`1+z`$. Its argument can be estimated as follows: $`|\mathrm{arg}(1+z)|\mathrm{arcsin}|z|=\mathrm{arcsin}\frac{q}{1+q}\theta `$. We plug this into $`\stackrel{~}{}`$:
$$\begin{array}{cc}\hfill \stackrel{~}{}(\epsilon )& (q+1)\left|(|u_t|^{2q},|u_t|^2)\right|\left\{\left(2(1\epsilon )\sqrt{\lambda (1\lambda )}+\mathrm{cos}\theta \right)\right|\lambda \beta (1\lambda )\alpha |\mathrm{sin}\theta \}\hfill \\ & +C(\alpha ,\beta ,\lambda ,\epsilon )|u_t|^{2q+1}u_t\hfill \end{array}$$
The curly bracket above can be made positive by suitably choosing $`\lambda `$ and $`\epsilon `$, namely we take $`\lambda =\mathrm{cos}^2\eta `$, we optimise for $`\eta `$ and we obtain the following condition, which is obviously fulfilled for small $`\epsilon >0`$ under Hypothesis 2.1 (remark that $`1/\mathrm{tan}\theta =\sqrt{2q+1}/q`$):
$$(1+\alpha \beta )|\beta \alpha |/\mathrm{tan}\theta +\epsilon (2\epsilon )/\mathrm{sin}^2\theta \mathrm{\hspace{0.17em}0},$$
(see Ginibre and Velo \[GV1\] for this argument, Mielke \[Mi\] has a slightly different formulation). We thus obtain
$$\begin{array}{cc}& \frac{1}{2}\left(\lambda \mathrm{d}u_t^2+\frac{(1\lambda )}{q+1}\mathrm{d}|u_t|^{q+1}^2\right)\hfill \\ & \epsilon \left(\lambda \mathrm{\Delta }u_t^2+\frac{(1\lambda )}{q+1}|u_t|^{2q+1}^2\right)\mathrm{d}t\hfill \\ & +\left(C_1u_t^2+C_2|u_t|^{q+1}^2\right)\mathrm{d}t+\left(C_3\xi _t^2|u_t|^q^2+C_4\xi _t^2\right)\mathrm{d}t\hfill \\ & +\left(C_5\mathrm{\Delta }u_t+C_6|u_t|^{2q+1}+C_7\right)u_t\mathrm{d}t\hfill \\ & +\mathrm{Re}\left(\lambda (u_t,\phi ^1(\phi \xi _t))+(1\lambda )(|u_t|^{2q}u_t,\xi _t)\right)\mathrm{d}w_t\hfill \\ & C\mathrm{d}t\frac{\epsilon }{2}\left(\lambda u_t^2+\frac{(1\lambda )}{q+1}|u_t|^{q+1}^2\right)\mathrm{d}t\hfill \\ & +\mathrm{Re}\left(\lambda (u_t,\phi ^1(\phi \xi _t))+(1\lambda )(|u_t|^{2q}u_t,\xi _t)\right)\mathrm{d}w_t,\hfill \end{array}$$
thanks to the following obvious inequality:
$$\begin{array}{c}\hfill \mathrm{\Delta }u_t^2\rho u_t^2+C\rho ^2u_t^2\end{array}$$
(4.11)
which holds for all $`\rho >0`$ and for some $`C>0`$. As before, we take expectations, integrate over $`t`$ and we use Gronwall’s inequality to find out the following bound:
$$\begin{array}{cc}& \mathrm{max}\{𝐄u_T^2,𝐄|u_T|^{q+1}^2\}\hfill \\ & \mathrm{max}\{C_{1,2},(𝐄u_0^2+𝐄|u_0|^{q+1}^2C_{1,2})e^{\epsilon T}+C_{1,2}\}.\hfill \end{array}$$
(4.12)
This and Lemma 4.3 prove Proposition 4.4.
We next consider solutions $`z(x,t)`$ to Eq.(4.2) with bounded initial condition. Proposition 4.4 on $`u(x,t)`$ implies the following:
###### Proposition 4.5
. For all $`p1`$, there is a $`C_{\mathrm{},p}`$ such that for all $`t>0`$, the following holds
$$\begin{array}{c}\hfill 𝐄z_t_{\mathrm{}}^pC_{\mathrm{},p}.\end{array}$$
(4.13)
Proof. By the bound (2.6), if $`d=1`$ then Proposition 4.4 implies the bound (4.13) for stopped solutions. If $`d=2`$ we need a bound in $`\mathrm{H}_{\delta ,y}^2`$. This is easily achieved with the help of a Gagliardo–Nirenberg inequality which we prove in Section 12:
###### Lemma 4.6
. Let $`f\mathrm{H}_{\mathrm{ul}}^3(𝐑^2)`$. For all $`K>0`$ there are $`C(K)`$, $`\eta `$ such that
$$\begin{array}{cc}& \phi _{\delta ,y}(x)\left|\mathrm{\Delta }\left(|f(x)|^{2q}f(x)\right)\overline{\mathrm{\Delta }f(x)}\right|dx\hfill \\ & \frac{1}{K}\phi _{\delta ,y}|^3f(x)|^2dx+C(K)\left(\underset{y}{sup}\phi _{\delta ,y}(x)|f(x)|^{2(q+1)}dx\right)^\eta .\hfill \end{array}$$
We use Inequality (4.11) (with $`u_t`$ replaced by $`u_t`$), Lemma 4.6, and the estimate (4.12) to bound the time derivative of $`\mathrm{\Delta }u_t_{\mathrm{L}_{\delta ,x}^2}^2`$:
$$\begin{array}{cc}\hfill \frac{1}{2}\mathrm{d}\mathrm{\Delta }u_t_{\mathrm{L}_{\delta ,x}^2}^2& \frac{1}{2}^3u_t_{\mathrm{L}_{\delta ,x}^2}^2\mathrm{d}t+(1+(1+\alpha ^2))\mathrm{\Delta }u_t_{\mathrm{L}_{\delta ,x}^2}^2\mathrm{d}t\hfill \\ & +C(1+\beta ^2)\underset{y}{sup}|u_t|^{q+1}_{\mathrm{L}_{\delta ,y}^2}^{2\eta }\mathrm{d}t+\frac{1}{2}\mathrm{\Delta }\xi _t_{\mathrm{L}_{\delta ,x}^2}^2\mathrm{d}t+(\mathrm{\Delta }u_t,\mathrm{\Delta }\xi _t)_{\delta ,x}\mathrm{d}w_t\hfill \\ \hfill & \left(\frac{1}{2}\mathrm{d}\mathrm{\Delta }u_t_{\mathrm{L}_{\delta ,x}^2}^2+C\right)\mathrm{d}t+(u_t,\phi ^1\mathrm{\Delta }(\phi \mathrm{\Delta }\xi _t))_{\delta ,x}\mathrm{d}w_t,\hfill \end{array}$$
where $`C`$ depends on the parameters in Eq.(2.1) (including $`\xi _t`$) and on $`u_t_{\mathrm{L}_{\delta ,x}^2}^p`$, $`|u_t|^{q+1}_{\mathrm{L}_{\delta ,x}^2}^\eta `$, and $`u_t_{\mathrm{L}_{\delta ,x}^2}^2`$ which satisfy Bounds (4.6) and (4.12) (or rather some extension of it to deal with the power $`\eta `$). By the usual Gronwall inequality, this proves (4.13) for stopped solutions (Eq.(4.5)).
We now choose a very large $`n_0C_{0,2}+C_{1,2}+C_{2,2}`$ and we let
$$_n=\{u:t<\tau (2n)\text{ s.t. }u_t_{\mathrm{}}>n_0+n\},$$
where $`\tau `$ is the stopping time from Eq.(4.4). By Tchebychev’s inequality, by Proposition 4.4 and (2.6) we have
$$\underset{n=1}{\overset{\mathrm{}}{}}𝐏(_n)\underset{n=1}{\overset{\mathrm{}}{}}n^2\underset{t\tau (2n)}{sup}𝐄(u_t_{\mathrm{}}^2)<\mathrm{}.$$
By the Borel–Cantelli Lemma, it means that almost surely only finitely many of the events $`_n`$ happen, and hence $`u_t_{\mathrm{}}`$ remains bounded as the cutoffs $`R`$ and $`M`$ in Eq.(4.5) are sent to infinity. Since $`\tau (R)>0`$ a.s. for all $`R>1`$ (by the a priori bound of Lemma 10.2) a uniform bound holds in a small interval of time and this can be iterated indefinitely. This implies uniform boundedness of $`z_t`$ and a similar argument holds for $`z_t_{\mathrm{}}^p`$, $`p>1`$.
Proof of Theorem 4.1By Proposition 4.4, it only remains to show that $`^mu_t^p`$ is bounded for $`m>1`$. Let $`p=2`$. We assume that it is true for $`m1`$ and we consider
$$\begin{array}{cc}& \frac{1}{2}\mathrm{d}^mu_t^2\hfill \\ & \frac{1}{2}^{m+1}u_t^2\mathrm{d}t+\left(1+\frac{1}{2}(1+\alpha ^2)\right)^mu_t^2\mathrm{d}t\hfill \\ & \mathrm{Re}(1+\mathrm{i}\beta )(^mu_t,^m(|u_t|^{2q}u_t))\mathrm{d}t+\frac{1}{2}^m\xi _t^2\mathrm{d}t+\mathrm{Re}(^mu_t,^m\xi _t)\mathrm{d}w_t\hfill \\ & \frac{1}{2}^{m+1}u_t^2\mathrm{d}t+\left(C_1^mu_t^2\mathrm{d}t+C_2u_t_{\mathrm{H}_{\delta ,y}^{m1}}^2+C_3\right)\mathrm{d}t\hfill \\ & +(1)^m\mathrm{Re}(u_t,\phi ^1^m(\phi ^m\xi _t))\mathrm{d}w_t,\hfill \end{array}$$
Using (4.11) (with $`u_t`$ replaced by $`^{m1}u_t`$), Proposition 4.5, and the recursion assumption, this can be bounded by:
$$\frac{1}{2}\mathrm{d}^mu_t^2\frac{1}{2}\left(^mu_t^2+C\right)\mathrm{d}t+(1)^m\mathrm{Re}(u_t,\phi ^1^m(\phi ^m\xi _t))\mathrm{d}w_t.$$
We then take expectations and integrate:
$$𝐄^mu_t^2𝐄^mu_0^2𝐄_0^t\left(^mu_s^2C\right)ds.$$
The case $`p2`$ is similar.
## 5 Invariant Measures
We now turn to the problem of the existence of an invariant measure for the process defined by Eq.(2.1). We construct here an explicit example of a smooth homogeneous random forcing which admits evident generalisations. Let $`\xi (x)`$ be a $`𝒞^{\mathrm{}}`$ almost periodic function on $`𝐑^d`$. Denoting $`T_y\xi (x)=\xi (x+y)`$, the set
$$G=\overline{\{T_y\xi :y𝐑^d\}}^\mathrm{L}_{}^{\mathrm{}}$$
is a compact group which can be endowed with the normalised Haar measure $`h`$. We denote by $`(G,_1,h)`$ the corresponding probability space ($`_1`$ is the sigma-algebra of Borel sets) and by $`\xi _y`$ the corresponding random variable. Let next $`w_\alpha (t)`$ be a standard Brownian motion (vanishing at $`0`$) on the probability space $`(𝒞_0(𝐑,𝐑),_2,𝐖)`$ ($`_2`$ is the sigma-algebra generated by the topology of uniform convergence on compact sets and $`𝐖`$ is the Wiener measure). Let $`z_\alpha (t)=y(t)y(0)`$ be a continuous process on $`(G,_1,h)`$ adapted to the filtration generated by the Brownian motion. We define the stochastic differential $`\xi _{y(0)+z_\alpha (t)}(x)\mathrm{d}w_\alpha (t)`$ on the probability space $`(\mathrm{\Omega },,𝐏)=(𝒞_0(𝐑,𝐑)\times G,_2\times _1,𝐖\times h)`$. Let $`\mathrm{\Xi }_\omega (x,t)`$ denote $`\xi _{y(0)+z_\alpha (t)}(x)w_\alpha (t)`$ where $`\mathrm{\Omega }\omega =(\alpha ,y(0))𝒞_0(𝐑,𝐑)\times G`$. By the nature of Haar measures, $`𝐏`$ is homogeneous in $`x`$, i.e. $`T_y^{}𝐏=𝐏`$ for all $`y𝐑^d`$ (see Vishik and Fursikov \[VF\] for a discussion of homogeneous measures).
Let $`\mathrm{\Phi }_\omega ^t`$ be the semi-flow generated by Eq.(2.1) with noise $`\mathrm{\Xi }_\omega (x,t)`$. Using Lemma 10.3, we can define a Markov semi-group $`𝒫_t`$ acting on $`𝒞_\mathrm{b}(\mathrm{L}_{\delta ,0}^2,𝐂)`$ by
$$\begin{array}{c}\hfill \left(𝒫_tf\right)(u)=_{\mathrm{L}_{\delta ,0}^2}f(\eta )𝐏\left(\mathrm{\Phi }_\omega ^t(u)\mathrm{d}\eta \right).\end{array}$$
(5.1)
$`𝒫_t`$ is a Markovian Feller semi-group (the Feller property follows from the continuity of $`\mathrm{\Phi }_\omega ^t`$). Its dual $`𝒫_t^{}`$ acts on probability measures over $`\mathrm{L}_{\delta ,0}^2`$ by
$$\begin{array}{c}\hfill \left(𝒫_t^{}\mu \right)(B)=_{\mathrm{L}_{\delta ,0}^2}𝐏\left(\mathrm{\Phi }_\omega ^t(u)B\right)\mu (\mathrm{d}u).\end{array}$$
(5.2)
We call $`\mu `$ an invariant measure for Eq.(2.1) if $`𝒫_t^{}\mu =\mu `$ for all $`t>0`$ (see Arnold \[Ar\]).
In this section, we prove the following Theorem, which is actually a simple consequence of the bounds derived in Section 4.
###### Theorem 5.1
. There exists at least one invariant measure $`\mu `$ for Eq.(2.1). This measure is homogeneous in $`x`$ and its support is contained in $`_{m0}\mathrm{H}_{\mathrm{ul}}^m`$.
Proof. We consider the family of measures $`\{\overline{\mu }_t\}_{t>0}`$, where
$$\overline{\mu }_t=\frac{1}{t}_0^t𝒫_s^{}𝜹_0ds,$$
$`𝜹_0`$ being the unit mass at $`0\mathrm{L}_{\delta ,0}^2`$. By Theorem 4.1, this family is tight in $`\mathrm{L}_{\delta ,0}^2`$ for all $`\delta >0`$. Namely for any $`\epsilon >0`$ there is a compact $`K_\epsilon \mathrm{L}_{\delta ,0}^2`$ such that $`\overline{\mu }_t(K_\epsilon )>1\epsilon `$. For the set $`K_\epsilon `$ we choose the ball of radius $`R(\epsilon )`$ in $`\mathrm{H}_{\mathrm{ul}}^m`$ ($`m>d/2`$) for sufficiently large $`R(\epsilon )`$ and the compactness follows from (2.5) and (2.7). By the Prokhorov Theorem, $`\{\overline{\mu }_t\}_{t>0}`$ is weakly precompact and thus there is at least one accumulation point $`\mu `$. By the standard Krylov–Bogolyubov argument $`\mu `$ is an invariant measure (see \[Ar, VF, DZ2\] for a detailed statement of these procedures).
Let $`t_n`$, $`n=1,2,\mathrm{}`$ be a sequence such that $`\mu =\mathrm{w}\mathrm{lim}_n\mathrm{}\overline{\mu }_{t_n}`$. Let $`B_R`$ be the ball of radius $`R`$ in $`\mathrm{H}_{\delta ,y}^m`$ and let $`f_{y,R}`$ be any bounded continuous function on $`\mathrm{L}_{\delta ,y}^2`$ vanishing on $`B_R`$ (which is a compact set). Since the topologies of $`\mathrm{L}_{\delta ,0}^2`$ and of $`\mathrm{L}_{\delta ,y}^2`$ are equivalent, this function is continuous on $`\mathrm{L}_{\delta ,0}^2`$. Obviously $`|f_{y,R}(\eta )\overline{\mu }_{t_n}(\mathrm{d}\eta )|<f_{\mathrm{}}\epsilon (R)`$ for all $`n`$ and $`y`$, where $`\epsilon (R)0`$ as $`R\mathrm{}`$. By weak convergence of $`\overline{\mu }_{t_n}`$ to $`\mu `$ this also holds for $`\mu `$ and hence the support of $`\mu `$ must be contained in $`_{y𝐑}\mathrm{H}_{\delta ,y}^m`$.
We next prove the homogeneity of $`\mu `$. Let $`f𝒞_\mathrm{b}(\mathrm{H}_{\mathrm{ul}}^m,𝐂)`$ and define the translation operator $`T_y`$ by $`T_yf(u)=f(T_yu)`$. We have
$$\begin{array}{cc}\hfill T_yf(\eta )\mu (\mathrm{d}\eta )& =\underset{n\mathrm{}}{lim}\frac{1}{t_n}_0^{t_n}\left(f(T_y\eta )𝐏(\mathrm{\Phi }_\omega ^t(0)\mathrm{d}\eta )\right)dt\hfill \\ & =\underset{n\mathrm{}}{lim}\frac{1}{t_n}_0^{t_n}\left(f(\eta )𝐏(T_y(\mathrm{\Phi }_\omega ^t(0))\mathrm{d}\eta )\right)dt\hfill \\ & =\underset{n\mathrm{}}{lim}\frac{1}{t_n}_0^{t_n}\left(f(\eta )𝐏(\mathrm{\Phi }_{T_y\omega }^t(T_y(0))\mathrm{d}\eta )\right)dt\hfill \\ & =\underset{n\mathrm{}}{lim}\frac{1}{t_n}_0^{t_n}\left(f(\eta )𝐏(\mathrm{\Phi }_\omega ^t(0)\mathrm{d}\eta )\right)dt=f(\eta )\mu (\mathrm{d}\eta ),\hfill \end{array}$$
where we have used the homogeneity of $`𝐏`$. Since the above holds for all $`f`$, it proves that $`\mu `$ is homogeneous and the proof of Theorem 5.1 is finished.
Remark. In the above construction of a tight family of measures, we could have considered any homogeneous initial measure $`\mu _0`$ supported by $`_{m0}\mathrm{H}_{\mathrm{ul}}^m`$ instead of $`𝜹_0`$.
## 6 Entropy Estimates
In this section, we define and estimate different notions of entropy for Eq.(2.1). We start with the topological entropy, then the measure-theoretic entropy and finally the $`\epsilon `$–entropy. All these quantities are extensive, hence we actually define their spatial densities. We define the spatial density of upper box-counting dimension as well.
To do so we first introduce the basic dynamical setup: let $`\mathrm{\Phi }_\omega ^t`$ ($`t>0`$) be the solution semi-flow to Eq.(2.1) for given noise parameter $`\omega `$ and let $`\theta ^t`$ be the shift semi-flow on $`\mathrm{\Omega }`$:
$$\mathrm{\Xi }_{\theta ^\tau \omega }(x,t)=\mathrm{\Xi }_\omega (x,t+\tau )\mathrm{\Xi }_\omega (x,t).$$
Let next $`S^t`$ be the semi-flow on $`\mathrm{L}_{\delta ,0}^2\times \mathrm{\Omega }`$ defined by
$$\begin{array}{cc}\hfill S^t:\mathrm{L}_{\delta ,0}^2\times \mathrm{\Omega }& \mathrm{L}_{\delta ,0}^2\times \mathrm{\Omega }\hfill \\ \hfill (u,\omega )& (\mathrm{\Phi }_\omega ^t(u),\theta ^t(\omega )).\hfill \end{array}$$
We consider the space $`\mathrm{H}_{\mathrm{ul}}^m`$ ($`m>d`$) endowed with the (weaker) topology of uniform convergence on the compact $`Q𝐑^d`$. By standard embeddings (see (2.6) and (2.7)) bounded sets of $`\mathrm{H}_{\mathrm{ul}}^m`$ are compact in $`\mathrm{L}_{}^{\mathrm{}}(Q)`$. Following Crauel et al. \[CDF\], we define the random attractor $`𝒜_\omega `$ as follows: Let $`B_R`$ be the ball of radius $`R`$, centre $`0`$ in $`\mathrm{H}_{\mathrm{ul}}^m`$ and let
$$\begin{array}{cc}\hfill 𝒜_\omega =& \overline{\underset{R>0}{}𝒜(\omega ,R)}^{\mathrm{H}_{\mathrm{ul}}^m},\hfill \\ \hfill 𝒜(\omega ,R)=& \underset{T>0}{}\overline{\underset{t>T}{}\mathrm{\Phi }_{\theta ^t\omega }^t(B_R)}^{\mathrm{H}_{\mathrm{ul}}^m}.\hfill \end{array}$$
By the estimates of Section 4, $`𝒜_\omega `$ is almost surely closed and bounded in $`\mathrm{H}_{\mathrm{ul}}^m`$ hence it is compact in $`\mathrm{L}_{}^{\mathrm{}}(Q)`$ for any bounded $`Q𝐑^d`$. Moreover the diameter of $`𝒜_\omega `$ in $`\mathrm{H}_{\mathrm{ul}}^m`$ is less than some $`R_\omega `$ with $`𝐏(\omega :R_\omega <\mathrm{})=1`$ and $`𝐄(\omega R_\omega )<\mathrm{}`$ (by Theorem 4.1). The following equivariance properties hold (we assume $`\theta ^tT_x=T_x\theta ^t`$ for all $`(x,t)𝐑^d\times 𝐑^+`$):
$$\begin{array}{cc}\hfill \mathrm{\Phi }_\omega ^t𝒜_\omega =& 𝒜_{\theta ^t\omega },\hfill \\ \hfill T_x𝒜_\omega =& 𝒜_{T_x\omega },\hfill \end{array}$$
(6.1)
and it contains the support of any invariant measure for $`S^t`$. Let next $`\mu `$ be an invariant measure in the sense of Section 5, namely a stationary measure for the Markov semi-group (5.1). We also assume that $`𝐏`$ is an invariant measure for $`\theta ^t`$. Then $`\mu \times 𝐏`$ is an invariant measure for the dynamical system $`(S^t,𝒳,)`$ where $`𝒳=\mathrm{L}_{\delta ,0}^2\times \mathrm{\Omega }`$ and $``$ is the associated sigma-algebra. More precisely, one has
$$𝐄\left(\omega \left(\mathrm{\Phi }_\omega ^t\right)^{}\mu \right)=\mu $$
(6.2)
(which is only a rephrasing of $`𝒫_t^{}\mu =\mu `$, see Eq.(5.2)). We introduce the following definitions:
###### Definition 6.1
. Let $`\tau >0`$, $`n𝐍`$, and $`Q𝐑`$. We define a pseudo-metric $`d_{\omega ,n,\tau ,Q}`$ on $`\mathrm{H}_{\mathrm{ul}}^m`$ by
$$d_{\omega ,n,\tau ,Q}(u,v)=\underset{k=0,\mathrm{},n1}{\mathrm{max}}\mathrm{\Phi }_\omega ^{k\tau }(u)\mathrm{\Phi }_\omega ^{k\tau }(v)_{\mathrm{L}_{}^{\mathrm{}}(Q)}.$$
Let $`N_{\omega ,n,\tau ,Q,\epsilon }`$ be the cardinality of a minimal $`(n,\epsilon )`$–cover of $`𝒜_\omega |_Q`$ (that is $`N_{\omega ,n,\tau ,Q,\epsilon }`$ is the least number of open sets whose diameter in the metric $`d_{\omega ,n,\tau ,Q}`$ is at most $`\epsilon `$ and whose union contains $`𝒜_\omega `$).
We define the cube $`Q_L=[\frac{1}{2}L,\frac{1}{2}L]^d`$. We are now able to prove the existence of the spatial density of topological entropy $`h_{\mathrm{top}}`$ for Eq.(2.1):
###### Proposition 6.2
. For all $`\tau >0`$ the following limit exists:
$$h_{\mathrm{top}}=\underset{\epsilon 0}{lim}\underset{L\mathrm{}}{lim}\frac{1}{L^d}\underset{n\mathrm{}}{lim}\frac{1}{n\tau }\mathrm{log}N_{\omega ,n,\tau ,Q_L,\epsilon }𝐏(\mathrm{d}\omega ),$$
(6.3)
This limit is independent of $`\tau >0`$.
Proof. The proof is similar to the deterministic case treated by Collet and Eckmann in \[CE2\] and is reproduced in Section 8.
Let $`𝒰=\{U_1,\mathrm{},U_k,\mathrm{}\}`$ be a countable (or finite) $`\mu `$–measurable partition of $`𝒜_\omega `$. For two partitions $`𝒰`$ and $`𝒱`$, we denote their refinement $`\{U_kV_{\mathrm{}}:U_k𝒰,V_{\mathrm{}}𝒱,\mu (U_kV_{\mathrm{}})>0\}`$ by $`𝒰𝒱`$. Moreover $`\mathrm{\Phi }_\omega ^\tau (𝒰)=\{\mathrm{\Phi }_\omega ^\tau (U_k):U_k𝒰\}`$ is a measurable partition of $`𝒜_{\theta ^\tau \omega }`$ whenever $`𝒰`$ is a measurable partition of $`𝒜_\omega `$. (Here $`\mathrm{\Phi }_\omega ^t`$ stands for the inverse of $`\mathrm{\Phi }_\omega ^t`$, namely $`\mathrm{\Phi }_\omega ^t(x)`$ is the set of all pre-images of $`x`$.)
###### Definition 6.3
. Let $`H_\mu (𝒰)`$ and $`H_\mu (𝒰|𝒱)`$ denote the entropy of a partition and the conditional entropy, both relative to a given measure $`\mu `$. They are defined as follows
$$\begin{array}{cc}\hfill H_\mu (𝒰)& =\underset{U𝒰}{}\mu (U)\mathrm{log}\mu (U),\hfill \\ \hfill H_\mu (𝒰|𝒱)& =\underset{U𝒰,V𝒱}{}\mu (UV)\mathrm{log}\left(\frac{\mu (UV)}{\mu (V)}\right).\hfill \end{array}$$
We adopt here the convention $`0\mathrm{log}0=0`$ therefore $`0<H_\mu (𝒰)\mathrm{log}card(𝒰)`$ (which is possibly infinite for countable $`𝒰`$). We also choose an arbitrary sequence $`\mathrm{\Sigma }_{\omega ,\epsilon }`$ of partitions of $`𝒜_\omega `$ in sets of diameter at most $`\epsilon `$ in the metric of $`\mathrm{L}_{}^{\mathrm{}}(Q_1)`$.
The second result in this section is the existence of the spatial density of measure-theoretic entropy $`h_\mu `$
###### Proposition 6.4
. For all $`\tau >0`$ the following limit exists:
$$h_\mu =\underset{\epsilon 0}{lim}\underset{L\mathrm{}}{lim}\frac{1}{L^d}\underset{n\mathrm{}}{lim}\frac{1}{n\tau }H_\mu \left(\underset{x𝐙^dQ_L}{}\underset{k=0}{\overset{n1}{}}\mathrm{\Phi }_\omega ^{k\tau }T_x(\mathrm{\Sigma }_{\theta ^{k\tau }T_x\omega ,\epsilon })\right)𝐏(\mathrm{d}\omega ).$$
(6.4)
It is independent of $`\tau >0`$ and of the particular choice of the sequence of partitions $`\mathrm{\Sigma }_{\omega ,\epsilon }`$.
Proof. Again, the proof is quite standard, see e.g. \[KH, LQ\] or Section 9.
We next introduce the notions of $`\epsilon `$–entropy $`_\epsilon `$ of Kolmogorov and Tikhomirov \[KT\] and of upper density of dimension $`d_{\mathrm{up}}`$.
###### Definition 6.5
. Let $`M_{\epsilon ,Q,\omega }`$ be the least cardinality of an open cover of $`𝒜_\omega `$ by sets of diameter less than $`\epsilon `$ in the metric of $`\mathrm{L}_{}^{\mathrm{}}(Q)`$ where $`Q`$ is compact (we call this an $`\epsilon `$–cover of $`𝒜_\omega |_Q`$). Let $`_\epsilon `$ be the Kolmogorov–Tikhomirov $`\epsilon `$–entropy defined by
$$_\epsilon =\underset{L\mathrm{}}{lim}\frac{\mathrm{log}M_{\epsilon ,Q_L,\omega }}{L^d}𝐏(\mathrm{d}\omega ),$$
et let $`d_{\mathrm{up}}(\omega )`$ be the upper density of dimension of $`𝒜_\omega `$:
$$d_{\mathrm{up}}=\underset{\epsilon 0}{lim\; sup}\frac{_\epsilon }{\mathrm{log}\epsilon ^1}.$$
The main results of the section are the following inequalities involving the different entropies just defined. Corresponding inequalities in finite dimensional dynamical systems are well-known \[KH\].
###### Theorem 6.6
. There is a $`\gamma <\mathrm{}`$ such that
$$h_\mu h_{\mathrm{top}}\gamma d_{\mathrm{up}}<\mathrm{}.$$
(6.5)
Before giving the proof of Theorem 6.6, we state a lemma which will prove useful later on.
###### Lemma 6.7
. There are $`C`$, $`\gamma `$ such that for all (sufficiently large) $`L`$ and all (sufficiently small) $`\epsilon >0`$, if $`uv_{\mathrm{L}_{}^{\mathrm{}}(Q_L)}\epsilon `$ then for $`t>0`$, $`L^{}=LC(1+t)\mathrm{log}1/\epsilon `$, one has
$$\mathrm{\Phi }_\omega ^t(u)\mathrm{\Phi }_\omega ^t(v)_{\mathrm{L}_{}^{\mathrm{}}(Q_L^{})}Ce^{\gamma t}\epsilon $$
almost surely.
Proof. Let $`u_t`$ and $`v_t`$ be two solutions to Eq.(2.1). By Lemma 10.3,
$$u_tv_t_{\mathrm{L}_{\delta ,0}^2}e^{\gamma t}u_0v_0_{\mathrm{L}_{\delta ,0}^2}$$
and moreover both $`u_t_{\mathrm{}}`$ and $`v_t_{\mathrm{}}`$ are bounded uniformly in time (see Proposition 4.5). Let $`K_t()`$ be the convolution kernel associated with the semi-group $`\mathrm{exp}(t)`$ (see (4.1)) and let $`r_s=u_sv_s`$. By Duhamel’s formula,
$$\begin{array}{cc}\hfill \left|r_t(x)\right|& \left|K_tr_0(x)\right|+\left|_0^tK_{ts}\left(𝒢_1(u_s,v_s)r_s+𝒢_2(u_s,v_s)\overline{r}_s\right)(x)ds\right|\hfill \\ \hfill & c_1e^{\gamma t}\left(\epsilon +\underset{|xy|^2Ct\mathrm{log}1/\epsilon }{sup}|r_0(y)|\right)\hfill \\ & +\underset{0st}{sup}\left(𝒢_1(u_s,v_s)_{\mathrm{}}+𝒢_2(u_s,v_s)_{\mathrm{}}\right)_0^t\frac{|K_{ts}|}{\sqrt{\phi _{\delta ,0}}}(\sqrt{\phi _{\delta ,x}}|r_s|)(x)ds\hfill \\ \hfill & c_1e^{\gamma t}\left(\epsilon +\underset{|xy|^2Ct\mathrm{log}1/\epsilon }{sup}|r_0(y)|\right)+c_2\sqrt{\phi _{\delta ,x}}|r_0|_2_0^te^{\gamma s}\frac{|K_{ts}|}{\sqrt{\phi _{\delta ,0}}}_2ds\hfill \\ \hfill & c_3(1+t)e^{(1+\gamma )t}\left(2\epsilon +\underset{|xy|^2Ct\mathrm{log}1/\epsilon }{sup}|r_0(y)|+\underset{|xy|C\mathrm{log}1/\epsilon }{sup}|r_0(y)|\right)\hfill \\ \hfill & \mathrm{\hspace{0.17em}4}c_3e^{(2+\gamma )t}\epsilon ,\hfill \end{array}$$
where in the last line we have assumed $`|x|\frac{1}{2}LC(1+t)\mathrm{log}1/\epsilon `$ (hence $`|y|\frac{1}{2}L`$) and used the assumption $`sup_{|y|L/2}|r_0(y)|\epsilon `$.
Proof of Theorem 6.6We split Theorem 6.6 into three independent statements, namely each one of the three inequalities in (6.5).
Proof of $`h_\mu h_{\mathrm{top}}`$We follow the most standard proof (originally by Misiurewicz, quoted in \[KH\]). We modify the partition $`\mathrm{\Sigma }_{\omega ,\epsilon }=\{\sigma _1,\mathrm{},\sigma _N\}`$ by “shrinking” each element, namely by replacing each $`\sigma _k`$ by a closed set $`U_k`$ with $`U_k\sigma _k`$ and we define $`U_0=𝒜_\omega \backslash _{k=1}^NU_k`$. We thus obtain a new partition $`𝒰_{\omega ,\epsilon }=\{U_0,\mathrm{},U_N\}`$ and an open cover $`𝒱_{\omega ,\epsilon }=\{U_1U_0,\mathrm{},U_NU_0\}`$. We assume that the $`U_k`$ have been chosen such that $`H_\mu (\mathrm{\Sigma }_{\omega ,\epsilon }|𝒰_{\omega ,\epsilon })𝐏(\mathrm{d}\omega )<1`$. Remark that
$$card\left(\underset{x𝐙^dQ_L}{}\underset{j=0}{\overset{n1}{}}\mathrm{\Phi }_\omega ^{j\tau }T_x(𝒰_{\theta ^{j\tau }T_x\omega ,\epsilon })\right)\mathrm{\hspace{0.17em}2}^{nL^d}card\left(\underset{x𝐙^dQ_L}{}\underset{j=0}{\overset{n1}{}}\mathrm{\Phi }_\omega ^{j\tau }T_x(𝒱_{\theta ^{j\tau }T_x\omega ,\epsilon })\right),$$
and by Definition 6.3
$$\begin{array}{cc}& H_\mu \left(\underset{x𝐙^dQ_L}{}\underset{j=0}{\overset{n1}{}}\mathrm{\Phi }_\omega ^{j\tau }T_x(𝒰_{\theta ^{j\tau }T_x\omega ,\epsilon })\right)\hfill \\ & \mathrm{log}card\left(\underset{x𝐙^dQ_L}{}\underset{j=0}{\overset{n1}{}}\mathrm{\Phi }_\omega ^{j\tau }T_x(𝒰_{\theta ^{j\tau }T_x\omega ,\epsilon })\right)\hfill \\ & \mathrm{log}card\left(\underset{x𝐙^dQ_L}{}\underset{j=0}{\overset{n1}{}}\mathrm{\Phi }_\omega ^{j\tau }T_x(𝒱_{\theta ^{j\tau }T_x\omega ,\epsilon })\right)+nL^d\mathrm{log}2.\hfill \end{array}$$
Consequently
$$\begin{array}{cc}& \underset{L\mathrm{}}{lim}\frac{1}{L^d}\underset{n\mathrm{}}{lim}\frac{1}{n\tau }H_\mu \left(\underset{x𝐙^dQ_L}{}\underset{j=0}{\overset{n1}{}}\mathrm{\Phi }_\omega ^{j\tau }T_x(𝒰_{\theta ^{j\tau }T_x\omega ,\epsilon })\right)𝐏(\mathrm{d}\omega )\hfill \\ & \underset{L\mathrm{}}{lim}\frac{1}{L^d}\underset{n\mathrm{}}{lim}\frac{1}{n\tau }\mathrm{log}card\left(\underset{x𝐙^dQ_L}{}\underset{j=0}{\overset{n1}{}}\mathrm{\Phi }_\omega ^{j\tau }T_x(𝒱_{\theta ^{j\tau }T_x\omega ,\epsilon })\right)𝐏(\mathrm{d}\omega )+\frac{\mathrm{log}C}{\tau }.\hfill \end{array}$$
Moreover, the difference between the original partition and the new one is small, namely:
$$\begin{array}{cc}& \underset{n\mathrm{}}{lim}\frac{1}{n\tau }H_\mu \left(\underset{j=0}{\overset{n1}{}}\mathrm{\Phi }_\omega ^{j\tau }(\mathrm{\Sigma }_{\theta ^{j\tau }\omega ,\epsilon })\right)𝐏(\mathrm{d}\omega )\hfill \\ & \underset{n\mathrm{}}{lim}\frac{1}{n\tau }H_\mu \left(\underset{j=0}{\overset{n1}{}}\mathrm{\Phi }_\omega ^{j\tau }(𝒰_{\theta ^{j\tau }\omega ,\epsilon })\right)𝐏(\mathrm{d}\omega )+\frac{1}{\tau }H_\mu (\mathrm{\Sigma }_{\omega ,\epsilon }|𝒰_{\omega ,\epsilon })𝐏(\mathrm{d}\omega ).\hfill \end{array}$$
Since all the above holds for arbitrarily large $`\tau >0`$ we get
$$h_\mu \underset{\epsilon 0}{lim}\underset{L\mathrm{}}{lim}\frac{1}{L^d}\underset{n\mathrm{}}{lim}\frac{1}{n\tau }\mathrm{log}card\left(\underset{x𝐙^dQ_L}{}\underset{j=0}{\overset{n1}{}}\mathrm{\Phi }_\omega ^{j\tau }T_x(𝒱_{\theta ^{j\tau }T_x\omega ,\epsilon })\right)𝐏(\mathrm{d}\omega ).$$
(6.6)
Let next $`\delta _{\omega ,\epsilon }`$ be the Lebesgue number of the cover $`𝒱_{\omega ,\epsilon }`$ (namely the largest $`\delta _{\omega ,\epsilon }>0`$ such that every ball of diameter $`\delta _{\omega ,\epsilon }`$ is contained in an element of $`𝒱_{\omega ,\epsilon }`$). Indeed $`\delta _{\omega ,\epsilon }`$ is also the Lebesgue number of $`_{j=0}^{n1}\mathrm{\Phi }_\omega ^{j\tau }(𝒱_{\theta ^{j\tau }\omega ,\epsilon })`$ with respect to the metric $`d_{\omega ,n,\tau ,Q_L}`$. Hence
$$card\left(\underset{x𝐙^dQ_L}{}\underset{j=0}{\overset{n1}{}}\mathrm{\Phi }_\omega ^{j\tau }T_x(𝒱_{\theta ^{j\tau }T_x\omega ,\epsilon })\right)M_{\delta _{\omega ,\epsilon },Q_L,\omega },$$
and this proves that the r.h.s. of (6.6) is less than $`h_{\mathrm{top}}`$.
Proof of $`h_{\mathrm{top}}\gamma d_{\mathrm{up}}`$The proof follows \[CE2\]. Let $`\rho >0`$ be such that $`_\epsilon \left(d_{\mathrm{up}}+\rho \right)\mathrm{log}1/\epsilon `$ for all $`\epsilon <\epsilon _0`$ and then let $`L_0=L_0(\epsilon ,\rho )`$ be such that all $`L>L_0`$ yield
$$\left|\frac{\mathrm{log}M_{\epsilon ,Q_L,\omega }}{L^d}𝐏(\mathrm{d}\omega )_\epsilon \right|\rho .$$
Let $`L^{}=L+C(T+1)\mathrm{log}(1/\epsilon )`$ and $`\epsilon ^{}=C^1\mathrm{exp}(\gamma T)\epsilon `$ (see Lemma 6.7). Let an $`\epsilon ^{}`$–cover of $`𝒜_\omega |_{Q_L^{}}`$ (in the sense of Definition6.5) be given. Then it is also a $`(T/\tau ,\epsilon )`$–cover (in the sense of Definition 6.1), hence
$$N_{\omega ,T/\tau ,\tau ,Q_L,\epsilon }M_{\epsilon ^{},Q_L^{},\omega },$$
from which follows
$$\begin{array}{cc}\hfill h_{\mathrm{top}}& =\underset{\epsilon 0}{lim}\underset{L\mathrm{}}{lim}\frac{1}{L^d}\underset{T\mathrm{}}{lim}\frac{1}{T}\mathrm{log}N_{\omega ,T/\tau ,\tau ,Q_L,\epsilon }𝐏(\mathrm{d}\omega )\hfill \\ & =\underset{\epsilon 0}{lim}\underset{L\mathrm{}}{lim}\frac{1}{L^d}\underset{T}{inf}\frac{1}{T}\mathrm{log}N_{\omega ,T/\tau ,\tau ,Q_L,\epsilon }𝐏(\mathrm{d}\omega )\hfill \\ & \underset{\epsilon 0}{lim}\underset{L\mathrm{}}{lim}\frac{1}{T}\frac{\mathrm{log}M_{\epsilon ^{},Q_L^{},\omega }}{L^d}𝐏(\mathrm{d}\omega )\hfill \\ & \underset{\epsilon 0}{lim}\underset{L\mathrm{}}{lim}\frac{1}{T}\left((d_{\mathrm{up}}+\rho )\mathrm{log}1/\epsilon ^{}+\rho \right).\hfill \end{array}$$
Since $`\mathrm{log}1/\epsilon ^{}=\gamma T+\mathrm{log}(C/\epsilon )`$, the limit $`T\mathrm{}`$ and $`\rho 0`$ leaves only $`\gamma d_{\mathrm{up}}`$ on the r.h.s. above.
Proof of $`d_{\mathrm{up}}<\mathrm{}`$We want to prove a bound on $`_\epsilon `$ of the form $`_\epsilon C\mathrm{log}1/\epsilon `$ for small $`\epsilon >0`$. To do so we use iteratively the following bound:
###### Lemma 6.8
. There are $`A,B_\omega ,C>0`$ such that for all $`L>0`$ and sufficiently small $`\epsilon >0`$, one has almost surely
$$M_{\epsilon ,Q_L,\omega }M_{2\epsilon ,Q_{L+C},\theta ^1\omega }A^{L^d}B_{\theta ^1\omega }^{1/\epsilon ^2}.$$
(6.7)
The proof of Lemma 6.8 is postponed to Section 7. Let $`\epsilon >0`$, $`L>0`$. Remember that there is an $`R_\omega `$ such that $`M_{R_\omega ,Q_L,\omega }=1`$. Let $`T_\omega `$ be the smallest integer larger than $`(\mathrm{log}2)^1\mathrm{log}(R_\omega /\epsilon )`$. By iterating $`T_\omega `$ times the bound (6.7), we obtain
$$M_{\epsilon ,Q_L,\omega }\underset{n=1}{\overset{T_\omega }{}}A^{(L+(n1)C)^d}B_{\theta ^n\omega }^{1/\epsilon ^2},$$
hence
$$_\epsilon =\underset{L\mathrm{}}{lim}\frac{\mathrm{log}M_{\epsilon ,Q_L,\omega }}{L^d}𝐏(\mathrm{d}\omega )𝐄(\omega T_\omega )\mathrm{log}AC(\mathrm{log}1/\epsilon +\mathrm{log}𝐄(\omega R_\omega )),$$
(by Jensen’s inequality) or, by Definition 6.5,
$$d_{\mathrm{up}}=\underset{\epsilon 0}{lim\; sup}\frac{_\epsilon }{\mathrm{log}1/\epsilon }C,$$
which is the bound we wanted to prove. With this inequality, the proof of Theorem 6.6 is finished.
## 7 Proof of Lemma 6.8
We give the proof for the notationally convenient case $`d=1`$. Let $`u`$ and $`v`$ be two orbits of Eq.(2.1) with initial conditions $`u_0`$ and $`v_0`$ such that $`u_0`$ and $`v_0`$ belong to $`𝒜_\omega `$. The difference $`r=uv`$ satisfies almost surely the equation
$$_tr=\left(1+(1+\mathrm{i}\alpha )_x^2\right)r+𝒢_1(u,v)r+𝒢_2(u,v)\overline{r},$$
(7.1)
where we have used the notation of Eq.(10.1).
Let $`\chi (x)`$ be a smooth and monotone function satisfying $`\chi (x)=1`$ if $`x1`$ and $`\chi (x)=0`$ if $`x2`$. We decompose the kernel of $`\mathrm{exp}(t)`$ into a low frequency part and a high frequency part:
$`K_t^{()}(x)`$ $`={\displaystyle \frac{1}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}e^{\mathrm{i}px+t(1(1+\mathrm{i}\alpha )p^2)}\chi (|p/p^{}|)dp,`$
$`K_t^{(+)}(x)`$ $`={\displaystyle \frac{1}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}e^{\mathrm{i}px+t(1(1+\mathrm{i}\alpha )p^2)}\left(1\chi (|p/p^{}|)\right)dp,`$
where $`p^{}>4`$ is a sufficiently large real number. We decompose the solutions $`r_t(x)`$ to Eq.(7.1) accordingly:
$`r_t(x)`$ $`=r_t^{()}(x)+r_t^{(+)}(x),`$
$`r_t^{()}(x)`$ $`=\left(K_t^{()}r_0\right)(x)+{\displaystyle _0^t}\left(K_{ts}^{()}(𝒢_1(u_s,v_s)r_s+𝒢_2(u_s,v_s)\overline{r}_s)\right)(x)ds,`$
$`r_t^{(+)}(x)`$ $`=\left(K_t^{(+)}r_0\right)(x)+{\displaystyle _0^t}\left(K_{ts}^{(+)}(𝒢_1(u_s,v_s)r_s+𝒢_2(u_s,v_s)\overline{r}_s)\right)(x)ds`$
The kernels $`K_t^{()}`$ and $`K_t^{(+)}`$ have some regularity and decay properties that we next describe: let the Bernstein class $`B_{R,k}`$ be the following set of functions:
$$B_{R,k}\{f\mathrm{L}_{}^{\mathrm{}}:f\text{ extends to an entire function},|f(z)|Re^{k|\mathrm{Im}z|}\}.$$
(7.2)
We have
###### Lemma 7.1
. For all $`p^{}>4`$, $`t>\frac{1}{2}`$, $`f\mathrm{L}_{}^{\mathrm{}}`$, $`K_t^{()}f`$ is in $`B_{R,2p^{}}`$ with $`R2C_0f_{\mathrm{}}`$. Moreover, for all $`n𝐍`$, there is a $`C_n>0`$ such that
$`|K_t^{()}(x)|`$ $`{\displaystyle \frac{C_n}{\sqrt{t}}}(1+x^2/t)^n.`$
$`|K_t^{(+)}(x)|`$ $`{\displaystyle \frac{C_n}{\sqrt{t}}}e^{(p^{})^2t/2}(1+x^2/t)^n.`$
The proof of Lemma 7.1 is omitted, see \[CE2, Ro\].
Pick a $`2\epsilon `$–cover of $`𝒜_\omega |_{Q_{L+C(\epsilon )}}`$ (which exists a.s. by compactness, see Definition 6.5) and let $`u`$ and $`v`$ belong to one of its elements. Then $`r_0=uv`$ satisfies $`|r_0(x)|2\epsilon `$ for $`|x|\frac{1}{2}(L+C(\epsilon ))`$. Define
$$\xi _y^{(n)}(x)=\frac{1}{(1+(xy)^2)^{n/2}}.$$
Remark that Lemma 10.3 also holds with $`\phi _y`$ replaced by $`\xi _y^{(n)}`$ ($`n2`$). Moreover by reproducing the proof of Lemma 6.7 using the bounds from Lemma 7.1 we obtain (for $`|x|L/2`$):
$$\begin{array}{cc}\hfill |r_1^{()}(x)|& |K_1^{()}r_0(x)|+C_0^1K_{1s}^{()}/\sqrt{\xi _0^{(n)}}_2\sqrt{\xi _y^{(n)}}r_s_2\hfill \\ & C\epsilon +2C\epsilon _0^1\frac{C_n}{\sqrt{1s}}e^{\gamma s}ds\hfill \\ & A\epsilon ,\hfill \end{array}$$
(7.3)
where $`A`$ depends on $`n`$ but not on $`p^{}`$ and
$$\begin{array}{cc}\hfill |r_1^{(+)}(x)|& |K_1^{(+)}r_0(x)|+C_0^1K_{1s}^{(+)}/\sqrt{\xi _0^{(n)}}_2\sqrt{\xi _y^{(n)}}r_s_2\hfill \\ & e^{(p^{})^2/2}\epsilon +2C\epsilon _0^1\frac{C_ne^{(p^{})^2(1s)/2}}{\sqrt{1s}}e^{\gamma s}ds\hfill \\ & B(p^{})\epsilon ,\hfill \end{array}$$
(7.4)
where $`B(p^{})0`$ as $`p^{}\mathrm{}`$. We choose $`p^{}`$ so large that $`B(p^{})<\frac{1}{2}`$.
We next use a result of Cartwright (see \[KT\], Eq.(191)): for all $`f`$ in the Bernstein class $`B_{R,2p^{}}`$ (see (7.2)), the following identity holds:
$$f(x)=\frac{\mathrm{sin}(8p^{}x)}{32(p^{})^2}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}(1)^nf(x_n)\frac{\mathrm{sin}(4p^{}(xx_n))}{(xx_n)^2},$$
(7.5)
where $`x_n=\frac{n\pi }{8p^{}}`$. Let $`f,g`$ be in $`B_{R,2p^{}}`$. A simple application of Eq.(7.5) shows that
$$fg_{\mathrm{L}_{}^{\mathrm{}}(Q_L)}C\underset{|n|[4p^{}L/\pi ]+4Cp^{}/(\epsilon \pi )}{sup}|f(x_n)g(x_n)|+\frac{1}{4}\epsilon .$$
Hence, among all the functions in $`B_{R_\omega ,2p^{}}`$ that are bounded by $`A\epsilon `$ in $`[\frac{1}{2}L,\frac{1}{2}L]`$ (by (7.3), $`r_1^{()}`$ is such a function), at most $`(4A)^{Cp^{}L}(4R_\omega /\epsilon )^{Cp^{}/\epsilon }`$ of them are $`\epsilon /2`$–separated on $`Q_L`$. By taking a ball of diameter $`\epsilon `$ around each of them, and repeating the operation for each element of the original $`2\epsilon `$–cover, we get an $`\epsilon `$–cover of $`\mathrm{\Phi }_\omega ^1(𝒜_\omega )|_{Q_L}=𝒜_{\theta ^1\omega }|_{Q_L}`$. The number of elements in this cover is at most
$$(4A)^{Cp^{}L}(4R_\omega /\epsilon )^{Cp^{}/\epsilon }M_{2\epsilon ,Q_{L+C},\omega }.$$
The proof of Lemma 6.8 is complete.
## 8 Proof of Proposition 6.2
We follow Collet and Eckmann’s proof \[CE2\], which is itself an adaptation of standard proofs of existence of the topological entropy, see e.g. \[KH\] and references therein. The proof of Proposition 6.2 is based on the following inequalities:
###### Lemma 8.1
. For all compacts $`Q`$, $`Q^{}`$, all $`m,n𝐍`$ and $`\epsilon >\epsilon ^{}>0`$ one has
$`N_{\omega ,n,\tau ,Q,\epsilon }`$ $`N_{\omega ,n,\tau ,Q,\epsilon ^{}},`$ (8.1)
$`N_{\omega ,n,\tau ,QQ^{},\epsilon }`$ $`N_{\omega ,n,\tau ,Q,\epsilon }N_{\omega ,n,\tau ,Q^{},\epsilon },`$ (8.2)
$`N_{\omega ,n+m,\tau ,Q,\epsilon }`$ $`N_{\omega ,n,\tau ,Q,\epsilon }N_{\theta ^{n\tau }\omega ,m,\tau ,Q,\epsilon },`$ (8.3)
Furthermore for any $`\tau ^{}<\tau `$ the following inequalities hold:
$$N_{\omega ,n,\tau ^{},Q_L,\epsilon }N_{\omega ,n,\tau ,Q_{f(L)},g(\epsilon )}N_{\omega ,n,\tau ^{},Q_{f(f(L))},g(g(\epsilon ))},$$
(8.4)
where $`f(L)=L+C(\tau +1)\mathrm{log}\epsilon ^1`$ and $`g(\epsilon )=c\mathrm{exp}(\gamma \tau )\epsilon `$ with $`C,c,\gamma `$ some constants.
Lemma 8.1 implies immediately that the limit in Eq.(6.3) exists: by subadditivity (8.3) and by invariance of $`𝐏`$ under $`\theta ^t`$, we get that
$$\mathrm{\Lambda }_1=\underset{n\mathrm{}}{lim}\frac{1}{n\tau }\mathrm{log}N_{\omega ,n,\tau ,Q_L,\epsilon }𝐏(\mathrm{d}\omega )$$
exists, it is non-increasing in $`\epsilon `$ and by further subadditivity (8.2)
$$\mathrm{\Lambda }_2=\underset{L\mathrm{}}{lim}\frac{1}{L^d}\mathrm{\Lambda }_1$$
also exists and is non-increasing in $`\epsilon `$ (by (8.1)). Hence the limit in Eq.(6.3) exists. By (8.4), it is independent of $`\tau `$.
Proof of Lemma 8.1The inequality (8.1) is obvious from the definitions. We prove (8.2) by making the observation that if $`\{A_1,\mathrm{},A_N\}`$ is an $`(n,\epsilon )`$–cover of $`𝒜_\omega |_Q`$ and $`\{B_1,\mathrm{},B_M\}`$ an $`(n,\epsilon )`$–cover of $`𝒜_\omega |_Q^{}`$, then $`\{A_jB_k:j=1,\mathrm{},N,k=1,\mathrm{},M\}`$ is an $`(n,\epsilon )`$–cover of $`𝒜_\omega |_{QQ^{}}`$.
Similarly if $`\{A_1,\mathrm{},A_N\}`$ is an $`(n,\epsilon )`$–cover of $`𝒜_\omega |_Q`$ and $`\{B_1,\mathrm{},B_M\}`$ an $`(m,\epsilon )`$–cover of $`𝒜_{\theta ^{n\tau }\omega }|_Q`$, then $`\{A_j\mathrm{\Phi }_\omega ^{n\tau }B_k:j=1,\mathrm{},N,k=1,\mathrm{},M\}`$ is an $`(m+n,\epsilon )`$–cover of $`𝒜_\omega |_Q`$ which proves (8.3).
The inequality (8.4) follows immediately from Lemma 6.7, since if $`D`$ is a set of diameter $`g(\epsilon )`$ in the metric $`d_{\omega ,n,\tau ,Q_{f(L)}}`$ then $`D`$ is a set of diameter at most $`\epsilon `$ in the metric $`d_{\omega ,n,\tau ^{},Q_L}`$.
Remark. The topology of $`\mathrm{L}_{}^{\mathrm{}}(Q)`$ is a simplifying choice (as far as Eq.(8.2) is concerned), but \[CE3\] have demonstrated that other topologies can be used as well.
## 9 Proof of Proposition 6.4
This proof is, like the proof of Proposition 6.2, based on subadditive bounds. We use well-known properties of the function $`H_\mu ()`$, see \[KH\], Chapter 4.3 (in particular Proposition 4.3.3). We recall that $`xx\mathrm{log}x`$ is concave, hence for any partition $`𝒰`$ and any $`t>0`$, the following holds:
$$H_\mu \left(\mathrm{\Phi }_\omega ^t(𝒰)\right)𝐏(\mathrm{d}\omega )H_\mu \left(\mathrm{\Phi }_\omega ^t(𝒰)𝐏(\mathrm{d}\omega )\right)=H_\mu (𝒰)$$
where we have used Eq.(6.2). We thus have
$$\begin{array}{cc}& H_\mu \left(\underset{k=0}{\overset{n+m1}{}}\mathrm{\Phi }_\omega ^{k\tau }(\mathrm{\Sigma }_{\theta ^{k\tau }\omega ,\epsilon })\right)𝐏(\mathrm{d}\omega )\hfill \\ & =H_\mu \left(\underset{k=0}{\overset{n1}{}}\mathrm{\Phi }_\omega ^{k\tau }(\mathrm{\Sigma }_{\theta ^{k\tau }\omega ,\epsilon })\right)𝐏(\mathrm{d}\omega )\hfill \\ & +H_\mu (\underset{k=n}{\overset{n+m1}{}}\mathrm{\Phi }_\omega ^{k\tau }(\mathrm{\Sigma }_{\theta ^{k\tau }\omega ,\epsilon })|.\underset{k=0}{\overset{n1}{}}\mathrm{\Phi }_\omega ^{k\tau }(\mathrm{\Sigma }_{\theta ^{k\tau }\omega ,\epsilon }))𝐏(\mathrm{d}\omega )\hfill \\ & H_\mu \left(\underset{k=0}{\overset{n1}{}}\mathrm{\Phi }_\omega ^{k\tau }(\mathrm{\Sigma }_{\theta ^{k\tau }\omega ,\epsilon })\right)𝐏(\mathrm{d}\omega )+H_\mu \left(\mathrm{\Phi }_\omega ^{n\tau }\underset{k=0}{\overset{m1}{}}\mathrm{\Phi }_{\theta ^{n\tau }\omega }^{k\tau }(\mathrm{\Sigma }_{\theta ^{(k+n)\tau }\omega ,\epsilon })\right)𝐏(\mathrm{d}\omega )\hfill \\ & H_\mu \left(\underset{k=0}{\overset{n1}{}}\mathrm{\Phi }_\omega ^{k\tau }(\mathrm{\Sigma }_{\theta ^{k\tau }\omega ,\epsilon })\right)𝐏(\mathrm{d}\omega )+H_\mu \left(\mathrm{\Phi }_\omega ^{}^{n\tau }\underset{k=0}{\overset{m1}{}}\mathrm{\Phi }_\omega ^{k\tau }(\mathrm{\Sigma }_{\theta ^{k\tau }\omega ,\epsilon })\right)𝐏(\mathrm{d}\omega ^{})𝐏(\mathrm{d}\omega )\hfill \\ & H_\mu \left(\underset{k=0}{\overset{n1}{}}\mathrm{\Phi }_\omega ^{k\tau }(\mathrm{\Sigma }_{\theta ^{k\tau }\omega ,\epsilon })\right)𝐏(\mathrm{d}\omega )+H_\mu \left(\underset{k=0}{\overset{m1}{}}\mathrm{\Phi }_\omega ^{k\tau }(\mathrm{\Sigma }_{\theta ^{k\tau }\omega ,\epsilon })\right)𝐏(\mathrm{d}\omega ),\hfill \end{array}$$
namely subadditivity in the time variable. We can prove subadditivity in the space variable in a similar way. Thus the first two limits in Eq.(6.4) exist. These limits are monotonically increasing as $`\epsilon 0`$, hence the third limit is well-defined.
We next prove that the limit is independent of the choice of $`\mathrm{\Sigma }_{\omega ,\epsilon }`$: let $`\mathrm{\Sigma }_{\omega ,\epsilon }`$ and $`\stackrel{~}{\mathrm{\Sigma }}_{\omega ,\epsilon }`$ be two different sequences, we get (by the Rokhlin inequality)
$$\begin{array}{cc}& |\underset{L\mathrm{}}{lim}\frac{1}{L^d}\underset{n\mathrm{}}{lim}\frac{1}{n\tau }H_\mu \left(\underset{\genfrac{}{}{0pt}{}{x}{𝐙^dQ_L}}{}\underset{k=0}{\overset{n1}{}}\mathrm{\Phi }_\omega ^{k\tau }T_x(\mathrm{\Sigma }_{\theta ^{k\tau }T_x\omega ,\epsilon })\right)\hfill \\ & \underset{L\mathrm{}}{lim}\frac{1}{L^d}\underset{n\mathrm{}}{lim}\frac{1}{n\tau }H_\mu \left(\underset{\genfrac{}{}{0pt}{}{x}{𝐙^dQ_L}}{}\underset{k=0}{\overset{n1}{}}\mathrm{\Phi }_\omega ^{k\tau }T_x(\stackrel{~}{\mathrm{\Sigma }}_{\theta ^{k\tau }T_x\omega ,\epsilon })\right)|\hfill \\ & H_\mu (\mathrm{\Sigma }_{\omega ,\epsilon }|\stackrel{~}{\mathrm{\Sigma }}_{\omega ,\epsilon })+H_\mu (\stackrel{~}{\mathrm{\Sigma }}_{\omega ,\epsilon }|\mathrm{\Sigma }_{\omega ,\epsilon })\hfill \end{array}$$
and the r.h.s. above vanishes as $`\epsilon 0`$ since these sequences generate the whole sigma-algebra of $`𝒜_\omega `$ in this limit.
We prove that Eq.(6.4) is independent of $`\tau `$ by using Lemma 6.7 and an argument similar to the one used in Section 8.
## 10 Uniqueness of Solutions
In this section, we use the Contraction Mapping Principle to prove uniqueness of solutions to Eq.(4.5) as well as estimates on the stopping times Eq.(4.4) (along the lines of Da Prato and Zabczyk, \[DZ1\] Chapter 7). We define a Banach space $`_{p,K,T}`$ of complex-valued predictable processes $`u`$ on the time interval $`[0,T]`$ with norm defined by
$$|u|_p=\left(\underset{0tT}{sup}𝐄(e^{Kt}u_t_{\mathrm{}}^p)\right)^{1/p}.$$
We prove existence and uniqueness of solutions to Eq.(4.3) in $`_{p,K,T}`$:
###### Lemma 10.1
. Let $`T>0`$, $`p>1`$, and $`M>1`$. For sufficiently large $`K`$, there exists a unique solution $`u_t_{p,K,T}`$ to Eq.(4.3) with initial data $`u_0`$.
Proof. We define the map $`:_{p,K,T}_{p,K,T}`$ by (see Eq.(4.3))
$$\left(X\right)_t=e^tu_0+_0^te^{(ts)}P_M(|X_s|)|X_s|^{2q}X_sds+_0^te^{(ts)}\xi _sdw_s.$$
If $`𝐄(u_0_{\mathrm{}}^p)<\mathrm{}`$ then obviously $``$ maps $`_{p,K,T}`$ into itself. We define
$$\begin{array}{cc}\hfill 𝒩(|x|^2)=& (b+\mathrm{i}\beta )P_M(|x|)|x|^{2q},\hfill \\ \hfill 𝒢_1(x,y)=& \frac{1}{2}\left(𝒩(|x|^2)+𝒩(|y|^2)+(|x|^2+|y|^2)_0^1𝒩^{}\left(t|x|^2+(1t)|y|^2\right)dt\right),\hfill \\ \hfill 𝒢_2(x,y)=& xy_0^1𝒩^{}\left(t|x|^2+(1t)|y|^2\right)dt.\hfill \end{array}$$
(10.1)
If $`X`$ and $`Y`$ are arbitrary elements of $`_{p,K,T}`$, then
$$\begin{array}{cc}& |(X)(Y)|^p\hfill \\ & =\underset{0tT}{sup}𝐄\left(_0^te^{(ts)(K)}\left(𝒢_1(X_s,Y_s)e^{Ks}\left(X_sY_s\right)+𝒢_2(X_s,Y_s)e^{Ks}\left(\overline{X_s}\overline{Y_s}\right)\right)ds_{\mathrm{}}^p\right)\hfill \\ & CM^{2qp}(Kc)^p|XY|^p.\hfill \end{array}$$
The map $``$ is thus a contraction on $`_{p,K,T}`$ if $`K>C^{}M^{2q}`$. This proves the existence and the uniqueness of a solution $`u_t`$ to Eq.(4.3).
To be able to treat solutions to Eq.(4.5) as solutions to Eq.(4.2) for some time, we use the following bounds on the stopping times $`\tau (R)`$ defined by Eq.(4.4):
###### Lemma 10.2
. There is a $`C>0`$ such that the following holds almost surely for all $`R>1`$:
$$\tau (R)CR^{2q}\mathrm{log}R.$$
Proof. This follows immediately from Lemma 10.1 since we can take any $`M>R`$.
We next show that solutions of Eq.(2.1) are also uniquely defined on $`\mathrm{L}_{\delta ,y}^2`$, using that bounded functions form a dense subset.
###### Lemma 10.3
. The semi-flow $`\mathrm{\Phi }_\omega ^t`$ extends almost surely to a bounded continuous semi-flow on $`\mathrm{L}_{\delta ,y}^2`$ for any $`\delta >0`$ and $`y𝐑^d`$.
Proof. We apply the non-propagation estimate of Ginibre and Velo \[GV1\]. Let $`u_0`$ and $`v_0`$ be two functions in $`\mathrm{L}_{\delta ,y}^2`$ and denote the corresponding solutions to Eq.(2.1) by $`u_t`$ and $`v_t`$. Their difference $`u_tv_t`$ satisfies (almost surely) the following inequality:
$$\begin{array}{cc}\hfill \frac{1}{2}_t\sqrt{\phi _{\delta ,y}}(u_tv_t)_2^2& (1+\frac{1}{2}\sqrt{1+\alpha ^2})\sqrt{\phi _{\delta ,y}}(u_tv_t)_2^2\hfill \\ & \mathrm{Re}(1+\mathrm{i}\beta )\phi _{\delta ,y}(\overline{u}_t\overline{v}_t)\left(|u_t|^{2q}u_t|v_t|^{2q}v_t\right).\hfill \end{array}$$
By \[GV1\] (Proposition 3.1), Hypothesis 2.1 implies that the last term above is negative. We thus get an estimate of the form $`u_tv_t_{\mathrm{L}_{\delta ,y}^2}\mathrm{exp}(ct)u_0v_0_{\mathrm{L}_{\delta ,y}^2}`$
This and Lemma 4.3 prove that $`\mathrm{\Phi }_\omega ^t`$ is uniformly bounded and continuous on $`\mathrm{L}_{\delta ,y}^2`$ for any $`\delta >0`$ and $`y𝐑^d`$ if we define $`u_t=lim_n\mathrm{}u_t^{(n)}`$ where $`u_0^{(n)}`$ is a Cauchy sequence of bounded functions approaching $`u_0`$.
## 11 Compact Embedding for Local Spaces
In this section, we give a proof of Relation (2.7) which is a trivial adaptation of \[Ad\], Theorem 6.53, p.174. More precisely we prove the embedding (2.7) to be Hilbert–Schmidt. Let $`\{e_n\}_{n𝐍}`$ be a complete orthonormal basis of $`\mathrm{H}_{\delta ,y}^{m+k}`$. Let $`\{Q_n\}_{n𝐍}`$ be a countable cover of $`𝐑^d`$ by balls of radius $`1`$. Let $`xQ_n`$, let $`\alpha m`$ and define the bounded linear operator $`D_x^\alpha `$ on $`\mathrm{H}_{\delta ,y}^{m+k}`$ by
$$D_x^\alpha (u)=^\alpha u(x).$$
Its norm is (by Sobolev embedding) bounded by
$$D_x^\alpha (u)_{\mathrm{H}_{\delta ,y}^{m+k}}^2\underset{0\alpha m}{\mathrm{max}}\underset{xQ_n}{sup}|^\alpha u(x)|^2\frac{C}{inf_{xQ_n}\phi _{\delta ,y}(x)}u_{\mathrm{H}_{\delta ,y}^{m+k}}^2.$$
By Riesz’ Lemma, $`D_x^\alpha ()=(v_x^\alpha ,)_{\mathrm{H}_{\delta ,y}^{m+k}}`$ for some vector $`v_x^\alpha `$ and
$$\underset{n=1}{\overset{\mathrm{}}{}}|^\alpha e_n(x)|^2=\underset{n=1}{\overset{\mathrm{}}{}}|(e_n,v_x^\alpha )_{\mathrm{H}_{\delta ,y}^{m+k}}|^2=v_x^\alpha _{\mathrm{H}_{\delta ,y}^{m+k}}^2.$$
Thus the Hilbert–Schmidt norm of the embedding map is
$$\underset{n=1}{\overset{\mathrm{}}{}}e_n_{\mathrm{H}_{\delta ^{},y}^m}^2=\underset{\alpha m}{}_{𝐑^d}v_x^\alpha _{\mathrm{H}_{\delta ^{},y}^{m+k}}^2\phi _{\delta ^{},y}(x)dxm\underset{n=1}{\overset{\mathrm{}}{}}_{Q_n}\frac{C\phi _{\delta ^{},y}(x)}{inf_{zQ_n}\phi _{\delta ,y}(z)}dx,$$
which is finite whenever $`\delta ^{}>\delta `$.
## 12 Proof of Lemma 4.6
The proof can be found in \[BGO, Mi\] and is summarised below. We decompose the plane into countably many sets $`Q(m,n)`$ of unit area and use the bounds $`\phi _{\delta ,y}(x)\mathrm{exp}(\delta |xy|)e\phi _{\delta ,y}(x)`$. For simplicity we assume $`\delta =1`$ and we drop it from our notation (if Lemma 4.6 is true for $`\delta =1`$ then it is true for all $`\delta >0`$ by scaling, possibly with different constants). We simply write $`_Df`$ for $`_Df(x)dx`$ for $`D𝐑^2`$. We have
$$\begin{array}{c}\hfill _{𝐑^2}\phi _y|\mathrm{\Delta }(|f|^{2q}f)\overline{\mathrm{\Delta }f}|C\underset{m,n}{}e^{|n|}_{Q(m,n)}|\mathrm{\Delta }f||f|^{2q1}\left(|f||\mathrm{\Delta }f|+|f|^2\right),\end{array}$$
(12.1)
where $`_mQ(m,n)\{x𝐑^2:n\frac{1}{2}|xy|n+\frac{1}{2}\}`$. We estimate each summand using Hölder and Gagliardo–Nirenberg inequalities. For any $`p,r`$ with $`p^1+r^1=1`$ and in particular for $`r=1+1/q`$ and $`p=1+q`$, we get:
$$\begin{array}{cc}& _{Q(m,n)}|\mathrm{\Delta }f||f|^{2q1}\left(|f||\mathrm{\Delta }f|+|f|^2\right)\hfill \\ & c_1\mathrm{\Delta }f_{2p}\left(f_{2pq/(p1)}^{2q}\mathrm{\Delta }f_{2p}+f_{2pq/(p1)}^{2q1}f_{4pq/(p+q1)}^2\right)\hfill \\ & c_2\mathrm{\Delta }f_{2p}\left(f_{2pq/(p1)}^{2q}\mathrm{\Delta }f_{2p}+f_{2pq/(p1)}^{2q1}\left(f_{2pq/(p1)}^{1/2}\mathrm{\Delta }f_{2p}^{1/2}\right)^2\right)\hfill \\ & =c_3\mathrm{\Delta }f_{2p}^2f_{2qr}^{2q}\hfill \\ & c_4^3f_2^{2(2q+2)/(2q+3)}f_{2(q+1)}^{2(q+1/(2q+3))}\hfill \\ & K^1^3f_2^2+c_5Kf_{2(q+1)}^{4q^2+6q+2}.\hfill \end{array}$$
By summing up all contribution to (12.1) we arrive at
$$\begin{array}{cc}& _{𝐑^2}\phi _y|\mathrm{\Delta }(|f|^{2q}f)\overline{\mathrm{\Delta }f}|\hfill \\ & CK^1\underset{m,n}{}e^{|n|}_{Q(m,n)}|^3f|^2+C^{}K\underset{m,n}{}e^{|n|}\left(_{Q(m,n)}|f|^{2(q+1)}\right)^\eta \hfill \\ & \stackrel{~}{C}K^1_{𝐑^2}\phi _y|^3f|^2+C^{\prime \prime }K\underset{n}{}ne^{|n|}\left(\underset{y}{sup}_{𝐑^2}\phi _y|f|^{2(q+1)}\right)^\eta \hfill \\ & =\stackrel{~}{C}K^1_{𝐑^2}\phi _y|^3f|^2+C^{\prime \prime \prime }K\left(\underset{y}{sup}_{𝐑^2}\phi _y|f|^{2(q+1)}\right)^\eta ,\hfill \end{array}$$
which proves Lemma 4.6.
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# Hydrodynamical Survey of First Overtone Cepheids
## 1. Introduction
Historically, s Cepheids denote a certain type of low amplitude Cepheids with almost sinusoidal light-curves. Recently, the large microlensing surveys EROS (Beaulieu et al. 1995), MACHO (Welch et al. 1995) and OGLE (Udalski et al. 1997) have confirmed unequivocally that these stars are overtone Cepheids. The vast majority are first overtone pulsators that coexist with a few second overtone Cepheids at the lower period end.
Fig. 1: Observational phase difference $`\mathrm{\Phi }_{21}`$ of Galactic first overtone Cepheids: light-curves \[mag\] (upper panel) and radial velocity curves (lower panel).
For a comparison between the observational data and the calculated model pulsations a Fourier decomposition provides an accurate quantitative representation. A salient feature in the Galactic first overtone Cepheid light-curve data (labelled with a superscript m) is a large and sharp drop of the Fourier phase difference $`\mathrm{\Phi }_{21}^m`$ as a function of period in the vicinity of the 3$`.^\mathrm{d}`$2 period. The upper panel of Fig. 1 shows the observational data summarized in Poretti (1994), and supplemented with V351 Cep and Anon C Mon (Moskalik, priv. comm.). The additional Fourier data are displayed as solid triangles in the left panel of Fig. 3. The quantity $`\mathrm{\Phi }_{31}^m`$ exhibits a more or less monotonic, but large $`2\pi `$ rise. The amplitude ratios $`R_{21}^m`$ and $`R_{31}^m`$ display a local minimum in the same vicinity.
A large set of Galactic Cepheid radial velocity data has recently become available (Kienzle et al. 1999, Krzyt et al. 2000). The phase difference $`\mathrm{\Phi }_{21}^v`$ for the radial velocity (superscript v) is plotted in the lower panel of Fig. 1, and the other Fourier coefficients are displayed as solid triangles in the left column of Fig. 3.
Rapid variations in the Fourier phases are not special to the first overtone Cepheids. Actually, one of the striking features of the classical (fundamental mode) Cepheids is a Hertzsprung progression of these phases, so named after the concomitant bump progression that Hertzsprung (1926) noticed in the shape of the light-curves. For the classical Cepheids the center of this progression lies in the vicinity of the 10 day period. It was conjectured by Simon & Schmidt (1976) that this progression might have its origin in the presence of a P<sub>0</sub>/P$`{}_{2}{}^{}=2`$ resonance between the fundamental mode of oscillation and the second overtone. This conjecture was later put on a solid mathematical basis with the help of the amplitude equation formalism (Buchler & Goupil 1984, Buchler & Kovács 1986, Kovács & Buchler 1989) and was confirmed with concomitant numerical hydrodynamical modelling (Buchler, Moskalik & Kovács 1990, Moskalik, Buchler & Marom 1992).
In fact, it is now well established mathematically that sharp features in the Fourier coefficients, such as those observed in Cepheids and BL Herculis stars, are due to the appearance of resonances of the excited mode with an overtone at certain pulsation periods (e.g.Buchler (1993), Buchler 2000). Conversely the lack of such structure as in RR Lyrae is indicative of the absence of resonances.
Subsequently, Antonello & Poretti (1986), from the behavior of $`\mathrm{\Phi }_{21}^m`$ and $`R_{21}^m`$ with period (Fig. 3) and from the analogy with the Fourier data of the fundamental Cepheids, have suggested that a similar resonance, viz. P<sub>1</sub>/P$`{}_{4}{}^{}=2`$ is operative in the first overtone Cepheids and is located near P<sub>1</sub> = 3$`.^\mathrm{d}`$2. However, Kienzle et al. (1999), on the basis of the corresponding radial velocity data, suggest that the resonance center lies at a much higher period, closer to 4$`.^\mathrm{d}`$6. This incongruity suggests that it is dangerous to guess the location of a resonance without proper theoretical input. We will discuss this point further below.
In contrast to the fundamental Cepheids, the first overtone Cepheid pulsators have only received scant theoretical attention. Aikawa et al. (1987) computed the radial velocity and light curves of 11 radiative overtone pulsator models with the specific purpose of reproducing the observations of SU Cas, but were not satisfied with their results. (One labels radiative models in which convective heat transport is disregarded). Later, Antonello & Aikawa (1993) calculated two short sequences of Cepheids in order to see if numerical hydrodynamic modelling would confirm the postulated role of the P<sub>1</sub>/P$`{}_{4}{}^{}=2`$ resonance in the vicinity of 3 days. Their results displayed some structure near the resonance, but failed to reproduce the observed structure in the Fourier coefficients, in particular the $`\mathrm{\Phi }_{21}^m`$ variation. The number of computed models was rather limited, and artificially enhanced opacities were used. Subsequently, Schaller & Buchler (1994), on the basis of an extensive study of radiative first overtone Cepheids with the OPAL opacities, reached the conclusion that radiative models cannot reproduce the observed structure of the Fourier coefficient $`\mathrm{\Phi }_{21}^m`$; a similar conclusion was also reached by Antonello & Aikawa (1995). This disagreement with observation came as a surprise considering how well the fundamental mode Galactic Cepheid pulsations can be modelled (e.g., Moskalik, Buchler & Marom 1992).
In the last few years a lot of effort has been devoted to including convection in the pulsation codes, and recently one of the major remaining challenges, namely the modelling of beat pulsations, has been met (e.g., Kolláth et al. 1998, Feuchtinger 1998). In this paper we apply the same convective codes to the study of first overtone pulsations.
## 2. Physical input
Linear and nonlinear models are calculated with the Vienna pulsation code (Feuchtinger 1999a), which solves the equations of radiation hydrodynamics together with a time-dependent model equation for turbulent convection. This code recently has been extended by a linear nonadiabatic normal mode analysis, details of which are presented in a separate paper. For comparison purposes some of the calculations that are described in this paper have also been performed in parallel with the Florida pulsation code (described in Kolláth, Buchler, Szabó & Csubry 2000). The latter uses a different numerical approach, but with only minor differences in the input physics. We have ascertained that the two codes give basically the same results.
For the Rosseland mean of the opacity we use the most recent OPAL tables (Iglesias & Rogers 1996) which are augmented by the Alexander & Ferguson (1994) low temperature opacities below 6000 K. The Eddington factor is set to 1/3.
Fig. 2: Period ratio $`P_4/P_1`$ versus pulsation period for radiative models. Open triangles refer to vibrationally stable models and filled circles to models which have a stable limit-cycle. Labels on top indicate the stellar mass of the corresponding sequence.
Our model sequences have constant mass and luminosity and an equilibrium effective temperature varying in steps of 100 K. They represent horizontal paths through the instability strip (IS). We adopt a mass–luminosity (ML) relation: $`\mathrm{log}(L/L_{})=0.79+3.56\mathrm{log}(M/M_{})`$, which is derived from the stellar evolution calculations of Schaller et al. (1992) which make use of the same OPAL opacity data. For a good coverage of the observed period range we vary the stellar mass between 4.25 and 6.5 $`M_{}`$ in steps of 0.25 $`M_{}`$. The chemical composition corresponds to a typical Galactic one of (X,Y,Z) = (0.70, 0.28, 0.02).
The Fourier decomposition of the resulting light- and radial velocity curves is calculated by a least squares fit with a standard Fourier sum (8 terms). Amplitude ratios $`R_{n1}=A_n/A_1`$ and phase differences $`\mathrm{\Phi }_{n1}=\mathrm{\Phi }_nn\mathrm{\Phi }_1`$ are then used for the comparison to the observed data. Following custom, a cos Fourier decomposition is used for the light-curve data, and a sin decomposition for the radial velocity data. Note further that we compute bolometric light variations, which are compared to V-band magnitudes. For the case of RR Lyrae stars it has been shown that the differences in the low order Fourier coefficients between bolometric and V light-curves are rather small, in particular for low amplitude first overtone pulsations (Dorfi & Feuchtinger 1999, Feuchtinger & Dorfi 2000). However, for metal-rich Galactic Cepheids this has to be checked by detailed radiative transfer calculations, which will be done in a companion paper.
For the transformation between theoretical and observed radial velocities we apply a constant projection and limb darkening factor ($`u_{\mathrm{obs}}=u_{\mathrm{cal}}/1.4`$) to the calculated velocity values (Cox 1980).
## 3. Radiative models
As a first step we reexamine the difficulties encountered by radiative pulsation models, i.e., models that for simplicity disregard all convection.
Fig. 4: Linear blue edges of radiative models (R-FBE and R-OBE) compared with those of convective models, series A (C-FBE and C-OBE). The labels on the right indicate the stellar masses.
As already discussed in the Introduction, a resonance between the first and the fourth overtone is responsible for the characteristic variations in the Fourier coefficients of both light- and radial velocity variations. The location of this resonance with respect to the pulsation period, which is of particular importance for the interpretation of nonlinear results, can best be determined from linear results. In Fig. 2 the period ratio $`P_4/P_1`$ is plotted as a function of $`P_1`$ for each sequence of constant mass, and filled circles denote models with a stable overtone limit-cycle. Models close to the resonance center ($`P_1/P_4=2`$), which fall within the region of stable overtone pulsation, appear between about P<sub>1</sub> = 4 and 5$`.^\mathrm{d}`$2.
The results of the nonlinear radiative survey are summarized in Fig. 3 which depicts the low order Fourier coefficients, on the left for the light-curves, and on the right for the radial velocity curves. The observational data are represented by filled triangles, the theoretical models by open circles with dotted lines connecting the models of each sequence. We recall that the sequences consist of models with a given mass and luminosity, with $`T_{\mathrm{ef}\mathrm{f}}`$ decreasing and $`P_1`$ increasing to the right.
Even though the overall picture is not at all disastrous, several severe problems are visible. First, and most strikingly, from the flatly distributed theoretical $`\mathrm{\Phi }_{21}^m`$ it is evident that the Z-shape of the observed data cannot be reproduced at all – a disagreement which has already been mentioned in the Introduction. In addition, the theoretical $`R_{21}^m`$ values are too low for periods greater than 4 days, and the $`R_{31}^m`$ are much too low overall. In contrast, the $`\mathrm{\Phi }_{31}^m`$ show reasonable agreement.
While the overall level of the pulsation amplitudes is set by pseudo-viscosity, it is interesting that the behavior of the amplitudes $`A_1^m`$ and $`A_1^v`$ as a function of P<sub>1</sub> follows the observations rather well.
For the radial velocity plots, the general agreement with observations is much better than for the light-curves. In particular, the calculated data fit the observed $`\mathrm{\Phi }_{21}^v`$ distribution. However, several models lie off the well defined observational distribution. The same discrepancy is also visible in all the other quantities. Below we show that the inclusion of convection gives better agreement.
In summary we thus corroborate the fact that radiative models are not able to reproduce satisfactorily the observational behavior of first overtone Cepheid pulsations.
## 4. Convective models
In the last few years it has become evident that the inclusion of convective energy transport is critical to the modelling of classical stellar pulsations, rather than just being necessary for stabilizing the models at low $`T_{\mathrm{ef}\mathrm{f}}`$. The unpleasant consequence is that several free parameters ($`\alpha `$’s) have to be added to the former parameter-free radiative pulsation models. Theory unfortunately provides no guidance for choosing the values of these parameters, and therefore a calibration with observational data becomes necessary (e.g., Stellingwerf 1984, Yecko et al. 1998, Feuchtinger 1999a).
Fig. 5: Instability strip boundaries for convective models (series A), bottom: in the Log L - Log $`T_{\mathrm{ef}\mathrm{f}}`$ plane, top: in the Log P<sub>1</sub> \- Log $`T_{\mathrm{ef}\mathrm{f}}`$ plane. From left to right: (first) overtone linear blue edge (OBE), fundamental linear blue edge (FBE), nonlinear overtone red edge (NORE), overtone linear red edge (ORE) and fundamental linear red edge (FRE); the labels on left refer to the stellar masses.
For the present investigation we use the convection model according to Kuhfuß (1986) and Gehmeyr & Winkler (1992) in the version of Wuchterl & Feuchtinger (1999). Essentially the same model has been adopted by the Florida pulsation code (cf. Kolláth et al. 2000), but with a slightly different parameterization. A summary of the free parameters (subsequently termed $`\alpha `$’s) and the interrelations between the two sets of parameters are given in Table 1. For details we refer to the above cited references.
In the following we present five series of calculations, A through E, whose $`\alpha `$’s are given in Table 1. In order to reduce the multidimensional parameter space to a reasonable set of $`\alpha `$’s, we have pursued the following strategy. The parameters $`\alpha _s`$, $`c_D`$ and $`\alpha _c`$ can be chosen to reduce the model to mixing length theory in the local static limit (Kuhfuß 1986, Wuchterl & Feuchtinger 1998), for the values $`\alpha _s=1/2\sqrt{2/3}`$, $`c_D=8/3\sqrt{2/3}`$ and $`\alpha _c=\alpha _s`$. The quantities $`\overline{\alpha }_s`$, $`\overline{\alpha }_c`$ and $`\overline{c}_D`$ in Table 1 are given relative to these ’standard’ values. We adopt the standard values in series A, and in addition set the mixing length parameter $`\alpha _{\mathrm{ML}}`$ to the widely used value of 3/2. The parameter of the turbulent viscosity $`\alpha _\mu `$ is used to adjust the pulsation amplitude. Turbulent pressure, overshooting, radiative losses and the convective flux limiter are disregarded in A. Series B investigates the effects of the flux limiter and series C the effects of radiative losses. Series D has much lower turbulent energy than series A and series E additionally includes the turbulent pressure and the turbulent flux. We wish to emphasize that the adopted choices of free parameters are by no means unique.
Instability Strip
First we examine the influence of convection on the blue edge for series A and compare it to the radiative models. In Fig. 4 the radiative linear blue edges (R-FBE and R-OBE) are drawn as solid lines, and the convective ones (C-FBE and C-OBE) as dotted lines. In contrast to the frequently adopted notion that convection is only important near the red edge of the IS (cf. however Stellingwerf 1984), both the fundamental and the first overtone blue edges are shifted toward higher temperatures (toward the left in the figure) by about 350 K and 150 K for the fundamental and first overtone pulsations, respectively.
The complete linear topography of the IS for series A is presented in Fig. 5. The dotted lines refer to fundamental mode pulsation and the solid lines to the first overtone, filled/open circles to blue edge/red edges. These edges are somewhat sensitive to the values of the $`\alpha `$’s (e.g., Yecko et al. 1998) and we show a comparison of the three series below.
The nonlinear first overtone red edge (NORE) is plotted as a dashed line in Fig. 5. It is located at considerably higher temperatures than the corresponding linear one. Slight smoothing has been applied because of the rather coarse steps in effective temperature. The low mass models (with M $`<`$ 5.5$`M_{}`$) that are located at the right side of the NORE are double-mode pulsators, whereas the more massive ones (M $`>`$ 5.5$`M_{}`$) pulsate in the fundamental mode. This modal change is the reason for the kink in the NORE (cf. Kolláth et al. 1998, Kolláth et al. 2000) for a detailed picture of the modal selection problem).
It is important to exercise considerable care that the computed overtone limit-cycles are indeed stable, and not just on a transient to either double-mode or to fundamental pulsations. These transients can be very long lasting and give an erroneous impression of steady behavior. A very efficient way of determining this stability with the ’analytical signal’ method is discussed in Kolláth & Buchler (2000).
As expected and already discussed earlier (Yecko et al. 1998 and Kolláth et al. (1998)), the fundamental and first overtone blue edges intersect at some point (at $``$ 7.5 $`M_{}`$). This is consistent with the observational fact that the overtone Cepheid periods exhibit an upper limit, which is around P<sub>1</sub> = 6 days for the Galaxy (with one single star found at 7$`.^\mathrm{d}`$57). The linear overtone period at the intersection point is 8$`.^\mathrm{d}`$9 here which is considerably higher than the observations suggest. However, the region above 6.5$`M_{}`$ where stable overtone pulsations are possible is very narrow, which reduces the observational likelihood of such long period first overtone Cepheids. Furthermore the linear growth rates are found to be very small, and the corresponding nonlinear models exhibit tiny amplitudes (around 0.03<sup>m</sup> for the 7 $`M_{}`$ sequence), since the pulsation amplitude scales with the square root of the growth-rate ($`A`$ $`\sqrt{\kappa }`$). From the nonlinear survey we find that the maximum overtone period lies close to the observed one only when the pulsation amplitudes are in general agreement with observed ones. Our efforts to adjust the $`\alpha `$’s so as to lower the period at the intersection point, reduce the growth-rates and the pulsation amplitudes too much.
Fig. 7: Left: Light curves and right: radial velocity curves for a sequence parallel to the blue edge. The light-curves are shifted vertically by 0.1 mag and the radial velocity curves by 7 km/s. The curves are labelled with the periods.
Light-Curves and Radial Velocities
Fig. 6 displays the light- and radial velocity curve data for series A. Again we use filled triangles for the observations, while open circles represent our full amplitude pulsating models.
The light-curve Fourier coefficients for series A, exhibited on the left of Fig. 6, show great improvement with respect to the radiative series. The theoretical $`\mathrm{\Phi }_{21}^m`$ distribution attracts immediate attention with a very conspicuous jump around P<sub>1</sub> = 3$`.^\mathrm{d}`$4, in contrast to all the radiative models. In fact, the last points of the sequences 5 through 7 fall in the range 0.0–1.0, way below the scale. Even though the magnitude of the jump is considerably higher than what is observed, our model series reproduces qualitatively the observational $`\mathrm{\Phi }_{21}^m`$ behavior. In addition, all other light-curve Fourier coefficients show good overall agreement with observations. The values of the 31 Fourier coefficients are practically the same for series A through E and we refer to Fig. 10 for their display. Compared to the radiative models there is an average increase in $`R_{31}^m`$ by almost a factor of 10, and for small periods, the convective models also display higher pulsation amplitudes and $`R_{21}^m`$ values.
The radial velocity data on the right of Fig. 6 show good overall agreement as well. In particular, the $`R_{21}^v`$ and $`\mathrm{\Phi }_{21}^v`$ distributions closely follow the observed ones, and they produce a much better match than the radiative models. For $`R_{31}^v`$ and $`\mathrm{\Phi }_{31}^v`$ a similar behavior occurs, even though the $`R_{31}^v`$ lie somewhat below the observed ones (cf. Fig. 10). However, the $`R_{31}^v`$ are tiny which decreases the relevance of this deviation. The only perhaps significant discrepancy appears in the calculated amplitudes which, for the higher pulsation periods, are larger than the observed ones. This is also reflected in the larger $`R_{21}^v`$.
We have used the observed overall value of the pulsation amplitude to calibrate the $`\alpha `$’s (in practice $`\alpha _\mu `$). When the amplitudes are increased beyond the observed values the jump in $`\mathrm{\Phi }_{21}^m`$ becomes increasingly weak and in disagreement with the observations.
Fig. 8: Period ratio $`P_4/P_1`$ versus pulsation period for convective models (series A). Open triangles denote vibrationally stable models. Filled/Open circles refer to models with a stable/unstable overtone limit-cycle. The labels on the right indicate the stellar masses.
The shapes of the calculated light- and radial velocity curves are displayed in Fig. 7 for a sequence of models running at 200 K distance parallel to the overtone blue edge.
Finally, we note that we have computed the same series A with the Florida convective (Lagrangean) pulsation code, and that the results are essentially identical. Despite the Lagrangean nature of the latter calculations the models show a very smooth behavior, in contrast to the radiative models for which the adaptive code is necessary (Buchler, Kolláth & Marom 1996) to give smooth light-curves (cf. also Sect. 5.1).
In summary we emphasize that the inclusion of convection is crucial for a successful quantitative modelling of the pulsational properties of first overtone Cepheids, in particular of the Fourier decomposition coefficients of the light- and radial velocity curves.
Location of Resonance
We return here to the important question of whether the resonance center is near P<sub>1</sub> = 3$`.^\mathrm{d}`$2 as suggested by the light-curves (Antonello & Poretti 1986) or near 4$`.^\mathrm{d}`$6 as the radial velocity data indicate (Kienzle et al. 1999).
First, we note that our calculations which used the Schaller et al. M–L relation ($`\mathrm{log}(L/L_{})=0.79+3.56\mathrm{log}(M/M_{})`$), reproduce the observed shift with period between the light-curve and the radial velocity curve $`R_{21}`$ and $`\mathrm{\Phi }_{21}`$. From our calculated linear period ratios we should therefore be able to locate the resonance center, and resolve this issue. (We stress that it is important to use the same code, i.e., the same differencing scheme and the same mesh to compare the hydrodynamics results to the linear periods). We note in passing that the relative differences between the nonlinear and the linear periods are at most +0.4%.
The linear period ratios P<sub>4</sub>/P<sub>1</sub> versus pulsation period for our convective series A are shown in Fig. 8. The filled circles denote models with a stable nonlinear overtone limit-cycle. Note that our nonlinear first overtone IS is very narrow. We shall return later (§6) to the importance of the narrowness. Only two of our mass sequences (5.5 and 5.75 $`M_{}`$) can undergo stable overtone pulsations with periods near the resonance center (in contrast to the radiative models of Fig. 2). The corresponding pulsation periods reveal the resonance to be located around P<sub>1</sub> = 4$`.^\mathrm{d}`$2 $`\pm 0.3`$, in fact very close to the value of 4$`.^\mathrm{d}`$6 that Kienzle et al. (1999) had conjectured.
Our calculations leave no doubt that the P<sub>1</sub>/P<sub>4</sub> = 2 resonance is responsible for the observed structure of the light and radial velocity Fourier coefficients, and that the resonance is located in the vicinity of P<sub>1</sub> = 4$`.^\mathrm{d}`$2.
It is somewhat surprising that the 2:1 resonance with the fourth overtone has such a pronounced effect on the Fourier data, because after all this overtone is so strongly damped. It has a relative damping rate per pulsation period of $`\kappa _4`$P$`{}_{0}{}^{}0.4`$ in the vicinity of the resonance, i.e., its amplitude would decay by 33% in one pulsation period.
Fig. 9: Effect of pulsation amplitude on the light-curve for a series A model located at the $`\mathrm{\Phi }_{21}^m`$ jump. Upper four solid lines, have decreasing turbulent viscosity, 0.25 (top) to 0.1 (bottom) in steps of 0.05. The lowest solid line shows the corresponding light-curve with the convective flux limiter included, the dashed line refers to the same model without the flux limiter.
## 5. Sensitivity to numerical and physical input
### 5.1. Lagrangean versus adaptive mesh
We mentioned in the preceding paragraphs that a comparison between convective Lagrangean and adaptive calculations reveals no differences, as far as first overtone Cepheid models are concerned. This comes as no surprise, because pulsation amplitudes are rather small, and no strong shock waves appear in the dynamics which would require a more elaborate numerical treatment. Moreover, the inclusion of convective energy transport considerably smoothes the sharp features in the combined H–He ionization zone which are a well known headache for radiative modelling.
However, a word of caution is necessary here. As already discussed in detail in Feuchtinger & Dorfi (1994) and Buchler, Kolláth & Marom (1996), adaptive models suffer from advection errors due to the non-Lagrangean motion of the cell boundaries. These errors are particularly severe in the interior where the cell-masses increase rapidly. In order to keep these errors small, the interior part of the model has to be treated as Lagrangean. The switching point between Lagrangean and adaptive zoning therefore has to be chosen with some care, as advection errors can considerably influence the dynamical behavior and ultimately the morphology of the light- and radial velocity curve. By comparing the adaptive results to Lagrangean results we checked in detail that our results are not vitiated by advection errors.
### 5.2. Radiation hydrodynamics versus equilibrium diffusion
A standard radiation diffusion equation for radiative transport is much more convenient and faster than a time-dependent treatment of radiative transfer (radiation hydrodynamics). Since both codes are available, it has seemed interesting to check whether the simplified diffusion was adequate for pulsational behavior. On the basis of the study of several sequences of models we find that, apart from small changes in the pulsation amplitudes, the results are essentially the same for both treatments. In particular no noticeable effect on the low order Fourier coefficients has been found. A radiation diffusion treatment is therefore fully adequate.
### 5.3. The M–L Relation
Our results do not depend sensitively on the chosen M–L relation as long as the latter puts the resonance in the right place. This is so because the agreement of the hydrodynamical results with the observations necessarily puts the resonance at the right place and thus fixes the zero-point of the M–L relation (Buchler et al. 1996). The properties of the models depend very little on the slope of the M–L relation because of the relatively narrow mass range of the overtone Cepheids.
### 5.4. Convection and the $`\alpha `$ parameters
In the following we discuss how several of the convective parameters influence the behavior of first overtone Cepheid models and in particular the Fourier coefficients of the light- and radial velocity curves.
Series B
A striking feature of the convective models of series A in Section 4 is the large jump of the $`\mathrm{\Phi }_{21}^m`$, and it is interesting to see whether the size of this jump can be decreased to observed values by changing the $`\alpha `$’s.
First of all it is instructive to investigate whether there are any peculiar features in the light-curve structure that are connected with that jump. Fig. 9 (solid line at the top) shows the light-curve of a model of series A which is located just to the left of that jump. The light-curve exhibits a shoulder on the rising branch that is absent in the observed light-curves. This shoulder appears only in models near the $`\mathrm{\Phi }_{21}^m`$ jump and no corresponding feature can be found in the radial velocity curve. If the pulsation amplitude of the model is increased beyond the observed value through a decrease in the turbulent viscosity, the shoulder becomes increasingly pronounced, as the lower solid lines indicate. Eventually a spike develops that is similar to the one found in the convective models of RR Lyrae stars (Feuchtinger 1999b).
In order to cure the problem of the spike Wuchterl & Feuchtinger (1998) capped the size of the correlations $`s^{}u^{}h^{}u^{}`$ to which both the source of turbulent energy and the convective flux are proportional (flux limiter). In series B we apply the same type of limiter to the first overtone Cepheid models. However, in contrast to the RR Lyrae models we use a higher value of $`\alpha _\mathrm{L}`$ = 3 instead of 1 which diminishes the effect of the flux limiter and hence only slightly changes the convective structure of the models. Because the limiter reduces the amount of convection and therefore also the dissipation, we need to increase the turbulent viscosity parameter $`\alpha _\mu `$ from 0.25 to 0.33 to maintain the same pulsation amplitudes.
The resulting change in the light-curve structure can be inferred from the bottom of Fig. 9 which plots the flux limited light-curve (solid line) as compared to the nonlimited case (dashed line). The Fourier analysis yields a drop of $`\mathrm{\Phi }_{21}^m`$ from 5.42 to 4.20 for the limited model. The comparison of the whole flux-limited sequence with observations is given in Fig. 10. The bottom panels show $`R_{31}`$, $`\mathrm{\Phi }_{31}`$ and $`A_1`$. It turns out that the inclusion of the flux limiter decreases the jump in $`\mathrm{\Phi }_{21}^m`$ considerably, while all other quantities remain almost unaffected. Clearly the best results are obtained when a flux limiter is included.
All our attempts to achieve the same effect as obtained with a flux limiter by using various combinations of $`\alpha `$’s have proved in vain. Essentially the same was found for RR Lyrae stars (Feuchtinger 1999b). This state of affairs is somewhat disconcerting because of the ad hoc nature of the flux limiter, and its cause may well be found in the oversimplified nature of our 1D treatment of turbulent convection.
Series C
Another effect that was omitted in the model series A of Section 4 concerns the decrease of turbulent kinetic energy through radiative losses. This effect is important when the radiative diffusion time scale becomes comparable to or smaller than the typical eddy rise time, i.e., when the Péclet number is small (This effect is treated differently in the Vienna code (Wuchterl & Feuchtinger 1998) and in the Florida code (Buchler & Kolláth 2000; Kolláth et al. 2000) who follow the recipe of Canuto & Dubikov 1998). A nonzero value of the corresponding parameter $`\gamma _r`$ causes both a decrease of the convective flux and the turbulent kinetic energy. In our sequence C we use $`\gamma _r`$ = 3.5. To compensate for the resulting decrease of dissipation and to avoid too large an instability strip and too large pulsation amplitudes, we increase the mixing length parameter $`\alpha _{\mathrm{ML}}`$ from 1.5 to 2 and the turbulent viscosity $`\alpha _\mu `$ from 0.25 to 0.35 (series C, see also Table 1). This yields approximately the same pulsation amplitudes as obtained without the Péclet correction.
The influence on the linear IS boundaries is shown in Fig. 12. The solid lines refer to models including radiative losses (C), dashed to the original sequence (A), and F and O denote the fundamental and first overtone mode, respectively. Both fundamental and overtone blue edges are shifted to the blue by the same amount of about 100K. In contrast, the average fundamental red edge shift of about 550 K to the blue edge is much larger than the corresponding 200 K for the overtone red edge. Considering the average linear IS widths (taken at 6 $`M_{}`$) we end up with 580 K for the first overtone and 780 K for the fundamental, compared to 700 K and 1200 K, respectively, for the series without radiative losses. Consequently, the inclusion of radiative losses has a differential effect on fundamental and first overtone growth rates, which is important for the calibration of the whole Cepheid picture (cf. Section 7).
The nonlinear results for series C are shown in Fig. 11 and compared to observed values. Even though the topology of the IS is changed considerably, the influence on the Fourier coefficients is not conspicuous. In particular the large jump of $`\mathrm{\Phi }_{21}^m`$ is only slightly reduced compared to series A in Fig. 6. Additionally, the position of that jump and also the maximum of $`\mathrm{\Phi }_{21}^v`$ remain at the same place. Bearing in mind that several constraints involving fundamental and double-mode pulsations have not been considered so far, such insensitivity is welcome because it provides leeway for matching additional constraints (cf. Section 7).
Fig. 12: Linear IS boundaries for convective model series C (which include radiative losses, solid lines) compared to series A (dashed lines) in the HR diagram.
Series D and E
The Kuhfuß standard choice for $`\alpha _\mathrm{s}`$, $`\alpha _\mathrm{c}`$ and $`c_\mathrm{D}`$ which gives the mixing length theory (MLT) limit in the local and static case, leads to rather high values of the turbulent kinetic energy $`e_t`$. For a typical hydrostatic initial model $`e_t`$ peaks around 0.55$`e`$ in the H ionization zone and at 0.25$`e`$ in the HeI zone, where $`e`$ denotes the internal energy. Dynamical effects might lead to even higher values of $`e_t`$ during some stages of the pulsation cycle (cf. Buchler, Yecko, Kolláth & Goupil 1999, Figs. 1 and 2). The corresponding convective Mach numbers $`\sqrt{2/3e_t}/c_s`$, where $`c_s`$ denotes the adiabatic sound-speed, reach values of about 0.7. Clearly one is close to the limit of validity of our convection model, which, by disregarding pressure fluctuations, assumes a convective element always to be in pressure equilibrium with its surroundings. It is therefore interesting to compute a model series with considerably lower $`e_t`$. This can be accomplished in different ways because several $`\alpha `$ parameters (viz. $`\alpha _{\mathrm{ML}}`$, $`\alpha _\mathrm{s}`$ and $`c_D`$) exhibit a strong influence on $`e_t`$.
In series D of Table 1 we increase the dissipation parameter $`c_D`$ by a factor of 4, which leads to an average reduction of $`e_t`$ by a factor of 3. At the same time we increase $`\alpha _c`$ to 1.5 times its original value (series A), which results in approximately the same convective flux structure. Moreover, in order to obtain the right pulsation amplitudes one needs to increase the turbulent viscosity. Despite these rather dramatic changes of the $`\alpha `$’s only minor changes in the pulsational properties of the models are found.
Series E includes both turbulent kinetic energy flux $`F_t`$ and turbulent pressure $`p_\mathrm{t}`$, but has the same $`\alpha `$’s as the low $`e_t`$ series D. The flux $`F_t`$ has only a small effect on the pulsation for reasonable values of $`\alpha _\mathrm{t}`$, i.e., as long as the convection zones do not invade the outer boundary. The turbulent pressure is also unimportant as long as it remains small compared to the gas pressure. The inclusion of these quantities thus causes neither significant changes in the topography of the instability strip nor in the Fourier coefficients.
For some choice of the parameters $`\alpha _{\mathrm{ML}}`$, $`\alpha _\mathrm{s}`$, $`\alpha _\mathrm{c}`$ and $`c_\mathrm{D}`$ the turbulent kinetic energy $`e_\mathrm{t}`$ is very large. Then, because $`p_\mathrm{t}=\alpha _\mathrm{p}\rho e_\mathrm{t}`$, the turbulent pressure can get as large or even larger than the gas pressure for the standard value of the parameter $`\alpha _\mathrm{p}`$ = 2/3. Since a much smaller value of $`\alpha _\mathrm{p}`$ does not seem appropriate in this picture (/eg Baker 1987) this suggests that it would be preferably to use sets of $`\alpha `$’s that yield a lower $`e_\mathrm{t}`$ profile and a reasonable $`p/p_t`$ ratio. On the other hand, such a problem might also reflect the limitations of the simple 1D model of convection that we use.
## 6. Width of the Instability Strip
We recall that the left (hot) side of the IS determined by the linear growth-rates (which change sign there), but that the red (cool) edge is determined by nonlinear effects, namely instability of the limit-cycles. At low masses (and luminosities) the overtone limit-cycles become unstable to double-mode pulsations, and at higher masses they turn into fundamental pulsations (Udalski et al. 1987, Kolláth et al. 2000).
The comparison of our calculated Fourier data with the observations suggests that the overtone IS must be very narrow. Indeed, Figs. 6, 10 and 11 show a strong tendency for the computed values of the $`\mathrm{\Phi }_{21}^m`$ (dotted lines) to climb above the observed values as the period of the models increases along each mass sequence, in particular the low mass sequences. Had we chosen $`\alpha `$’s that yield a much broader IS then the disagreement of the computed values with the observations would have been severe.
It is somewhat puzzling that the observations show practically no low amplitude overtone Cepheids (Fig. 3), neither in light nor in radial velocity, and neither at the blue edge nor at the red edge. Of course there is some observational bias against low amplitude pulsators but we do not believe that it can account for the observed deficiency. In Buchler, Kolláth & Feuchtinger (2000) we show that the build-up of the pulsation amplitude can be delayed by stellar evolutionary effects. But this happens only on the redward entry into the IS. Another possibility is that the behavior of the growth-rates with $`T_{\mathrm{ef}\mathrm{f}}`$ is much steeper than our calculations indicate. If this were the reason it would point to an inadequacy of the simple 1D treatment of convection that we use.
## 7. Fundamental mode pulsators
Even though our first overtone Cepheid models display good agreement with observations, this tells only one part of the story. A comprehensive model for Galactic Cepheids will have to reproduce the observed behavior of the complete modal behavior (fundamental, overtone and double-mode pulsations) throughout the whole IS. Accordingly, further constraints such as the Hertzsprung progression of the Fourier coefficients of the fundamental Cepheid light- and radial velocity curves (connected with the P<sub>0</sub>/P$`{}_{2}{}^{}=2`$ resonance), or the location and properties of the double-mode pulsations need to be included. Such a calibration is beyond the scope of this paper. There is no a priori guarantee that our adopted parameter sets, which give good results for first overtone Cepheids, also work for fundamental Cepheids. We thought it useful to ascertain that with our $`\alpha `$’s the fundamental mode models are at least reasonably good. On the basis of a few sequences of models we find that even though the agreement is not perfect, the main features in the Fourier coefficients can be reproduced. There is therefore hope that future work will be able to determine a set of $`\alpha `$’s that will yield a comprehensive picture of the Galactic Cepheids.
## 8. Low metallicity Cepheids
The Magellanic Clouds are thought to be metal-deficient compared to the Galaxy, and presumably so are the SMC and LMC Cepheids. Nevertheless, the observed characteristics of these Cepheids (e.g., stellar parameters, pulsation amplitudes, position of resonances, double-mode behavior, etc.) are very close to those of their Galactic siblings. However, current models show a strong metallicity (Z) dependence that is in conflict with the observed behavior (e.g., Buchler, Kolláth, Beaulieu & Goupil 1996, Buchler 2000). This issue will be addressed in detail in a forthcoming paper.
## 9. Summary and conclusions
In this paper we have addressed the modelling of Galactic first overtone Cepheids with two different state-of-the-art stellar pulsation codes. Both codes include a treatment of time-dependent convective energy transfer, viz. the Vienna and the Florida codes. A reexamination of radiative models with an adaptive mesh and radiation hydrodynamics code reveals no improvement when compared with the simpler Lagrangean radiative diffusion code. In particular, the conspicuous Z-shape of the $`\mathrm{\Phi }_{21}^m`$ with period cannot be reproduced with radiative modelling.
In contrast, we demonstrate that with the inclusion of convective energy transport it is possible to reproduce the observed behavior of Galactic first overtone Cepheids. The Schaller et al. M–L relation that we have used here puts both the overtone P<sub>1</sub>/P<sub>4</sub>=2 and the fundamental P<sub>0</sub>/P<sub>2</sub>=2 resonances in approximately the right places as the agreement between the calculated and the observed Fourier data show. With a slight adjustment of the M–L relation the agreement with the observations could be further improved. In particular our models reveal that the P<sub>1</sub>/P$`{}_{4}{}^{}=2`$ resonance which is responsible for the structure in the Fourier coefficients, is located at pulsation periods in the vicinity of $`P_1`$= 4$`.^\mathrm{d}`$2, as conjectured by Kienzle et al. (1999) on the basis of their radial velocity data.
## 10. Acknowledgements
This work has been supported by NSF (grant AST 9819608) and by OTKA (T-026031).
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# 1 Introduction
## 1 Introduction
The generally accepted motivation for baryon asymmetric Universe is the observed absence of the macroscopic amounts of antimatter up to the scales of clusters of galaxies, which probably extends on all the part of the Universe within the modern cosmological horizon . The modern cosmology relates this baryon asymmetry of the Universe to the process of baryosynthesis., i.e. to the creation of baryon excess in very early Universe . In the homogeneous baryon asymmetric Universe the Big Bang theory predicts exponentially small fraction of primordial antimatter. Therefore, any non exponentially small amount of antimatter in the modern Universe is the profound signature for new phenomena, related to the existence of antimatter domains and leads to the respective predictions for antinuclear component of galactic cosmic rays.
The most recent analysis finds that the size of possible antimatter domains in baryon symmetrical Universe should be only few times smaller than the modern cosmological horizon to escape the contradictions with the observed gamma ray background . The distribution of antibaryon excess, corresponding to relatively small ($`<\mathrm{\hspace{0.17em}10}^5`$) volume occupied by it, can arise in inflational models with baryosynthesis and is compatible with all the observational constraints on the annihilation of antimatter in the baryon dominated Universe . The size and amount of antimatter domains is related to the parameters of models of inhomogeneous baryosynthesis (see for review ). With the account for all possible mechanisms for inhomogeneous baryosynthesis, predicted on the base of various and generally independent extensions of the standard model, the general analysis of possible domain distributions is rather complicated. But the main point of the existing mechanisms of baryosynthesis, important for our aims, is that all of them can lead to inhomogeneity of baryon excess generation and even to generation of antibaryon excess in some regions of space, when the baryon excess, averaged over the whole space, being positive (see reviews in ).
On the other hand, EGRET data on diffuse gamma background show visible peak around $`E_\gamma \mathrm{\hspace{0.17em}70}`$ MeV in gamma spectrum, which fact can be naturally explained by the decays of $`\pi ^0`$-mesons, produced in nuclear reactions. Interactions of the protons with gaseous matter in the Galaxy shift the position of such a peak to higher values of gamma energy due to $`4`$-momentum conservation. At the same time the secondary antiprotons, produced in the cosmic ray interactions with interstellar gas, are too energetic and their annihilation also cannot explain the observational data.
The above consideration draws attention to the model with antimatter globular cluster existing in our Galaxy, which cluster can serve as a permanent source of antimatter due to (anti)stellar wind or (anti)Supernova explosions. The isolated antimatter domain can not form astronomical object smaller than globular cluster . The isolated anti-star can not be formed in the surrounding matter since its formation implies the development of thermal instability, during which cold clouds are pressed by hot gas. Pressure of the hot matter gas on the antimatter cloud is accompanied by the annihilation of antimatter. Thus anti-stars can be formed in the surrounding antimatter only, what may take place when such surrounding has at least the scale of globular cluster. One can expect to find antimatter objects among the oldest population of the Galaxy , in the halo, since owing to strong annihilation of antimatter and matter gas the formation of secondary antimatter objects in the disk component of our Galaxy is impossible. So in the estimation of antimatter effects we can use the data on the spherical component of our Galaxy as well as the analogy with the properties of the old population stars in globular clusters and elliptical galaxies. The total mass of such cluster(s) is constrained from below by the condition of antimatter domain survival in the surrounding baryonic matter because small antimatter domains completely annihilate in the early Universe before the stage of galaxy formation. The upper limit on the total mass of antimatter can be estimated from the condition, that the gamma radiation from annihilation of antimatter with galactic matter gas does not exceed the observed galactic gamma background. The expected upper limit on cosmic antihelium flux from antimatter stars in our Galaxy was found only factor of two below the modern level of sensitivity in experimental cosmic antihelium searches . In the first approximation the integral effect we study depends on the total mass of the antimatter stars and does not depend on the amount of globular clusters. The only constraint is that this amount does not exceed the observed number of galactic globular clusters (about 200).
Assume that antimatter globular cluster, moves along elliptical orbit in the halo. The observed dispersion of velocity of globular clusters is $`<v>\mathrm{\hspace{0.17em}300}`$ km/s and of the long axis of their orbits is $`<r>\mathrm{\hspace{0.17em}20}`$ kpc. This gives $`T\mathrm{\hspace{0.17em}2}10^{15}`$ s as the order of the magnitude for the period of orbital motion of the cluster in the Galaxy. The period the cluster moves along the dense region of the disk with the mean half–width $`D\mathrm{\hspace{0.17em}100}`$ pc depends on the angle at which the orbit crosses the plane of the disk and is of the order
$$t_d\frac{D}{<v>}10^{13}s.$$
This means that the cluster spends not more than 1% of the time in the dense region of galactic disk, where the density of gas is of the order of $`n_H^{disk}\mathrm{\hspace{0.17em}1}`$ cm<sup>-3</sup>, moving the most time in the halo with much lower density of the matter gas $`n_H^{halo}\mathrm{\hspace{0.17em}5}10^4`$ cm<sup>-3</sup>. Therefore, we can neglect the probability to find the cluster in the disk region and consider the case when the source of the antimatter is in the halo.
One could expect two sources of the annihilation gamma emission from the antimatter globular cluster. The first one is the annihilation of the matter gas captured by the antimatter stars. Another source is the annihilation of the antimatter, lost by the antimatter stars, with interstellar matter gas. It is clear that the gamma flux originating from the annihilation of the matter gas on the antimatter stars surface is negligible. Really, an antimatter star of the Solar radius $`R=R_{}`$ and the Solar mass $`M=M_{}`$ captures matter gas with the cross section
$$\sigma \pi R\left(R+\frac{2GM}{v^2}\right)410^{22}cm^2,$$
so that the gamma luminosity of cluster of $`10^5`$ stars does not exceed $`L_\gamma M_510^{29}erg/s`$, where $`M_5`$ is the relative mass of the cluster in units $`10^5M_{}`$, $`M_{cl}=M_510^5M_{}`$. Such a low gamma luminosity being in the halo at the distance of about $`10`$ kpc results in the flux $`F_\gamma \mathrm{\hspace{0.17em}10}^{13}`$ (ster$``$ cm$`{}_{}{}^{2}`$ s)<sup>-1</sup> of $`1000`$ MeV gamma rays near the Earth, what is far below the observed background. This explains why the antimatter star itself can be rather faint gamma source elusive for gamma astronomy and shows that the main contribution into galactic gamma radiation may come only from the annihilation of the antimatter lost by the antistars with the galactic interstellar gas.
There are two sources of an antimatter pollution from the (anti-)cluster: the (anti-)stellar wind and the antimatter Supernova explosions. In both cases the antimatter is expected to be spread out over the Galaxy in the form of positrons and antinuclei. The first source provides the stationary in-flow of antimatter particles with the velocities in the range from few hundreds to few thousands km/s to the Galaxy. The (anti)Supernova explosions give antimatter flows with velocities order of 10<sup>4</sup> km/s. The relative contributions of both these sources will be estimated further on the base of comparison with the observational data assuming that all the contribution into diffuse gamma background comes from the antimatter annihilation with the interstellar matter gas. We assume in present paper that the chemical content to be dominated by anti-hydrogen and consider the contribution from the annihilation of the antiprotons only.
We consider the quasi-stationary case, provided by the presence of a permanent source of the antimatter. The assumption about stationarity strongly depends on the distribution of magnetic fields in the Galaxy, trapping charged antiparticles, annihilation cross section and on the distribution of the matter gas. We shall see that the assumption about stationarity is well justified by existing experimental data and theoretical models.
We carried out a careful consideration of the possibility to reproduce the observed spectrum of diffuse gamma background, suggesting the existence of maximal possible amount of the antimatter in our Galaxy. We showed that the predicted gamma spectrum is consistent with the observations. In this case the integral amount of galactic antimatter can be estimated, which estimation leads to definite predictions for cosmic antinuclear fluxes , accessible for cosmic ray experiments in the nearest future .
## 2 The model of galactic antimatter annihilation.
In this section we shall show that one can consider the antiproton annihilation in the halo as a stationary process and the distribution of the antiprotons does not depend practically on position and motion of the globular cluster of antistars.
One of the most crucial points for the considered model is the annihilation cross section of the antiprotons. In difference to the inelastic cross section of the $`pp`$ collisions, the cross section in the $`\overline{p}p`$ annihilation steeply grows as kinetic energy of the antiprotons goes to zero. This growth leads to the obvious fact that the main contribution into gamma flux must come from the annihilation of the slowest antiprotons. Therefore we need to have reliable estimation for the annihilation cross section of the antiprotons at low kinetic energies. Existing theoretical models based mainly on the partonic picture of the hadronic interactions are definitely invalid for $`\overline{p}p`$ annihilation at low energies and we used experimental data both for evaluation of the annihilation cross section as well as for the final state configuration.
At small energies the cross section must be proportional to the inverse power of the antiproton velocity. To find this dependence we have to match the available experimental data on $`\sigma _{ann}`$ with this expected behavior. As it follows from data , obtained at CERN-LEAR, the dependence $`\sigma _{ann}v^1`$ is valid already for laboratory antiproton momenta $`p_{lab}\mathrm{\hspace{0.17em}1000}`$ MeV/c. The annihilation cross section is the difference between total and inelastic ones, $`\sigma _{ann}\sigma _{tot}\sigma _{el}`$. Thus, at $`P_{lab}\mathrm{\hspace{0.17em}300}MeV/c`$ we used data from for the total and elastic cross sections and at momenta less than $`300`$ MeV/c we used the dependence
$$\begin{array}{ccc}\sigma _{ann}\left(P<\mathrm{\hspace{0.17em}300}\text{MeV/c}\right)\hfill & =\hfill & \sigma _0C\left(v^{}\right)/v^{}\hfill \\ \sigma _{el}\hfill & =\hfill & const,\hfill \end{array}$$
(1)
for annihilation and elastic cross sections, respectively, where $`v^{}`$ is the velocity of the antiproton in the $`\overline{p}p`$ center-of-mass system. Additional Coulomb factor $`C(v^{})`$ gives large increase for the annihilation cross section at small velocities of the antiproton and is defined by the expression :
$$C\left(v^{}\right)=\frac{2\pi v_c/v^{}}{1\mathrm{exp}\left(2\pi v_c/v^{}\right)},$$
(2)
where, $`v_c=\alpha c`$, with $`\alpha `$ and $`c`$ being the fine structure constant and the speed of light, respectively.
Using the experimental data on the $`\overline{p}p`$ annihilation cross section we found that value $`\sigma _0`$ in (1) is equal to:
$$\sigma _0=\sigma _{ann}^{exp}\left(P=\mathrm{\hspace{0.17em}300}\text{MeV/c}\right)=160\text{mb}.$$
We used the spherical model for halo with $`z`$ axis directed to North Pole and $`x`$ axis directed to the Solar system. We parametrized the number density distribution of interstellar hydrogen gas $`n_H(r,z)`$ along $`z`$ direction as:
$$\begin{array}{ccc}n_H\left(z\right)\hfill & =\hfill & n_H^{halo}+\mathrm{\Delta }_H\left(z\right),\hfill \\ & & \\ \mathrm{\Delta }_H\left(z\right)\hfill & =\hfill & \frac{n_H^{disk}}{1+\left(z/D\right)^2},\hfill \end{array}$$
(3)
with $`n_H^{halo}=\mathrm{\hspace{0.17em}5}10^4`$ cm<sup>-3</sup> being the hydrogen number density in the halo, $`n_H^{disk}=\mathrm{\hspace{0.17em}1}`$ cm<sup>-3</sup> being the hydrogen number density in the disk and $`D=\mathrm{\hspace{0.17em}100}`$ pc being the half-width of the gaseous disk. We chose here the hydrogen number density in the halo in suggestion that $`\mathrm{\hspace{0.17em}90}\%`$ of the halo mass is a non-baryonic dark matter. Such a distribution of the matter gas is to large extent the worst case for our aims since the matter density along $`z`$ axis falls slowly and visible fraction of the antiprotons will annihilate sufficiently far of the galactic disk plane. Nevertheless, as we shall see, even in this case the picture is still quasi-stationary and the antiproton number density in the halo is practically not disturbed by the annihilation in the dense regions.
The validity of the stationary approximation depends on the interplay of the life-time of the antiprotons to the annihilation and their confinement time in the Galaxy. To evaluate the antiproton confinement time we used the results of the ”two–zone” leaky box model (LBM) . The authors of considered the spectra of secondary antiprotons produced in collisions of the cosmic ray protons with interstellar gas. If to compare the antiproton spectrum, obtained in , one easily observes that shape of the spectrum beautifully reproduces the observational data on $`\overline{p}/p`$ ratio. But the predicted total normalization is lower by factor $`2÷3`$ than the data. Owing to the fact that confinement time enters as a common factor in the predicted $`\overline{p}/p`$ ratio, we found necessary factor, performing the fit to the observational data. Experimental points have been taken from where references on the data can be found. The data on $`\overline{p}/p`$ ratio we used have been collected in balloon experiments and region of low kinetic energies, $`E_{kin}\mathrm{\hspace{0.17em}100}`$ MeV, is strongly affected by the heliosphere . To avoid this influence we removed from the fit two the most left points in Fig.1. Solid curve in Fig.1(a) represents the ”two-zone” LBM predictions for the $`\overline{p}/p`$ ratio, multiplied by the fitted factor $`K=2.58`$, which factor increases the confinement time for slow antiprotons in the Galaxy up to $`5.510^8`$ years. Dashed curve is the phenomenological fit in the form $`R(E)=aE^{b+c\mathrm{lg}E}`$, which we plotted for comparison. The shapes of both curves match fairly. Fig.1(b) shows the resulting antiproton confinement times for Galaxy as whole (solid) and for disk only (dashed).
Fig.2 shows the antiproton life-time to the annihilation (a) and the free path length of the antiprotons (b) versus their distance of the galactic plane, $`z`$ for three values of the antiproton velocity. In the stationary case to compensate the annihilation of the antiprotons with matter gas the number density of the antiprotons must satisfy the equation:
$$\frac{d^2n_{\overline{p}}}{dEdt}=I_{\overline{p}}\left(E\right)v\sigma \left(v\right)n_H\frac{dn_{\overline{p}}}{dE}.$$
(4)
The solution of this equation is:
$$\frac{dn_{\overline{p}}}{dE}=I_{\overline{p}}\left(E\right)t_{ann}\left(E\right)\left(1e^{t/t_{ann}}\right),$$
(5)
with $`t_{ann}=\left[v\sigma (v)n_H\right]^1`$ being the life-time of the antiprotons relative to the annihilation.
From Fig.2(a) we can conclude that for antiprotons with velocities 10<sup>3</sup> km/s (stellar wind) the confinement time in the halo, starting from distancies $`z\mathrm{\hspace{0.17em}2}`$ kpc, is less than their annihilation time. Thus, from (5) we obtain for the halo:
$$n\left(E\right)I_{\overline{p}}\left(E\right)T_{conf}.$$
(6)
In the gaseous disk the situation is just opposite. The antiprotons annihilate with high rate and their life-time to the annihilation is much less than the time necessary to escape the Galaxy volume.
Other words, the antiprotons are storaging in the halo during the confinement time $`\mathrm{\hspace{0.17em}5}10^8`$ yrs increasing the gamma flux by factor $`T_{conf}`$. We can also conclude that during large confinement time the antiprotons are being spread over the halo with constant number density not depending on the position of the antistars cluster and under usual acceleration mechanisms in the halo their energy spectrum comes to the stationary form. Additionally from Fig.2(a) we see that the ”storaging” volume is order of the volume of the halo $`V_{halo}=\mathrm{\hspace{0.17em}4}\pi R_{halo}^3/3`$ when the region with $`T_{conf}>>T_{ann}`$ is restricted by $`|z|\mathrm{\hspace{0.17em}2}`$ kpc. Thus intensive annihilation takes place within the volume $`V_{ann}\pi R_{halo}^2\mathrm{\hspace{0.17em}4}kpc`$. The ratio of these two volumes is order of
$$\frac{V_{ann}}{V_{halo}}\frac{4kpc}{4/3R_{halo}}20\%$$
and the annihilation of the antiprotons in the gaseous disk practically does not affect the number density of the antiprotons in the Galaxy as whole.
The above consideration provides quasi-stationary distribution of antimatter in the halo and, as results, constant number density of the antiprotons in the galactic halo. Fig.2(b) shows $`z`$ dependence of free path length of the antiprotons at three values of their velocity.
## 3 Diffuse gamma flux.
The gamma flux arriving from the given direction is defined by the well known expression:
$$J_\gamma \left(E_\gamma \right)=_0^L𝑑l\psi (E_\gamma ,r,z).$$
(7)
The integration must be performed up to the edge of the halo $`L=\alpha _xR_{}+\sqrt{R_{halo}^2R_{}^2\left(1\alpha _x^2\right)}`$ with $`\alpha _x`$ being the cosine of the line-of-sight to the $`x`$ axis, directed from the Galaxy Center to the Sun and lying in the plane of the Solar orbit.
Function $`\psi (E_\gamma )`$ in (7) is the intensity of gamma sources along the observation direction $`l`$ in assumption of isotropic distribution of gamma emission. This function is defined as:
$$\begin{array}{ccc}\psi (E_\gamma ,r,z)\hfill & =\hfill & _{E_{min}}^{\mathrm{}}𝑑Ev\left(E\right)\sigma _{ann}\left(E\right)n_H(r,z)n_{\overline{p}}(E,r,z)W(E_\gamma ;E)\hfill \\ & & \\ W(E_\gamma ;E)\hfill & =\hfill & \frac{dn_\gamma (E_\gamma ;E)}{dE_\gamma dO}.\hfill \end{array}$$
(8)
To simulate the gamma energy spectrum and angular distribution $`W(E_\gamma ;E)`$ we used the Monte Carlo technics. The experimental data on the $`\overline{p}p`$ annihilation at rest (see Table) have been used to simulate the probabilities of different final states. In practice, the approximation of the annihilation at rest is valid with very good accuracy up to laboratory momenta of the incoming antiprotons about $`0.5`$ GeV because at these laboratory momenta the kinetic energy of the antiproton is still order of magnitude less than the twice antiproton mass. The simulation of the distributions of final state particles has been performed according to phase space in the center-of-mass of the $`\overline{p}p`$ system. PYTHIA 6.127 package has been used to perform the subsequent decays of all unstable particles. Momenta of stable particles ($`e^\pm `$, $`p/\overline{p}`$, $`\mu ^\pm `$, $`\gamma `$ and neutrinos) have been boosted in the laboratory reference frame. The resulting average number of $`\gamma `$’s per annihilation is
$$<n_\gamma >=𝑑\mathrm{\Omega }𝑑E_\gamma W(E_\gamma ;E)=3.93\pm \mathrm{\hspace{0.17em}0.24}$$
and agrees with experimental data.
In the stationary case we can put that annihilation rate in the halo is being constantly compensated by the permanent source of the antiprotons. But, owing to the fact that the antiprotons annihilation rate in the gaseous disk is much greater than in halo, we need to take into account the dependence of the antiproton density on $`z`$ coordinate. Fig.2(b) demonstrates that free path length of the slowest antiprotons is comparable with half-width of the disk $`D`$. To take this effect into account we have to consider the annihilation with disk gas. For given value of $`z`$ we have:
$$\frac{dn_{\overline{p}}(z,E)}{dz}=\sigma _{ann}\left(E\right)\mathrm{\Delta }_H\left(z\right)n_{\overline{p}}(z,E).$$
(9)
The differential equation (9) can be easily solved and results the following antiproton number density distribution along $`z`$ axis:
$$n_{\overline{p}}(z,E)=n_0\mathrm{exp}\left\{\sigma _{ann}\left(E\right)_z^{z_{max}}𝑑z^{}\mathrm{\Delta }_H\left(z^{}\right)\right\},$$
(10)
where, $`z_{max}=L\alpha _z`$ is the maximal value of $`z`$ coordinate, defined by the edge of the halo, and $`n_0`$ is the antiproton number density far from the disk.
The next point we need to consider is the antiproton energy spectrum. As it will be shown further, the stellar wind from antistars has to give more significant contribution in the antimatter pollution from the anticluster. The original distribution of the stellar wind particles has a Gaussian form peaking at velocities $`v\mathrm{\hspace{0.17em}500}`$ km/s . The interplanetary shocks accelerate emitted particles and the resulting stellar cosmic rays flux becomes proportional to $`J_{SW}vE_{kin}^2`$ in the range of kinetic energies up to $`\mathrm{\hspace{0.17em}100}`$ MeV . Additional acceleration occurs in the interstellar plasma and, as we believe, produces the observable spectrum of the galactic cosmic rays $`vE_{kin}^{2.7}`$. Both the acceleration mechanisms are being defined by the collisionless shocks in interplanetary or Galaxy plasmas and are charge-independent. One has to take into account also the relative movement of the hypothetical antistars cluster with velocity $`\mathrm{\hspace{0.17em}300}`$ km/s as well as the similar velocities of the matter gas defined by the gravitational field of the Galaxy. Thus, one can expect that minimal relative velocity of the antiprotons from (anti)stellar wind and the matter gas is something about $`v_{min}\mathrm{\hspace{0.17em}600}700`$ km/s. Following the above consideration, we chose the antiproton spectrum in the halo (far from regions with high matter gas density) to be similar to the galactic cosmic-rays proton spectrum in the whole range of the antiproton energies:
$$n_{\overline{p}}\left(E,z>>D\right)\left(\frac{1\text{GeV}}{E_{kin}}\right)^{2.7},$$
(11)
with the normalization:
$$_{E_{min}}^{\mathrm{}}n_{\overline{p}}\left(E,z>>D\right)𝑑E=n_0.$$
Actually, reasonable variation of the form of the antiproton flux does not affect significantly the total normalization and changes only the gamma spectrum at higher energies. The main contribution in the integrated antiproton number density comes from the slowest antiprotons owing to fast growth of the annihilation cross section with decrease of the velocity. We don’t consider in present paper the contribution in the gamma flux from the annihilation of the secondary antiprotons produced in the collisions of the cosmic-ray protons with interstellar gas. This effect must give the main contribution at higher energies of gammas and needs careful investigation of the deceleration mechanisms in the halo.
If we assume that all the gamma background at high galactic latitudes is defined by the antiproton annihilation, we have the only free parameter in our model - the minimal velocity of the antiprotons $`v_{min}`$. Therefore, for given value $`v_{min}`$ the integrated number density of the antiprotons in the halo $`n_0`$ can be found from comparison with the observational data on diffuse gamma flux. If we choose the minimal velocity of the antiprotons order of the velocity of the stellar wind, $`v_{SW}\mathrm{\hspace{0.17em}1000}`$ km/s, being equivalent to kinetic energy of the antiprotons $`E_{kin}^{SW}\mathrm{\hspace{0.17em}5.2}`$ keV, we obtain the necessary integral number density of the antiprotons $`n_0`$ to be equal to:
$$n_0^{SW}5.010^{12}\text{cm}^3.$$
(12)
Fig.3(a,b) demonstrates the resulting differential gamma distribution in the Galactic North Pole direction in comparison with EGRET data in the range $`10E_\gamma \mathrm{\hspace{0.17em}1000}`$ MeV. The peak of $`\pi ^0`$ decay is clearly seen both in calculations as well as in experimental distributions. Fig.3(c) shows the charged multiplicity distribution in the annihilation model described above. The comparison with the experimental points taken from serves as additional confirmation of our calculations.
We also performed calculations for two other values of the minimal velocity of the antiprotons $`v_{disp}=\mathrm{\hspace{0.17em}300}`$ km/s and for the velocity of the (anti)matter thrown out by the Supernovae, $`v_{SN}=\mathrm{\hspace{0.17em}2}10^4`$ km/s. The respective necessary values of the integral antiproton number density are:
$$\begin{array}{ccc}n_0^{disp}\hfill & \hfill & 2.010^{12}\text{cm}^3\hfill \\ & & \\ n_0^{SN}\hfill & \hfill & 6.010^{11}\text{cm}^3.\hfill \end{array}$$
(13)
Thus, one can see that necessary integral antiproton density in the halo practically linearly depends on minimal velocity of the antiprotons in the range $`300v\mathrm{\hspace{0.17em}10}^4`$ km/s. Note, that the approximation about annihilation at rest is valid for all the range of above minimal velocities and the resulting gamma spectrum does not change its form at such a variation of $`v_{min}`$.
## 4 Discussion and Conclusion
Let us estimate the intensity of the antiproton source and, as result, the total mass of the hypothetical globular cluster of antistars for three values of the minimal antiproton velocity: $`v_{disp}`$, $`v_{SW}`$ and $`v_{SN}`$. The first case assumes that antiprotons have been decelerated and travel in the halo with velocities equal to the velocity dispersion defined by the galactic gravitational field. The second value of $`v_{min}`$ is the order of the speed of the fast stellar wind and the third case is the velocity of the particles blown off by the Supernova explosion without possible deceleration.
If we integrate over the volume of the whole halo and take into account the antiproton storaging in the halo during the confinement time, we obtain for the integral intensity of the antiproton source $`\dot{M}\left(n_0m_pV_{halo}\right)/t_{conf}`$. For above three variants of the minimal velocity of the antiprotons and $`t_{conf}\mathrm{\hspace{0.17em}5}10^8`$ years from (12) and (13) we obtain the following values of the necessary antiproton source intensity:
$$\begin{array}{ccc}\dot{M}^{disp}\hfill & \hfill & 3.010^9M_{}/yr\hfill \\ & & \\ \dot{M}^{SW}\hfill & \hfill & 8.510^9M_{}/yr\hfill \\ & & \\ \dot{M}^{SN}\hfill & \hfill & 1.010^7M_{}/yr\hfill \end{array}$$
(14)
From the analogy with elliptical galaxies in the case of constant mass loss due to stellar wind one has the mass loss $`10^{12}M_{}`$ per Solar mass per year. In the case of stellar wind we find for the mass of the anticluster:
$$M_{clu}^{SW}210^4M_{}.$$
(15)
To estimate the frequency of Supernova explosions in the antimatter globular cluster the data on such explosions in the elliptical galaxies were used , what gives the mean time interval between Supernova explosions in the antimatter globular cluster $`\mathrm{\Delta }T_{SN}\mathrm{\hspace{0.17em}1.5}10^{15}M_5^1`$ s. For $`M_5>\mathrm{\hspace{0.17em}1}`$ this interval is smaller than the period of the orbital motion of the cluster, and one can use the stationary picture considered above with the change of the stellar wind mass loss by the $`\dot{M}f_{SN}M_{SN}`$, where $`f_{SN}=\mathrm{\hspace{0.17em}6}10^{16}M_5`$ s<sup>-1</sup> is the frequency of Supernova explosions and $`M_{SN}=\mathrm{\hspace{0.17em}1.4}M_{}`$ is the antimatter mass blown off in the explosion. Following the theory of Supernova explosions in old star populations only the supernovae of the type I (SNI) take place, in which no hydrogen is observed in the expanding shells. In strict analogy with the matter SNI the chemical composition of the antimatter Supernova shells should include roughly half of the total ejected mass in the internal anti-iron shell with the velocity dispersion $`v_i\mathrm{\hspace{0.17em}8}10^8`$ cm/s and more rapidly expanding $`v_e\mathrm{\hspace{0.17em}2}10^9`$ cm/s anti-silicon and anti-calcium external shell. The averaged effective mass loss due to Supernova explosions gives the antinucleon flux $`\dot{N}\mathrm{\hspace{0.17em}10}^{42}M_5s^1`$, but this flux contains initially antinuclei with the atomic number $`A\mathrm{\hspace{0.17em}30}\mathrm{\hspace{0.17em}60}`$, so that the initial flux of antinuclei is equal to $`\dot{A}(2\mathrm{\hspace{0.17em}3})10^{40}M_5s^1`$. Due to the factor $`Z^2A^{2/3}`$ in the cross section the annihilation life-time of such nuclei is smaller than the cosmic ray life-time, and in the stationary picture the products of their annihilation with $`Z<\mathrm{\hspace{0.17em}10}`$ should be considered. With the account for the mean multiplicity $`<N>\mathrm{\hspace{0.17em}8}`$ of annihilation products one obtains the effective flux $`\dot{A}_{eff}(1.5\mathrm{\hspace{0.17em}2.5})10^{41}M_5`$ s<sup>-1</sup>, being an order of magnitude smaller than the antiproton flux from the stellar wind.
If to take the antimatter stellar wind as small as the Solar wind $`(\dot{M_{}}=\mathrm{\hspace{0.17em}10}^{14}M_{}`$ yr$`{}_{}{}^{1})`$ this corresponds to the antiproton flux by two orders of magnitude smaller than one chosen above in (14), and the antimatter from Supernova should play the dominant role in the formation of galactic gamma background. For the Supernova case we have for the mass of the anticluster the value
$$M_{clu}^{SN}\mathrm{\hspace{0.17em}4.0}10^5M_{},$$
which value agrees with the estimation . If we assume that significant fraction of the antiprotons from stellar wind is decelerated up to $`v_{disp}`$ the respective mass of the globular cluster of antistars can be reduced up to
$$M_{clu}^{disp}\mathrm{\hspace{0.17em}7}10^3M_{}.$$
It is necessary to make small remark. Namely, in principle, one cannot exclude that the secondary antiprotons produced in $`pp`$ collisions can be decelerated in the halo magnetic fields up to velocities order of few hundreds km/s. In this case they will also give contribution in the diffuse gamma flux annihilating with the matter gas and the calculations performed in present paper are valid in this case also.
In conclusion we can say that the hypothesis on the existence of antimatter globular cluster in the halo of our Galaxy does not contradict to either modern particle physics models or observational data. Moreover, the Galactic gamma background measured by EGRET can be explained by antimatter annihilation mechanism in the framework of this hypothesis. If the mass of such a globular cluster is of order of $`10^4÷10^5M_{}`$, we can hope that other signatures of its existence like fluxes of antinuclei can be reachable for the experiments in the nearest future. The analysis of antinuclear annihilation cascade is important in the realistic estimation of antinuclear cosmic ray composition but seems to be much less important in its contribution into the gamma background as compared with the effect of antimatter stellar wind. This means that the gamma background and the cosmic antinuclei signatures for galactic antimatter are complementary and the detailed test of the galactic antimatter hypothesis is possible in the combination of gamma ray and cosmic ray studies.
Acknowledgements. The authors acknowledge the COSMION Seminar participants for useful discussions. The work was partially carried out in framework of State Scientific Technical Programme ”Astronomy. Fundamental Space Research”, Section ”Cosmoparticle Physics”. One of the authors (M.Kh.) expresses his gratitude also to COSMION-ETHZ and AMS-EPICOS collaborations for permanent support.
Table. Relative probabilities of $`\overline{p}p`$ annihilation channels.
$$\begin{array}{cccc}& & & \\ & & & \\ \text{Channel}\hfill & \hfill \text{Rel. prob.},\%& \text{Channel}\hfill & \hfill \text{Rel. prob.},\%\\ & & & \\ & & & \\ & & & \\ \pi ^+\pi ^{}\pi ^0\hfill & \hfill 3.70& 2\pi ^+2\pi ^{}\eta \hfill & \hfill 0.60\\ \rho ^{}\pi ^+\hfill & \hfill 1.35& \pi ^0\rho ^0\hfill & \hfill 1.40\\ \rho ^+\pi ^{}\hfill & \hfill 1.35& \eta \rho ^0\hfill & \hfill 0.22\\ \pi ^+\pi ^{}2\pi ^0\hfill & \hfill 9.30& 4.99\pi ^0\hfill & \hfill 3.20\\ \pi ^+\pi ^{}3\pi ^0\hfill & \hfill 23.30& \pi ^+\pi ^{}\hfill & \hfill 0.40\\ \pi ^+\pi ^{}4\pi ^0\hfill & \hfill 2.80& 2\pi ^+2\pi ^{}\hfill & \hfill 6.90\\ \omega \pi ^+\pi ^{}\hfill & \hfill 3.80& 3\pi ^+3\pi ^{}\hfill & \hfill 2.10\\ \rho ^0\pi ^0\pi ^+\pi ^{}\hfill & \hfill 7.30& K\overline{K}\mathrm{\hspace{0.17em}0.95}\pi ^0\hfill & \hfill 6.82\\ \rho ^+\pi ^{}\pi ^+\pi ^{}\hfill & \hfill 3.20& \pi ^0\eta ^{}\hfill & \hfill 0.30\\ \rho ^{}\pi ^+\pi ^+\pi ^{}\hfill & \hfill 3.20& \pi ^0\omega \hfill & \hfill 3.45\\ 2\pi ^+2\pi ^{}2\pi ^0\hfill & \hfill 16.60& \pi ^0\eta \hfill & \hfill 0.84\\ 2\pi ^+2\pi ^{}3\pi ^0\hfill & \hfill 4.20& \pi ^0\gamma \hfill & \hfill 0.015\\ 3\pi ^+3\pi ^{}\pi ^0\hfill & \hfill 1.30& \pi ^0\pi ^0\hfill & \hfill 0.06\\ \pi ^+\pi ^{}\eta \hfill & \hfill 1.20& & \\ & & & \end{array}$$
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# Are There Three Peaks in the Power Spectra of GX 339-4 and Cyg X-1?
## 1 Introduction
Observations by the EXOSAT and later by the Ginga satellites during 1980’s and early 1990’s, and more recently by the Rossi X-ray Timing Explorer (RXTE) have revealed a rich variety of X-ray variability behaviour in neutron star and black hole candidate (BHC) systems. Neutron star systems in general, and the so-called Z-sources in specific , have exhibited a set of relatively narrow features in their X-ray variability power spectral densities (PSD). These features, referred to as quasi-periodic oscillations (QPO), phenomenologically appear to fall predominantly into one of four classes. At low Fourier frequencies (see van der Klis 1995 for a review) there are the so-called normal branch oscillations (NBO) with $`f5`$$`20`$ Hz and the horizontal branch oscillations (HBO) with $`f15`$$`60`$ Hz. A second set of QPO, typically occurring in pairs at high frequencies $`f2001200`$ Hz, are referred to as the “lower-frequency kHz QPO” and the “upper-frequency kHz QPO” . Similar high frequency features have been observed in lower-luminosity neutron star systems, specifically the so-called atoll sources . The atoll systems also have exhibited low-frequency features that in some ways are similar to the HBO (see Hasinger & van der Klis 1989, Homan et al. 1998).
A wide variety of features ranging from narrow (see, for example, Nowak, Wilms, & Dove 1999; hereafter NWD) to broad “noise components” (see, for example, the discussion of van der Klis 1994a, van der Klis 1994b) have been observed in BHC systems as well. Frequently, the PSDs associated with these systems are flat at low frequencies, show a low frequency break into a steeper $`f^1`$$`f^2`$ spectrum, and often exhibit a low frequency QPO at frequencies slightly above the break (see Wijnands & van der Klis 1999, and references therein). It has been noted that the location of the break and the frequency of the QPO are often correlated , both in neutron stars and BHC systems.
In a recent work, Psaltis, Belloni, & van der Klis (hereafter PBK) have presented a further analogy between the features observed in neutron star and in BHC systems. Specifically, they have shown that the HBO frequency is apparently correlated with the lower-frequency kHz QPO frequency (Fig. 1 of PBK). With suitable identifications (in part relying upon the break frequency-QPO frequency correlation of Wijnands & van der Klis , as discussed by PBK) to features observed in atoll and BHC systems, this correlation is seen to extend over three decades in Fourier frequency and encompass Z-sources, atoll sources, and BHC sources. If these apparent correlations exist because of an underlying common physical mechanism, then all these variability phenomena are somehow intrinsic to the accretion flow, and do not explicitly rely upon the presence or absence of either a hard surface or a magnetic field in these systems.
A correlation is also observed between the lower-frequency kHz and upper-frequency kHz QPO. Some models (e.g., Miller, Lamb, & Psaltis 1998) associate this difference frequency with being nearly equal to the spin-frequency of the neutron star (in contrast to the theories of Psaltis & Norman , Stella & Vietri 1998, Merloni et al. 1999, and Stella et al. 1999 discussed further below). PBK, however, speculated about whether the observed lower-frequency kHz QPO/upper-frequency kHz QPO correlation would also extend to low frequencies and BHC systems in a manner similar to the putative HBO/lower-frequency kHz QPO correlation. In this work we reexamine a set of observations of the BHC GX 339$``$4 in order to search for evidence of a “third QPO” in this system. We present such evidence in §2. In §3, we also present evidence for multiple broad-peaked features in the PSD of Cyg X-1. We discuss in §4 the significance of these fits as regards the hypothesized QPO correlations. Furthermore, in light of the observed time lags and degree of linear correlation (i.e., the coherence function) between the soft and hard X-ray variability, we discuss whether the individual “QPO fit-components” to the PSD actually represent physically distinct processes in the accretion flow. We summarize our results in §5.
## 2 Power Spectral Densities for GX 339$``$4
In a previous work (NWD), we presented timing analysis for eight separate RXTE observations of GX 339$``$4. The faintest observation, which in terms of 3–9 keV flux is a factor $`2.5`$ fainter than the next brightest observation, clearly shows larger amplitude and characteristically lower-frequency variability than the other observations. Furthermore, the faintest observation shows time lags between its soft and hard X-ray variability \[1989, 1992, 1997, 1999a\] that are significantly shorter than those exhibited by the other observations. The remaining observations, which span a factor of $`2`$ in terms of 3–9 keV flux, have relatively similar timing properties. Specifically, the shapes and amplitudes of their PSDs are all comparable, and each shows a narrow PSD peak between 0.26–0.34 Hz. In addition, their Fourier frequency-dependent time lags span less than a factor of two at any given Fourier frequency.
In the frequency range of 0.1–30 Hz, NWD were relatively successful in fitting the PSDs to a functional form that consisted of a power law approximately $`f^1`$, plus two additional Lorentzian components of the form
$$P(f)=\pi ^1\frac{R^2Qf_0}{f_0^2+Q^2(ff_0)^2},$$
(1)
where $`f_0`$ is the resonant frequency of the Lorentzian, $`Q`$ is the quality factor ($`f_0/\mathrm{\Delta }f`$, where $`\mathrm{\Delta }f`$ is the full-width-half-maximum of the Lorentzian), and $`R`$ is the fit amplitude (root mean square variability, rms $`=R[1/2\mathrm{tan}^1(Q)/\pi ]^{1/2}`$, i.e. rms $`=R`$ as $`Q\mathrm{}`$). However, due to signal-to-noise considerations, NWD were unable to obtain an adequate fit above $`30`$ Hz.
As the seven brightest observations of GX 339$``$4 presented in NWD are so intrinsically similar, we decided to average these observations to form a composite PSD. The PSD, utilizing energy channels covering 0–21.9 keV (i.e., channels A–D from NWD summed) were calculated and averaged together using the normalization of Leahy et al. (note that this normalization weights the signal PSD by the count rate of the observation, with the average count rates here ranging from 517–877 cps), the noise was subtracted (see Nowak et al. 1999a; and references therein), and then the PSD was renormalized to the normalization of Belloni & Hasinger using the mean count rate of the summed observations. For this normalization, integrating over Fourier frequency yields the mean square variability relative to the square of the mean of the lightcurve. Error bars on the PSD were calculated as in Nowak et al. \[1999a\]
We fit this PSD with a model that consisted of a zero frequency-centered Lorentzian ($`A_0/[1+(f/f_0)^2]`$), plus four additional QPO features of the form of eq. 1. The two lowest frequency QPO features were constrained to be harmonics of one another (this allowed us to fit the “asymmetric” nature of this feature that was commented upon in NWD), but otherwise all QPO amplitudes, widths, and frequencies were allowed to be fit parameters. Our best fit model is also presented in Fig. 1. The fit yielded $`\chi ^2=177`$ for 75 degrees of freedom. The fit parameters (with subscripts 0, 1, h, 2, 3, for the zero frequency-centered Lorentzian, lowest frequency QPO, its harmonic, and higher frequency QPOs) are presented in Table 1. Error bars are 90% confidence level, i.e. $`\mathrm{\Delta }\chi ^2=2.71`$ for one interesting parameter.
We have also applied the above fit to the faintest observation (observation 5) from NWD. Here we are able to fit the zero frequency-centered Lorentzian, the low frequency QPO (plus harmonic), and the middle frequency QPO. The results are also presented in Table 1, and the data are shown in Fig. 2. The high-frequency QPO, if present, is essentially unconstrained (although in Table 1 we give the error bars for the “local minimum” in $`\chi ^2`$). Fig. 2 also shows the residual noise level. This level corresponds to the expected amplitude of postive 1-$`\sigma `$ fluctuations above the mean value of the Poisson noise PSD, and therefore is indicative of the minimum PSD amplitude at which a signal can be detected (see Nowak et al. 1999a, and references therein). This residual noise level, which scales inversely proportionally to the count rate and to the square root of the integration time, shows the difficulty of detecting the presence of a “third QPO feature” in the high frequency PSD of a single observation. The composite PSD shown in Fig. 1 has a residual noise level $`9`$ times lower than that for the faintest observation of GX 339$``$4.
## 3 Power Spectral Densities for Cyg X-1
Nowak et al. \[1999a\] discussed observations of Cyg X-1 in its low luminosity/spectrally hard state. In that work, we presented doublely broken power law fits to the PSD, with the low frequency PSD being essentially flat, the middle-frequency PSD being approximately $`f^1`$, and the high-frequency PSD being approximately $`f^2`$. Although the *fractional* deviation of the data from the fits was quite small, the reduced $`\chi ^2`$ ranged from $`4`$$`16`$ due to the excellent statistics achievable with RXTE. Here we refit the PSD with the functional form discussed in §2. The fit parameters for two separate energy channels (0–4 keV and 14–45 keV; see Nowak et al. 1999a) are also presented in Table 1, and the fits are shown in Fig. 3.
The Cyg X-1 PSDs show subtle, but statistically significant, features. Broad peaks are seen at $`0.3`$ Hz, 2 Hz, 10 Hz, and 50 Hz. The nearly flat portion of the PSD at $`<1`$ Hz shows some sign of a peak. We use the same model that we fit to GX 339$``$4, namely a zero frequency-centered Lorentzian, a QPO and its harmonic, plus two additional QPO, and obtain reasonable fits. The reduced $`\chi ^2`$ for these fits were significantly lower ($`\mathrm{\Delta }\chi ^2>150`$ for eight extra degrees of freedom) than that for the doublely broken power law fits discussed in Nowak et al. \[1999a\]. Note that the lowest frequency ‘QPO’ is identified with the portion of the PSD at $`<1`$ Hz, and its harmonic is identified as having both larger amplitude and greater $`Q`$ value, in contrast to the fits to GX 339$``$4.
Although these fits are a significant improvement over the doublely broken power law fits of Nowak et al. \[1999a\], there are large residuals, especially near 0.2 Hz. We thus considered another fit consisting of five QPO features plus one harmonic (or possibly sub-harmonic, as the QPO associated with the higher of the two related frequencies had a larger amplitude). The results for these fits are presented in Table 2, and the fits are also shown in Fig. 3. This model resulted in significant improvements to the fits ($`\mathrm{\Delta }\chi ^250`$—90, for four extra degrees of freedom); however, compared to the fits presented in Table 1, there was relatively little change in the parameters for the highest frequency QPO fit components.
In these models, what was fit as a zero frequency-centered Lorentzian is now fit as a strong, broad QPO plus a weaker, narrow QPO. Whether these fit components should be regarded as separate physical features, or whether the improvement in the fit is merely an indication of the inadequacy of the zero frequency-centered Lorentzian in describing the low frequency PSD is, of course, unclear. Also along these lines, we note that although we have constrained the features at $`0.8`$ Hz and 1.6 Hz to be harmonics of each other, the fits are slightly improved ($`\mathrm{\Delta }\chi ^2=7.0`$ for the 0–4 keV band and $`\mathrm{\Delta }\chi ^2=8.9`$ for the 14–45 keV band) if the frequencies are allowed to float freely. Note also that there are still positive residuals at $`>10^3`$ Hz. The above points, as we discuss further below, serve to highlight the fact that when the PSD features are as broad and as subtle as they are in both Cyg X-1 and GX 339$``$4 the intepretation of what is and is not to be considered a “QPO” becomes much more ambiguous. This needs to be borne in mind in the discussion below of the correlations between fit features.
## 4 Discussion
The evidence for a third, high-frequency feature in the composite PSD of GX 339$``$4 is clear. Removing this feature, the $`\chi ^2`$ increases to $`442`$. Here we should note that to some extent the fits are affected by the fact that separate observations with slightly varying PSD properties have been averaged together. Some broadening of the “QPO features” is expected as a result, and the $`Q`$ values for features fit to the summed PSD should be considered as lower limits. Likewise, the fitted frequencies and amplitudes of the variability features should be viewed as indicative of an average value, but not strictly applicable to any of the individual PSD used in the combined observation. Furthermore, their error bars should be viewed as lower limits. Along these lines, we note that although the presence of a third feature is indicated, it is almost as well-modelled ($`\chi ^2=207`$) by a zero frequency-centered Lorentzian ($`[1+(f/f_b)^2]^1`$) with a break frequency of $`f_b=23\pm 3`$ Hz.
The question arises as to what extent this third variability feature might be an artifact of averaging together seven separate observations. This is of some concern; however, we note that in terms of well-measured properties, the seven individual PSDs are very similar to each other. Using the normalization of Belloni & Hasinger , at frequencies $`<10`$ Hz (where all seven individual PSDs have good statistics) the variance of the noise subtracted PSDs is $`<25\%`$. Specifically,
$$P_s(f)^1\left(\underset{i}{}[P_s(f)P_i(f)]^2\right)^{1/2}<0.25,$$
(2)
at all frequencies $`f10`$ Hz, where $`P_i(f)`$ are the noise subtacted PSD of the individual observations, and $`P_s(f)`$ is the noise subtracted PSD of the summed observation. The variance of the mean is lower by a factor of $`\sqrt{6}`$ and therefore is $`<10\%`$. This result is essentially unchanged if we weight the observations by count rate. The fact that some variation exists from observation to observation means that the error bars that we used should be viewed as lower limits, and the derived $`\chi ^2`$ should be viewed as upper limits. Also as regards possible systematic variations, NWD showed the broad features at $`0.35`$ and $`2.5`$ Hz have fitted frequencies that ranged only from $`0.3`$$`0.4`$ Hz and $`2`$$`3`$ Hz, respectively, among the individual observations *and* among separate energy bands within each of these observations. Given the similarities among the PSDs, it would be unusual for the strong third variability feature present in the combined observation to be an artifact of the PSD averaging. (We note that the third variability feature discussed below for Cyg X-1 was for a single, short observation.)
Taking all three features in GX 339$``$4 as real and not an artifact of the PSD summation, these features are relatively broad (even considering possible systematic broadening due to averaging) compared to high frequency QPO seen, for example, in neutron stars. Statistically, they all have $`Q>0`$ at a high significance level; however, they are all far from narrow, periodic features. The QPO at the higher frequency end of the correlation presented by PBK all tend to be much narrower features (in terms of FWHM) than those presented here. Nevertheless, if we take these features and identify the low-frequency QPO with the HBO, the middle-frequency QPO with the lower-frequency kHz QPO, and the high-frequency QPO with the upper-frequency kHz QPO, then they approximately fall along the correlations suggested by PBK. We show the above data overlaid on the suggested correlation of PBK in Fig. 4. This is true for the composite PSD as well for the PSD of the faintest observation from NWD, keeping in mind that for this latter observation the error bars on the highest frequency feature are for the local minimum and that the presence of this latter feature is not strongly constrained. We also note that the correlation between the frequency, $`f_0`$, of the zero frequency-centered Lorentzian and the frequency, $`f_1`$, of the lowest frequency QPO is in the same sense as the break frequency/QPO frequency correlation discussed by Wijnands & van der Klis .
The fact that the “HBO” and “lower-frequency kHz QPO” of GX 339$``$4 appear to lie approximately along the suggested correlation was already noted by PBK. What is new here is that there appears to be a third QPO that extends the lower-frequency kHz QPO/upper-frequency kHz QPO correlation downward by nearly two orders of magnitude in Fourier frequency. The high-frequency QPO of GX 339$``$4 appears to lie slightly below the extrapolation of the lower-frequency kHz QPO/upper-frequency kHz QPO correlation; however, assuming the correlation to be $`f^\alpha `$, a slope change of $`\delta \alpha 0.2`$ is all that is required between the low-frequency GX 339$``$4 point and the high-frequency trend. Furthermore, some theories (e.g., the relativistic precession theory of Stella & Vietri 1998; Merloni et al. 1999; Stella et al. 1999) suggest that one should not expect a single power law over the whole range of the putative correlation. Specific predicted frequency correlations for the relativistic precession theory can be found in Stella et al. .
The features fit to the Cyg X-1 PSD are slightly more problematic to fit onto this trend. In light of the fact that the best fit consists of five QPO features plus one harmonic, there are the questions of which fit components are to be considered as “true QPO” (as opposed to merel a fit artifact) and which are to be associated with which feature in the putative correlation. In Fig. 4, we have taken the fits of Table 2 and associated $`\mathrm{QPO}_5`$ with the upper-frequency kHz QPO, $`\mathrm{QPO}_4`$ with the lower-frequency kHz QPO, and $`\mathrm{QPO}_3`$— the sub-harmonic to the stronger 0.16 Hz feature— with the HBO. With these associations, the features fit to the high energy (14–45 keV) PSD of Cyg X-1 agree with suggested trend, whereas the features fit to the low energy (0–4 keV) PSD of Cyg X-1 deviate from the trend, mostly due to the extremely low frequency of the “lower-frequency kHz QPO”. We note, however, that the frequency-break of this component occurs at roughly the same location in the low and high energy PSDs. The fact that the feature in the low energy PSD has an extremely low $`Q0.2`$ leads to a very low fit-frequency in order to yield a break in the same location as the somewhat narrower feature ($`Q0.6`$) in the high energy PSD. If one performs a joint fit to the 0–4 keV and 14–45 keV PSDs, constraining the QPO frequencies (but not the amplitudes or widths) to be the same in both bands, a reasonable fit is obtained with $`\chi ^2=156`$ for 111 degrees of freedom. The resulting fit frequencies are comparable to the frequencies fit to the 14–45 keV band in Table 2.
It is difficult to test the correlations suggested by PBK since most of the $`Q`$ values for the fitted features are very low, making the features very subtle. Another, more serious problem, is the question of identification of the features. One alternative fit to the Cyg X-1 PSDs would be to place a zero frequency-centered Lorentzian with break frequency $`0.01`$ to remove the low frequency residuals seen in Fig. 3, and then to identify $`\mathrm{QPO}_1`$ and $`\mathrm{QPO}_2`$ as the HBO and a harmonic, $`\mathrm{QPO}_3`$ and $`\mathrm{QPO}_\mathrm{h}`$ as the lower-frequency kHz QPO and a harmonic, $`\mathrm{QPO}_4`$ as the upper-frequency kHz QPO, and $`\mathrm{QPO}_5`$ as a new, “fourth QPO”. \[The theoretical existence of a fourth such frequency has been suggested by Psaltis & Norman (2000), for example.\] These frequencies would also fit on the suggested trend of PBK, except now $`\mathrm{QPO}_4`$ in the low energy PBK would be too low a frequency to fall on the trend for the upper-frequency kHz QPO. *At present, there is no unambiguous, rigorous method of associating a fitted feature with the suggested trend of PBK.*
Even given these caveats regarding correlations between the various fit-components, it is tempting to associate each fit-component of the PSD with a “resonance” in the variability properties of the accretion flow (see Psaltis & Norman 2000). This is counter to the models of Kazanas et al. and Poutanen & Fabian , for example, that essentially postulate a single “response” for the variability properties of the accreting system. NWD phenomenologically elaborated upon the concept of each PSD fit-component corresponding to a separate physical mechanism by further postulating that each of these system “responses” had a separate “driver”, uncorrelated with the variability-drivers of the other PSD components<sup>1</sup><sup>1</sup>1Even if a system has a set of independent responses to a source of input fluctuation or noise, for example the accretion disk responses discussed by Psaltis & Norman , the net outputs will be perfectly correlated (i.e., have a coherence function of unity) if they are responding to the same source of noise fluctuations. This point is discussed in further detail by Bendat & Piersol , Vaughan & Nowak , and NWD.. The net Fourier frequency-dependent phase lags (or, equivalently, time lags$`=`$phase lags/$`2\pi f`$) and coherence function \[$`\gamma ^2(f)`$, a measure of the degree of *linear* correlation between two lightcurves; see Bendat & Piersol 1986, Vaughan & Nowak 1997\] between soft and hard variability were then given by a combination of the phase lags and coherence functions for each individual fit-component of the PSD (see NWD, §4 and Fig. 8; Vaughan & Nowak 1997).
NWD noted that for the case of GX 339$``$4 wherever one PSD fit component dominated (for example, near 0.3 Hz in Fig. 1, where the low-frequency QPO dominates), the coherence function would be near unity and there would be an approximately flat Fourier phase lag shelf when comparing soft and hard variability. Wherever two PSD fit components crossed one another, there would be a slight dip in the coherence function and a transition from one characteristic Fourier phase lag to another (see Fig. 5). Fig. 1 shows that at any given Fourier frequency, typically two PSD fit components dominate. Ignoring the other two fit components, at any given frequency we can then calculate the phase lag between soft and hard X-ray variability for the remaining two PSD fit components by simultaneously fitting the measured net phase lag and the measured net coherence function<sup>2</sup><sup>2</sup>2Note that there is still an ambiguity in the values of the Fourier phases, as there are two independent solutions for a two-component fit to the measured phase lag and coherence function. Specifically, see Fig. 8 of NWD, where reflecting the phases of the individual PSD components through the net measured phase lag yields and equally valid solution. In our solution that follows, at low Fourier frequency we choose the phase lag for the zero frequency-centered Lorentzian to be the value closest to the net measured phase lag. Then, progressing to higher Fourier frequencies, we choose solutions that are most nearly continuous with the lower frequency values..
We have carried out such a procedure for the second brightest GX 339$``$4 observation (which showed the strongest QPO at $`0.3`$ Hz) presented in NWD. We have fit the 0–21.9 keV PSD of this observation with the same type of model as discussed in Table 1. For this non-composite PSD fit, constraints are again weakest for the “third QPO” fit-component. Furthermore, we have *assumed* this PSD shape and amplitude for both the 0–3.9 keV and 10.8–21.9 keV PSDs. (Fitting each energy band individually, no constraints can be made on the “third QPO”.) We then followed the procedures outlined in §4 of NWD (see specifically eqs. 4 and 5) for fitting the phase lag and coherence function, except that instead of fitting four constant (as a function of Fourier frequency) phases over the entire data set, we fit two phases at any given Fourier frequency. Specifically, at any given Fourier frequency, we fit the phase lags for the two strongest PSD fit-components, *assuming the other two PSD fit components to be negligible*. Error bars do not account for uncertainties in the PSD fit parameters, and only represent that portion of the error due to the uncertainty in the measured phase lag and measured coherence function. The results for this decomposition (which, as we discuss above, is *not* unique) are presented in Fig. 5.
Fig. 5 is presented more for illustrative purposes rather than as a completely rigorous and definitive fit. The things to note about this possible decomposition are the following. If the QPO fit components each represent a separate, independent physical process, then the deviations from unity coherence can be understood within the context of such a model as well. At low frequency the net phase/time lags are most characteristic of the zero frequency-centered Lorentzian, at middle frequencies they are most representative of the low frequency QPO, and at moderately high frequencies they are most representative of the middle frequency QPO. At $`0.3`$ Hz, there is a recovery of the coherence function to near unity values that would be associated with the low frequency QPO becoming dominant in the PSD. Although it is difficult to accurately measure the high frequency phase/time lags and coherence function, the fact that the middle and high frequency fit components to the PSD have comparable amplitudes at $`f>10`$ Hz could explain the strong loss of coherence at these frequencies (which is characteristic of Cyg X-1 as well; see Nowak et al. 1999a). Although the net phase lag indicates that hard photons lag the soft photons, it is possible for portions of the individual variability components to show exactly the opposite behavior.
## 5 Summary
The main results of this work can be summarized as follows.
* A composite power spectrum of GX 339$``$4 is well-fit by a model that consists of a zero frequency-centered Lorentzian, a moderate width ($`Q1`$$`2`$) quasi-periodic oscillation and its harmonic, plus two additional, broad ($`Q0.5`$–1) QPOs.
* The frequencies of these three QPO components approximately agree with the correlations among horizontal branch oscillations/lower-frequency kHz QPOs/upper-frequency kHz QPOs suggested by Psaltis, Belloni, & van der Klis .
* Fits to the PSD of Cyg X-1 are more problematic. Multiple QPO components provide the best fits to the PSD as a function of observed energy band. Depending upon the identifications made, the high energy band PSD fits can be made consistent with the correlations suggested by Psaltis, Belloni, & van der Klis . The low energy band PSD fits show deviations from these correlations, although a joint fit of the low- and high-energy PSD does allow for a single set of frequencies consistent with the suggested correlations.
* One can construct a phenomenological model where the net Fourier frequency-dependent phase lag and coherence between soft and hard X-ray variability can be expressed as the combination of phase lags and coherences from the individual fit components to the PSD.
* If such a model for the phase lags and coherences is correct, then although the net phase lag indicates that soft variability leads the hard variability, the trend may be exactly opposite for some of the individual variability components. A *requirement* of this model is that although the individual components have coherent variability between soft and hard photons, the variability components are incoherent with one another. This not only requires that each component be representative of a different “response” within the accretion flow, but also that each component has a separate variability “driver”
The correlations between the different fit-components seems at least to be very suggestive. It is tempting to associate all of these quasi-periodic features observed in both neutron star and black hole systems with a common set of phenomena intrinsic to the accretion flow itself, and not to the properties of the compact object such as the presence of a surface or a magnetic field. Such may not be the case in reality, however. The presence of what appear to be four distinct features (the zero frequency-centered Lorentzian and three broad, quasi-periodic features) in the PSD of GX 339$``$4, coupled with suggestive evidence from the Fourier frequency-dependent phase lags and coherence function (see also NWD), at the least seems to be arguing for models with multiple, distinct sources of variability (e.g., Psaltis & Norman 2000), and against models with a single “type” of variability (e.g., Kazanas et al. 1997, Poutanen & Fabian 1999).
This research has been supported by NASA grants NAG5-3225 and NAG5-4731. Dimitrios Psaltis kindly provided the data for Fig. 4 and valuable comments. I also would like to thank Andrew Hamilton for useful discussions, the hospitality of the Aspen Center for Physics, the participants of the “X-ray Probes of Relativistic Astrophysics” workshop for many useful and stimulating discussions, and the hospitality of P. Coppi, C. Bailyn, and the Yale Astronomy department while this work was being completed.
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# Persistent current in metals with a large dephasing rate
## Abstract
In a weakly disordered metal electron interactions are responsible for both decoherence of the quasi-particles as well as for quantum corrections to thermodynamic properties. We consider electrons which are interacting with two-level-systems. We show that the two-level-systems enhance the average equilibrium (“persistent”) current in an ensemble of mesoscopic rings. The result supports the recent suggestion that two puzzles in mesoscopic physics may be related: The low temperature saturation of the dephasing time and the high persistent current in rings.
Quantum interference effects play a crucial role in the low temperature properties of normal metals. Prominent examples are weak localization and the associated magnetoresistance. Recently it was suggested that two of the unresolved problems in the physics of mesoscopic metals may have a common solution: The large value of the persistent current in mesoscopic rings and the low temperature saturation of the dephasing rate which is seen in magnetoresistance measurements.
The first problem is the large value of the persistent current in rings. Lévy et al measured the nonlinear response to a magnetic field of an ensemble of $`10^7`$ mesoscopic copper rings. The measured signal corresponds to a current $`II_0\mathrm{sin}(2\pi \varphi /\varphi _0)`$ circulating in each ring. $`\varphi `$ is the magnetic flux which penetrates each ring and $`\varphi _0=h/e`$ is the flux quantum. For temperatures in the mK regime the amplitude was $`|I_0|0.3`$nA per ring, which is of the order of one elementary charge in the time $`\tau _D`$ an electron needs to diffuse around the ring, $`|I_0|0.6e/\tau _D=0.6eE_c/\mathrm{}`$. Here $`E_c=\mathrm{}/\tau _D=\mathrm{}D/L^2`$ is the Thouless energy, $`D`$ is the electron diffusion constant, and $`L`$ is the circumference of the ring. Similar results were reported in Refs..
Theory, when neglecting electron interactions, predicts a current that is of the order $`Ie\delta /\mathrm{}`$, where $`\delta `$ is the average distance of single particle levels at the Fermi energy. With the parameters of the experiment, $`\delta /k0.2`$mK and $`E_c/k25`$mK, the current obtained is about two orders of magnitude too small. Electron interactions first seemed to improve the situation. For Coulomb interaction it was found that $`Ie\mu ^{}E_c/\mathrm{}`$, where $`\mu ^{}`$ is a dimensionless number that characterizes the strength of the interaction in the Cooper channel. However estimates of $`\mu ^{}`$ gave a value which is an order of magnitude too small when comparing it with the experiment. Surprisingly an enhancement of the current was also reported in presence of a moderate concentration of magnetic impurities, with $`Ie(E_c/\mathrm{})\mathrm{min}(\mathrm{}/\tau _s,E_c)/kT`$, where $`1/\tau _s`$ is the spin-flip scattering rate. This can become larger than the current coming from the Coulomb interaction, however since the temperature dependence is different from the one observed this mechanism has not been considered as a possible explanation of the experiment in Ref..
The second problem concerns the phase coherence of the electrons. Whereas it is expected that the dephasing rate goes to zero in the zero temperature limit many experiments show a saturation at low temperature. Usually this saturation is attributed to the presence of magnetic impurities or to heating. However, recently a saturation of the dephasing time has been observed, also after excluding these possibilities. Several attempts have been made to explain the low temperature saturation of the dephasing time . It has been argued by Altshuler et al that non-equilibrium electromagnetic noise can decohere the electrons without heating them. Originally, this non-equilibrium noise was suggested to be due to external radiation which couples into the samples. On the other hand dephasing could also occur due to internal noise. In this case a saturation of the dephasing time could also occur in equilibrium. Experimental evidence is in favor of an internal dephasing mechanism , however it is open if equilibrium or non-equilibrium processes dominate.
Recently Kravtsov and Altshuler have extended earlier work on the effect of a high frequency electromagnetic field in mesoscopic rings and have shown that non-equilibrium noise leads to a directed non-equilibrium current. They then suggested that both the “large” currents observed in and the strong dephasing are related and non-equilibrium phenomena.
In this paper we demonstrate that also for the system in thermal equilibrium an enhanced persistent current is expected if there is an additional electron interaction which gives also rise to strong dephasing. For the particular model involving two-level-systems (TLS) we find (1) a diamagnetic current in the low magnetic field limit (2) a temperature dependence which is close to the experimentally observed one (3) an amplitude which depends on the concentration of TLS. In the following we first recall some of the theoretical concepts concerning the persistent currents. We then estimate the persistent current coming from TLS and, finally, relate the persistent current amplitude and the dephasing rate.
The equilibrium current in a ring which is penetrated by a magnetic flux $`\varphi `$ is calculated by taking the derivative of the thermodynamic potential, $`I(\varphi )=\frac{}{\varphi }\mathrm{\Omega }(\mu ,\varphi ).`$ For simplicity we do not discuss the subtle questions concerning differences between the canonical $`F(N,\varphi )`$ and the grand canonical thermodynamical potential $`\mathrm{\Omega }(\mu ,\varphi )`$. In an ensemble of weakly disordered rings the disorder configuration will change from ring to ring, so in order to calculate the average persistent current of an ensemble of rings one has to average over disorder, $`I(\varphi )I(\varphi )_{\mathrm{dis}}`$. For non-interacting electrons the grand canonical potential, and therefore the persistent current, depends only exponentially weak on the magnetic flux and one finds only a small persistent current.
The situation changes in presence of interactions. As an example take the Coulomb interaction as in Ref. and consider the classical expression for the Coulomb energy,
$$H=\frac{1}{2}d𝐫d𝐫^{}v(𝐫𝐫^{})\delta n(𝐫,\varphi )\delta n(𝐫^{},\varphi ).$$
(1)
This quantity depends on the magnetic flux $`\varphi `$ even on average since the fluctuations of the electron density are magnetic flux dependent and may be written as
$$\delta n(𝐫,\varphi )\delta n(𝐫^{},\varphi )_{\mathrm{dis}}=\underset{m}{}A_m\mathrm{cos}(4\pi m\varphi /\varphi _0)$$
(2)
with
$$A_m=\frac{4N(ϵ_F)}{𝒱}\frac{\mathrm{sin}^2(k_F|𝐫𝐫^{}|)}{(k_F|𝐫𝐫^{}|)^2}kT\underset{\omega >0}{}\sqrt{\frac{\omega }{E_c}}\mathrm{e}^{m\sqrt{\omega /E_c}},$$
(3)
compare Ref.. Here $`\omega =2\pi nkT`$ are bosonic frequencies and $`𝒱`$ is the volume and $`N(ϵ_F)`$ is the density of states of the Fermi level. As a technical remark we would mention that Eq.(3) is obtained by evaluating the diagram with one particle-particle propagator (cooperon). The harmonics of the persistent current $`I(\varphi )=_mI_m\mathrm{sin}(4\pi \varphi m/\varphi _0)`$ are finally found as $`I_m=16\mu _0e/\pi m^2\tau _D`$ with $`\mu _0=N(ϵ_F)d𝐫v(𝐫)\mathrm{sin}^2(k_Fr)/(k_Fr)^2`$. Including the exchange energy reduces the current by a factor two, and higher orders in the interaction reduce the interaction amplitude, $`\mu _0\mu ^{}\mu _0/[1+\mu _0\mathrm{ln}(ϵ_F\tau _D)]`$.
When opening an additional interaction channel one will find an additional contribution to the persistent current. In Ref. this has been demonstrated for magnetic impurities. Here we consider the interaction of conduction electrons with nonmagnetic impurities, which we assume to couple to the electron density. The Hamiltonian is of the form
$$\widehat{H}_{\mathrm{int}}=dx\widehat{n}(𝐱)\widehat{V}(𝐱).$$
(4)
The operator $`\widehat{V}(𝐱)`$ that is due to the impurities will be specified more explicitly below. To second order in this interaction one finds a correction to the free energy which is the sum of a Hartree and a Fock like term, $`\delta \mathrm{\Omega }=\delta \mathrm{\Omega }_\mathrm{H}+\delta \mathrm{\Omega }_\mathrm{F}`$, which are given by ($`\beta =1/kT`$)
$`\delta \mathrm{\Omega }_\mathrm{H}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _0^\beta }d\tau {\displaystyle d𝐱d𝐱^{}\widehat{n}(𝐱)\widehat{n}(𝐱^{})}`$ (6)
$`\times \left[\widehat{V}(𝐱,\tau )\widehat{V}(𝐱^{},0)\widehat{V}(𝐱)\widehat{V}(𝐱^{})\right]`$
$`\delta \mathrm{\Omega }_\mathrm{F}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{s,s^{}}{}}{\displaystyle _0^\beta }d\tau {\displaystyle d𝐱d𝐱^{}\mathrm{\Psi }_s^{}(𝐱,\tau )\mathrm{\Psi }_s^{}(𝐱^{},0)}`$ (8)
$`\times \mathrm{\Psi }_s(𝐱,\tau )\mathrm{\Psi }_s^{}^{}(𝐱^{},0)\widehat{V}(\mathrm{x},\tau )\widehat{\mathrm{V}}(\mathrm{x}^{},0),`$
where $`\mathrm{\Psi }_s^{}(𝐱,\tau )`$ and $`\mathrm{\Psi }_s(𝐱,\tau )`$ are operators for fermions with spin $`s`$ and the brackets $`\mathrm{}`$ are the thermal average . If $`\widehat{V}(𝐱)`$ describes pure potential scattering, then $`\widehat{V}(𝐱)`$ is a c-number with the result that $`\delta \mathrm{\Omega }_\mathrm{H}=0`$. $`\delta \mathrm{\Omega }_\mathrm{F}0`$ but does not depend on the magnetic flux which can be traced back to the fact that $`\widehat{V}(𝐱,\tau )=\widehat{V}(𝐱)`$ is static. This can become different if the impurity has an internal degree of freedom. Consider a TLS, realized by an impurity which sits in a double well potential with minima at $`𝐫`$ and $`𝐫+𝐝`$ which are nearly degenerate in energy. We write the scattering potential as $`\widehat{V}(𝐱)=V[\widehat{n}_A\delta (𝐱𝐫)+\widehat{n}_B\delta (𝐱𝐫𝐝)]`$. $`\widehat{n}_A`$ and $`\widehat{n}_B`$ are the number operators for the impurity in the relevant potential minimum. Since the impurity is in either of these minima $`\widehat{n}_A+\widehat{n}_B=1`$. We further characterize the impurity by the asymmetry $`ϵ`$ and a tunneling amplitude $`\mathrm{\Delta }`$ between between the two minima, so the impurity Hamiltonian is
$$\widehat{H}_{\mathrm{imp}}=\left(\begin{array}{cc}ϵ& \mathrm{\Delta }\\ \mathrm{\Delta }& ϵ\end{array}\right).$$
(9)
The Hartree energy (6), which is nonzero in this model, may be interpreted from the point of view of both the electrons and the impurities. From the electronic point of view the electron impurity interaction gives rise to an effective electron-electron interaction: Comparing Eqs.(1) and (6) one realizes that the Coulomb interaction is replaced by an effective interaction
$`v(𝐱𝐱^{}){\displaystyle _0^\beta }d\tau \left\{\widehat{V}(𝐱,\tau )\widehat{V}(𝐱^{},0)\widehat{V}(𝐱)\widehat{V}(𝐱^{})\right\}`$ (10)
due to the defects. From the impurity point of view the coupling to the conduction electrons changes the level asymmetry,
$$\left(\begin{array}{cc}ϵ& \mathrm{\Delta }\\ \mathrm{\Delta }& ϵ\end{array}\right)\left(\begin{array}{cc}ϵ+V\widehat{n}(𝐫)& \mathrm{\Delta }\\ \mathrm{\Delta }& ϵ+V\widehat{n}(𝐫+𝐝)\end{array}\right),$$
(11)
which then changes the free energy as given in Eq.(6) to second order in $`V`$.
We discuss the persistent current first in the most simple situation, where we neglect the tunnel splitting $`\mathrm{\Delta }`$. In this case $`\widehat{V}(𝐱,\tau )`$ is static so $`\delta \mathrm{\Omega }_\mathrm{F}`$ remains independent of magnetic flux, as in the case of “normal” disorder. Here and below we will therefore concentrate on the Hartree energy. Using the relation $`\widehat{n}_A+\widehat{n}_B=1`$ and averaging over “normal” disorder we can rewrite the Hartree energy as
$`\delta \mathrm{\Omega }_\mathrm{H}_{\mathrm{dis}}`$ $`=`$ $`|V|^2\delta n^2(𝐫,\varphi )_{\mathrm{dis}}\left(1{\displaystyle \frac{\mathrm{sin}^2(k_Fd)}{(k_Fd)^2}}\right)`$ (13)
$`\times {\displaystyle _0^\beta }\mathrm{d}\tau \{\widehat{n}_A(\tau )\widehat{n}_A(0)\widehat{n}_A\widehat{n}_A\}.`$
If the TLS asymmetry is large, $`|ϵ|>kT`$, then $`\widehat{n}_A(\tau )\widehat{n}_A(0)\widehat{n}_A\widehat{n}_A=0`$ and therefore $`\delta \mathrm{\Omega }_\mathrm{H}_{\mathrm{dis}}=0`$. For a TLS with a small asymmetry, $`|ϵ|<kT`$ one finds $`\widehat{n}_A(\tau )\widehat{n}_A(0)\widehat{n}_A\widehat{n}_A=1/4`$ so that $`\delta \mathrm{\Omega }_\mathrm{H}_{\mathrm{dis}}0`$ and a persistent current results. From the integration over $`\tau `$ it follows that the current coming from a single defect is proportional to the inverse temperature, in full analogy to the persistent current from a magnetic impurity. For the system with a finite density of TLS the asymmetry $`ϵ`$ will not be a constant, instead there will be a distribution of asymmetries. Using eq.(2) and below we determine the persistent current as
$$I\frac{8}{\pi }\frac{c_{\mathrm{act}}N(ϵ_F)V^2F}{kT}\frac{e}{\tau _D},$$
(14)
where $`F=1\mathrm{sin}^2(k_Fd)/(k_Fd)^2`$ and $`c_{\mathrm{act}}`$ is the concentration of TLS with $`ϵ<kT`$ and therefore is active in producing a persistent current. Assuming a flat distribution of asymmetries between zero and $`ϵ_{\mathrm{max}}>kT`$, the concentration of active TLS is proportional to the temperature, $`c_{\mathrm{act}}=ckT/ϵ_{\mathrm{max}}`$, which then cancels the inverse temperature dependence of the persistent current of a single defect. The current is diamagnetic in contrast to the paramagnetic current from the repulsive Coulomb interaction. The amplitude of the current is of the diffusive scale, $`Ie/\tau _D`$, as for the Coulomb interaction. The dimensionless prefactor $`\mu ^{}`$ is replaced by the factor $`\mu _{\mathrm{TLS}}=cFN(ϵ_F)V^2/ϵ_{\mathrm{max}}`$. which should be of order one if this mechanism is relevant for the currents observed in Ref.. Assuming an atomic scattering cross section of the TLS and the factor $`F1`$ this requires a density of states of TLS that is comparable to the density of states of the metallic host and therefore of the order $`10^{18}/\mathrm{Kcm}^3`$. At 100mK this corresponds to a concentration of active two-level-systems of about 2ppm which is not a small number but, in principle not impossible . For the assumed distribution of asymmetries the temperature dependence of the persistent current is only due to the temperature dependence of the local density fluctuations, see eq.(3), and is therefore identical to the temperature dependence of the persistent current from the Coulomb interaction. The latter has been shown to agree well with the experiment in Ref.. Finally it is important to discuss spin-orbit scattering, since in the gold or copper rings in the experiments the spin-orbit rate is large. Following Refs. we find that spin-orbit scattering reduces the persistent current due to the mechanism discussed here by a factor four, but the sign remains diamagnetic.
Let us now allow a finite tunnel splitting $`\mathrm{\Delta }`$, i.e. spontaneous transitions of the impurity between the two minima. The correlation function that is relevant for the persistent current, i.e. the impurity susceptibility, is given by
$$_0^\beta d\tau \widehat{n}_A(\tau )\widehat{n}_A(0)\widehat{n}_A\widehat{n}_A=\{\begin{array}{cc}\frac{1}{4}\frac{1}{kT}\hfill & \\ \frac{1}{4}\frac{\mathrm{\Delta }^2}{ϵ^2+\mathrm{\Delta }^2}\frac{1}{\sqrt{ϵ^2+\mathrm{\Delta }^2}}\hfill & \end{array},$$
(15)
in the two limits where $`ϵ^2+\mathrm{\Delta }^2<(kT)^2`$ and $`ϵ^2+\mathrm{\Delta }^2>(kT)^2`$. Whereas for static defects with $`\mathrm{\Delta }=0`$ the correlation function is non-zero only in the high temperature limit, $`kT>ϵ`$, the correlation function for dynamic defects is non-zero even in the zero temperature limit, so these defects contribute to the persistent current even for $`T0`$. We calculate the persistent current under the assumption of a flat distribution of $`ϵ`$ between zero and $`ϵ_{\mathrm{max}}`$ and a distribution of $`\mathrm{\Delta }`$ that is proportional to $`1/\mathrm{\Delta }`$ between $`\mathrm{\Delta }_{\mathrm{min}}`$ and $`\mathrm{\Delta }_{\mathrm{max}}`$. We then find $`I(e/\tau _D)F\mathrm{}/(\tau _{\mathrm{TLS}}ϵ_{\mathrm{max}})`$ as before when we neglected the tunnel splitting. $`\mathrm{}/\tau _{\mathrm{TLS}}cN(ϵ_F)V^2`$ is the electron scattering rate off the TLS.
Finally we discuss the relation of the persistent current and dephasing. In Ref. it has been demonstrated that TLS lead to dephasing with a rate that is temperature independent in a certain range of temperature. Both the persistent current amplitude and the dephasing rate are hard to estimate for a given material since they depend on the concentration of TLS and the distribution of $`ϵ`$ and $`\mathrm{\Delta }`$. It is therefore of interest to relate the two quantities, in order to reduce the number of unknown parameters. Notice that in order to have dephasing there must be real transitions between two impurity states, and one finds that the defects with $`kT>\sqrt{ϵ^2+\mathrm{\Delta }^2}>\mathrm{}/\tau _\varphi `$ are most effective for dephasing. On the other hand all defects with $`kT>\sqrt{ϵ^2+\mathrm{\Delta }^2}`$ and even some with $`kT<\sqrt{ϵ^2+\mathrm{\Delta }^2}`$ contribute to the persistent current, see Eq.(15). We cannot therefore give a general relation between dephasing rate and persistent current amplitude. We can, however, as shown below, give such a relation for our special choice of the distribution of $`ϵ`$ and $`\mathrm{\Delta }`$. The dephasing rate has been estimated as
$$\frac{1}{\tau _\varphi }\{\begin{array}{cc}\mathrm{\Delta }_{\mathrm{max}}F/(ϵ_{\mathrm{max}}\tau _{\mathrm{TLS}}\lambda )\hfill & \mathrm{if}\mathrm{}/\tau _\varphi <\mathrm{\Delta }_{\mathrm{max}}<kT\hfill \\ \mathrm{\Delta }_{\mathrm{max}}(F/\mathrm{}\lambda ϵ_{\mathrm{max}}\tau _{\mathrm{TLS}})^{1/2}\hfill & \mathrm{if}\mathrm{\Delta }_{\mathrm{max}}<\mathrm{}/\tau _\varphi <kT\hfill \end{array}$$
(16)
with $`\lambda =\mathrm{ln}(\mathrm{\Delta }_{\mathrm{max}}/\mathrm{\Delta }_{\mathrm{min}})`$. The persistent current amplitude, $`I\mu _{\mathrm{TLS}}(e/\tau _D)`$ with $`|\mu _{\mathrm{TLS}}|F\mathrm{}/(ϵ_{\mathrm{max}}\tau _{\mathrm{TLS}})`$, is therefore of the order
$$|\mu _{\mathrm{TLS}}|\{\begin{array}{c}\lambda (\mathrm{}/\tau _\varphi )/\mathrm{\Delta }_{\mathrm{max}}\hfill \\ \lambda (\mathrm{}/\tau _\varphi )^2/\mathrm{\Delta }_{\mathrm{max}}^2\hfill \end{array}$$
(17)
in the two limits considered. For example for the gold sample of Ref. $`\mathrm{}/\tau _\varphi 2`$mK below $`500`$mK. If the constant dephasing rate is from the mechanism we consider, then the lowest measured temperature ($`40`$mK) is an upper limit for $`\mathrm{\Delta }_{\mathrm{max}}`$, and leads to the estimate $`|\mu _{\mathrm{TLS}}|>\lambda /20`$.
The dephasing rate for low temperature, $`\mathrm{}/\tau _\varphi <kT<\mathrm{\Delta }_{\mathrm{max}}`$, is proportional to $`T`$ and given by $`1/\tau _\varphi FkT/(ϵ_{\mathrm{max}}\tau _{\mathrm{TLS}}\lambda )`$. Here one finds $`|\mu _{\mathrm{TLS}}|\lambda (\mathrm{}/\tau _\varphi )/kT`$, which depends only on one unknown parameter, $`\lambda `$. A dephasing rate which is linear in $`T`$ has been observed in various three-dimensional and two-dimensional samples. The values which were reported for $`\tau _\varphi `$ at $`10`$K are around $`\tau _\varphi 10^{12}`$sec – $`510^{10}`$sec, which corresponds to $`(\mathrm{}/\tau _\varphi )/kT210^2`$$`1`$. Also from these considerations it seems rather reasonable that the parameter $`|\mu _{\mathrm{TLS}}|`$ can reach values of order one.
In this paper we estimate the persistent current linear in the concentration of TLS and we neglect Kondo correlations. Kondo physics has been suggested as a possible solution of the dephasing problem in Ref.. The persistent current, of course, will be modified by Kondo correlations, however it is beyond the scope of this paper to estimate this quantitatively. We also do not attempt to give an exhaustive discussion of the limit of high concentration of impurities. However it is clear that our theory will overestimate the current when the phase coherence time $`\tau _\varphi `$ becomes of the order of, or shorter than, the diffusive time $`\tau _D`$. The related problem for the persistent current coming from magnetic impurities has been discussed in Ref..
In summary we demonstrated that the interaction of conduction electrons with impurities induces a persistent current. Under reasonable assumptions we find a temperature dependence that is set by the diffusive scale. The most crucial point however is the current amplitude here given by $`I|\mu _{\mathrm{TLS}}|e/\tau _D`$. The dimensionless parameter $`\mu _{\mathrm{TLS}}`$ depends on the concentration of the TLS, so a reliable estimate of the current amplitude is difficult. Experimentally the interrelationship of dephasing and persistent current may be checked by measuring the persistent current for different materials. For silver, where no saturation of the dephasing time has been observed, we expect a smaller persistent current than in gold or copper where the dephasing time saturates at low temperature. The sign of the current may help to decide if non-equilibrium fluctuation suggested in Ref. or the equilibrium electron-impurity interactions studied here dominate the current: For a system with strong spin-orbit interactions Ref. predicts a paramagnetic current, whereas we found a diamagnetic current.
We acknowledge stimulating discussions with U. Eckern and financial support by the DFG through SFB 484 and Forschergruppe HO/955.
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# 1 Introduction
## 1 Introduction
Recent works in string theory have proposed some models as non-perturbative formulations. Especially, the IIB matrix model has been considered as a constructive definition of the type IIB superstring theory. This model is zero-volume limit of the ten-dimensional large $`N`$ supersymmetric Yang-Mills theory and defined by the following action,
$$S=\frac{1}{g^2}\mathrm{tr}\left(\frac{1}{4}[A_\mu ,A_\nu ]^2+\frac{1}{2}\overline{\psi }\mathrm{\Gamma }^\mu [A_\mu ,\psi ]\right),$$
(1)
where $`A_\mu `$ and $`\psi `$ are $`N\times N`$ traceless Hermitian matrices. The interesting feature of this model is that the space-time coordinates are considered as the eigenvalues of these matrices. Then, we expect that the fundamental issues including the dimensionality and the quantum gravity can be understood by studying the dynamics of the model.
To take the continuum limit ($`N\mathrm{}`$) for the IIB Matrix model, a sensible double scaling limit should be determined dynamically. The scaling property of the model has been studied with the light-cone string field Hamiltonian of the type IIB superstring theory. The scaling property of the two important quantities of the model, the string scale ($`\alpha ^{}`$) and the string coupling constant ($`g_{str}`$) are determined as follows,
$`\alpha ^{}`$ $``$ $`g^2N^{a+b}=g^2N^\gamma Constant,`$
$`g_{str}`$ $``$ $`N^{ab}.`$ (2)
For the finite value of the string coupling constant ($`g_{str}`$), one have a restriction, $`a=b`$. In the IIB matrix model, the exponent ($`\gamma `$) plays an important role that one can take the large $`N`$ limit as the continuum limit for the IIB Matrix model. This exponent is determined dynamically.
For studying the dynamical aspects of the model, some Monte Carlo simulations have been performed. In Ref., the existence of the large $`N`$ limit of bosonic $`SU(N)`$ Yang-Mills matrix model for $`D>2`$ has been discussed analytically. Then, in Ref., the bosonic model has been studied with a analytical method, the $`1/D`$ expansion, as well as a numerical one. The numerical results support the $`1/D`$ expansion as an effective tool detecting the large $`N`$ scaling behavior of the model. The leading term of the $`1/D`$ expansion at $`D>3`$ suggests that the exponent ($`\gamma `$) takes a value of 1. We study the model in ten dimensions and reconfirm the expected scaling behavior by a numerical method with a larger size matrix.
In analogy with the two-dimensional model, the Eguchi-Kawai model, we study the scaling property of the Wilson loop in ten-dimensions. The area law of the Wilson loop operator has been found in the four-dimensional model in Ref.. In this article, we also obtain that the area dependence of the Wilson loop obeys the area law in ten dimensions. We calculate the string tension from the area law of the Wilson loop. We thus consider that the scaling property of the string tension is estimated in ten-dimensional model.
This paper is organized as follows. In section 2, we review the model and some perturbative analysis. In section 3, we show the numerical results and the scaling property of the model. Then, we present the data which show the existence of the double scaling limit in the model. Finally, in section 4, we summarize and discuss our numerical results.
## 2 Large $`N`$ behavior of correlation functions
First, let us remind the perturbative arguments of the bosonic model and describe briefly the large $`N`$ behavior of the correlation functions of the gauge fields based on .
The bosonic model of the IIB matrix model is given by
$$S_{bosonic}=\frac{1}{g^2}\mathrm{tr}[A_\mu ,A_\nu ][A^\mu ,A^\nu ],$$
(3)
where $`A_\mu `$ are $`N\times N`$ Hermitian matrices representing the ten-dimensional gauge fields. The coupling constant ($`g`$) is nothing but a scale parameter and is absorbed with the rescaling of the gauge field ($`A_\mu `$) as $`A_\mu \frac{1}{\sqrt{g}}A_\mu `$.
The Schwinger-Dyson equation is given by
$$0=𝑑A\frac{}{A}\left(\mathrm{tr}\left(A_\mu e^{S_{bosonic}}\right)\right),$$
(4)
leading to the relation of the correlation function,
$$<\mathrm{tr}([A_\mu ,A_\nu ]^2)>=D(N^21)g^2.$$
(5)
For the estimation of the large $`N`$ behavior of the correlation function, $`<\mathrm{tr}([A_\mu ,A_\nu ]^2)>`$, the matrices ($`A_\mu `$) are decomposed into the diagonal parts ($`X_\mu ^i`$) and the off diagonal parts ($`\stackrel{~}{A_\mu }^{ij}`$), and eq.(5) is taken up to the second order of the off diagonal elements ($`\stackrel{~}{A_\mu }^{ij}`$),
$$<\mathrm{tr}[A_\mu ,A_\nu ]^2>2<\mathrm{tr}[X_\mu ,\stackrel{~}{A}_\nu ][X_\mu ,\stackrel{~}{A}_\nu ]><\mathrm{tr}[X_\mu ,\stackrel{~}{A}_\nu ][X_\nu ,\stackrel{~}{A}_\mu ]>.$$
(6)
Counting the order of the diagonal parts, the large $`N`$ behavior of the leading term of the correlation function counting with the order of the diagonal parts is given by
$$<\mathrm{tr}[A_\mu ,A_\nu ]^2>g^2N^2.$$
(7)
For the finite value of the correlation function, an upper limit of the typical scale of the extend of the space-time, for example can be $`R^2=<\mathrm{tr}A^2>`$, is suggested as $`R^2N^{1/2}`$ perturbatively and from the $`1/D`$ expansion the large $`N`$ behavior is shown as
$$R^2gN^{1/2}.$$
(8)
In the similar manner the correlation functions can be calculated perturbatively.
Next, we consider the Wilson loop operator in the IIB matrix model. The Wilson loop operator and the large $`N`$ behavior have been studied with the light-cone string field theory of the type IIB superstring. The Wilson loop operator ($`w(C)`$) is defined as,
$$w(C)=\mathrm{tr}(v(C)),$$
(9)
where $`C`$ denotes the closed path and $`v(C)`$ is defined as $`v(C)=U_\mu \mathrm{}U_\mu =P_C\mathrm{exp}(i𝑑\sigma k_\mu A_\mu )`$ in the bosonic model. The matrices ($`U_\mu `$) are considered as the unitary matrices,
$$U_\mu =\mathrm{exp}(i𝑑lA_\mu ).$$
(10)
In the ordinary lattice gauge field theory, the expectation value of the Wilson loop operator which spreads a large area behaves as follows,
$$w(I,J)\mathrm{exp}(KI\times J),$$
(11)
where $`I`$ and $`J`$ are the side lengths of the rectangular loop and $`K`$ denotes the string tension. In the same analogy, we study the Wilson loop operator of the bosonic model of the IIB matrix model.
From the scaling relation of two-dimensional Eguchi-Kawai model,
$$g^2NConstant,$$
(12)
It is expected that the similar scaling relation also holds in the IIB matrix model as
$$\alpha ^{}g^2N^\gamma Constant.$$
(13)
The exponent ($`\gamma `$) should be determined dynamically from the model. For the large $`N`$ limit, the parameter ($`g^2N^\gamma `$) must be fixed in the IIB matrix model in the same manners as the parameter ($`g^2N`$) must be fixed in the Eguchi-Kawai model.
We notice that the bosonic model is equivalent to the $`D>2`$ Eguchi-Kawai model in the weak coupling limit. Since the $`U(1)^D`$ symmetry rotates all the eigenvalues by the same angle, the following expansion is valid in the weak coupling region,
$$U_\mu e^{i\alpha _\mu }e^{iA_\mu },$$
(14)
where $`\alpha _\mu `$ take constant values due to the $`U(1)^D`$ symmetry and $`A_\mu `$ are small. The bosonic model action can be obtained by expanding the action of the ten-dimensional Eguchi-Kawai model in terms of $`A_\mu `$. When the higer order terms of $`A_\mu `$ can be neglected, we can obtain the area law of the Wilson loop operator as
$$w(I\times J)=<\mathrm{tr}(U_\mu \mathrm{}U_\mu )_{I\times J}><\mathrm{tr}(e^{iA_\mu }\mathrm{}e^{iA_\mu })_{I\times J}>\mathrm{exp}(K(I\times J)).$$
(15)
In Ref., the area dependence of the Wilson loop operator has been measured and found the area law in $`D=4`$.
In following section, we will show that the area law holds in the bosonic model of the IIB matrix model in $`D=10`$.
## 3 Monte Carlo simulation and Large N behavior of the bosonic model
For numerical simulation of the model, we consider the partition function,
$$Z=𝑑Ae^{S_{bosonic}},$$
(16)
where the action $`S_{bosonic}`$ is given by eq.(3) and the measure of gauge fields is defined by
$$dA=\underset{\mu =1}{\overset{10}{}}\left[\underset{i=1}{\overset{N}{}}\underset{j=i}{\overset{N}{}}d(A_\mu ^{ij})\right].$$
(17)
The action is quadratic with respect to each component, which means that we can update each component by generating gaussian random number in the heat-bath and the Metropolis algorithm.
To confirm the large $`N`$ behavior of the model, we measure the following expectation values,
$`R^2=<{\displaystyle \frac{1}{N}}\mathrm{tr}(A^2)>,`$
$`<{\displaystyle \frac{1}{N}}\mathrm{tr}([A_\mu ,A_\nu ]^2)>.`$ (18)
From the perturbative calculation, the large $`N`$ behaviors are shown as
$`R^2gN^{1/2},`$
$`<{\displaystyle \frac{1}{N}}\mathrm{tr}([A_\mu ,A_\nu ]^2)>g^2N.`$ (19)
In Fig.1 and Fig.2, we show the numerical results of the extent of space-time ($`R^2`$) and the correlation function ($`<\mathrm{tr}([A_\mu ,A_\nu ])^2>`$) for $`D=10`$ with $`N=16,32,48,64,128`$, respectively. We obtain the large $`N`$ behavior as
$`R^2gN^{0.5(1)},`$
$`<{\displaystyle \frac{1}{N}}\mathrm{tr}([A_\mu ,A_\nu ]^2)>g^2N^{1.0(1))}.`$ (20)
By the simulation using the larger size matrix, the numerical result get close to the expected results.
Then, we consider the Wilson loop operator,
$$w(C)=<\frac{1}{N}\mathrm{tr}(e^{iA_\mu }\mathrm{}e^{iA_\nu })>.$$
(21)
We take the loop ($`C`$) as the rectangular ($`I\times J`$) where we select any two direction ($`\mu ,\nu `$) in ten dimensions.
We show the measurement results of the loop operator in Fig.3. Since the dependence of the direction of the rectangular is not found, we consider that in the bosonic model the isotropy of the ten-dimensional space-time is not broken down spontaneously.
Then, we calculate the large $`N`$ behavior by the Wilson loop operator. From the numerical results, the Wilson loop operator, $`<\frac{1}{N}\mathrm{tr}(e^{iA_\mu }\mathrm{}e^{iA_\nu })_{I\times J}>`$, closes to the exponential curve with the large size area. It means that the Wilson loop operator obeys the area law eq.(15).
Then, we can obtain the string tension. In Fig.4, we plot the string tension ($`K`$).
We find the large $`N`$ behavior of the string tension ($`K`$),
$$Kg^2N^{1.07(1)}.$$
(22)
From the string tension ($`K`$), the string scale ($`\alpha ^{}`$) is estimated as
$$\alpha ^{}1/Kg^2N^{1.07(1)}=Constant.$$
We remark that the string field theory of the light-cone frame suggests
$$\alpha ^{}g^2N=Constant.$$
(23)
We consider that the numerical result closes to the analytic one and that the four-dimensional model has the same scaling property.
Furthermore the numerical result suggests that the large $`N`$ behavior of the square root of the string tension approximately equal to the inverse of the extent of the space-time.
$$K^{1/2}g^1N^{0.54(1)}R^1.$$
(24)
It means that the Planck scale of the theory has the same scaling property of the extent of the space time in ten-dimensions. It also holds on the two-dimensional model and the four-dimensional model.
## 4 Summary and Discussion
Let us summarize the main points made in our calculation. We confirm that the Wilson loop operator in the ten-dimensional bosonic model obeys the area law similar to the two and four-dimensional model. For the scaling behavior of the bosonic model of the IIB matrix model, our numerical estimation is
$$\alpha ^{}g^2N^{1.07(1)}=Constant.$$
(25)
Our results show the ten-dimensional space-time extends with $`g^{1/2}N^{0.25(10)}`$ and the extent scale of the space-time approximately equal to the Planck scale ($`l`$), $`lK^{1/2}`$. Furthermore, these results are consistent to the suggestion from the string field theory on light-cone frame and the $`1/D`$ expansion. From the numerical results, we consider that the ten-dimensional bosonic model and the four-dimensional model have the same scaling property.
In this article, we calculate only the bosonic model of the IIB matrix model. For the future work, we are also considering the numerical simulation of the full model including the fermionic term. The supersymmetric four-dimensional model has been studied in Ref.. Optimistically, we expect that we can simplify the fermionic term with the perturbative calculation. In Ref., it is claimed that the model with the 1-loop effective action of the IIB matrix model produces the four-dimensional space-time from the ten-dimensional space-time. We thus make preparations the calculation of the modified model including the fermionic term. The improved supersymmetric model including the fermionic term is studied.
Acknowledgements
We would like to thank T.Yukawa, N.Ishibashi, Y.Kitazawa and H.Kawai. Furthermore, we are grateful to F.Sugino, N.Tsuda, S.Oda and especially J.Nishimura for fruitful discussions and advice. We are also grateful to the members of the KEK theory group. Numerical calculations were performed using the NEC SX4 (Tokai University) and the originally designed cluster computer for quantum gravity and strings, CCGS (KEK).
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# Untitled Document
PUPT-1932 ITEP-TH-23/00 NSF-ITP-00-43 hep-th/0005204
MONOPOLES AND STRINGS IN NONCOMMUTATIVE GAUGE THEORY
David J. Gross <sup>1</sup>, Nikita A. Nekrasov <sup>2</sup>
<sup>1</sup> Institute for Theoretical Physics, University of California Santa Barbara CA 93106
<sup>2</sup> Institute for Theoretical and Experimental Physics, 117259 Moscow, Russia
<sup>2</sup> Joseph Henry Laboratories, Princeton University, Princeton, New Jersey 08544
e-mail: gross@itp.ucsb.edu, nikita@feynman.princeton.edu
Abstract
We study some non-perturbative aspects of noncommutative gauge theories. We find analytic solutions of the equations of motion, for noncommutative U(1) gauge theory, that describe magnetic monopoles with a finite tension string attached. These solutions are non-singular, finite and sourceless. We identify the string with the projection of a D-string ending on a D3-brane in the presence of a constant $`B`$-field.
05/00
1. Introduction
Recently there has been a revival of interest in field theories on noncommutative spaces , especially those that emerge as various limits of M theory compactifications . The latest circumstances in which such theories were found involve D-branes in the presence of a background Neveu-Schwarz $`B`$-field . The interest in such theories is motivated by many analogies between noncommutative gauge theories and large $`N`$ ordinary non-abelian gauge theories , and also by the many features that noncommutative field theories share with open string theory .
In this paper we study some non-perturbative dynamical objects in noncommutative gauge theory, specifically four dimensional gauge theory with an adjoint Higgs field. The theory depends on a dimensionfull parameter $`\theta `$ which enters the commutation relation between the coordinates of the space: $`[x,x]i\theta `$. We treat only the bosonic fields, but these could be a part of a supersymmetric multiplet, with $`𝒩=2`$ supersymmetry or higher. Such field theories arise on the world volume of D3-branes in the presence of a background constant $`B`$-field along the D3-brane.
A D3-brane can be surrounded by other branes as well. For example, in the Euclidean setup, a D-instanton could approach the D3-brane. In fact, unless the D-instanton is dissolved inside the brane the combined system breaks supersymmetry . The D3-D(-1) system can be rather simply described in terms of a noncommutative $`U(1)`$ gauge theory - the latter has instanton-like solutions . However, it turns out that the “topology” of the combined system is non-trivial (despite the fact that the notion of a “point” on a noncommutative space makes very little sense, the non-triviality of topology can be detected), and it is this topology that supports the instantons .
Another, perhaps even more interesting situation, is that of a D-string that ends on a D3-brane. The endpoint of the D-string is a magnetic charge for the gauge field on the D3-brane. In the commutative case, in the absence of the $`B`$-field, the D-string is a straight line, orthogonal to the D3-brane. It projects onto the D3-brane in the form of a singular source, located at the point where the D-string touches the D3-brane. From the point of view of the D3-brane this is a Dirac monopole, with energy density that diverges at the origin.
The situation changes drastically when the $`B`$-field is turned on. One can trade a constant background $`B`$-field with spatial components for a constant background magnetic field. The latter pulls the magnetic monopoles with the constant force. As a consequence, the D-string bends , in order for its tension to compensate the magnetic force. It projects to the D3-brane as a half-line with finite tension. It is quite fascinating to see, as we shall explicitly verify, that the shadow of this string is seen by the noncommutative gauge theory. The $`U(1)`$ noncommutative gauge theory has a monopole solution, that is everywhere non-singular, and whose energy density localizes along a half-line, making up a semi-infinite string. We should stress that the non-singularity of the solution is the non-perturbative in $`\theta `$ property, it couldn’t be seen by the expansion in $`\theta `$ around the Dirac monopole .
The fact that all the fields involved are non-singular, and that the solution is in fact a solution to the noncommutative version of the Bogomolny equations everywhere, makes us suspect that the string in the monopole solution is an intrinsic object of the gauge theory. As such, one could expect that the noncommutative gauge theory holographically describes strings as well. This statement is further supported by the fact that in the limit of very large $`B`$-field (the limit which must be well described by the non-commutative gauge theory ) the D-string almost lies on top of the D3-brane, practically dissolving in it.
Finally, by applying S-duality one could map the solution we found into the solution describing the electric flux tube, represented by the fundamental string . In this way one may hope to arrive at the description of the confining strings in the noncommutatve Yang-Mills theories. Notice however, that the S-duality maps the theory with the spatial noncommutativity to that of the temporal noncommutativity, with all its surprises , in addition to the strong coupling .
The outline of this paper is as follows. In Section 2 we consider some general features of noncommutative field theory, and discuss how it is convenient to work in the Fock space in which the coordinates are expresses as creation and annihilation operators. In Section 3 we construct the Green’s function of the Laplace operator on noncommutative spaces, which illustrates the smearing of space induced by the noncommutativity of the coordinates. We also give a brief introduction to noncommuative gauge theories.
In Section 4 we set up the equations for BPS solutions of four dimension noncommutative gauge theories. We review Nahm’s construction of commutative monopoles, exhibit the SU(2) monopole, as well as the Dirac monopole in this framework. Section 5 is devoted to the construction of the explicit solution of the BPS equations for the $`U(1)`$ noncommutative gauge theory coupled to a scalar field. The properties of the solution are analyzed in Section 6. We conclude, in Section 7, with a discussion of the implications of our analysis.
Acknowledgements.
We would like to thank D. Bak, S. Cherkis, D. -E. Diaconescu, S. Giddings, A. Hashimoto, K. Hashimoto, N. Itzhaki, I. Klebanov, T. Piatina, A. Polyakov, A. Schwarz, S. Shatashvili, and K. Selivanov for discussions. Our research was partially supported by NSF under the grant PHY94-07194, in addition, research of NN was supported by Robert H. Dicke fellowship from Princeton University, partly by RFFI under grant 00-02-16530, partly by the grant 00-15-96557 for scientific schools. NN is grateful to ITP, UC Santa Barbara, CIT-USC Center, and CGTP at Duke University for their hospitality during various stages of this work.
2. Noncommutative Field Theory
Consider space-time with coordinates $`x^i`$, $`i=1,\mathrm{},d`$ which obey the following commutation relations:
$$[x^i,x^j]=i\theta ^{ij},$$
where $`\theta ^{ij}`$ is a constant asymmetric matrix of rank $`2rd`$. By noncommutative space-time one means the algebra $`𝒜_\theta `$ generated by the $`x^i`$ satisfying (2.1), together with some extra conditions on the allowed expressions of the $`x^i`$. The elements of $`𝒜_\theta `$ can be identified with ordinary functions on $`𝐑^d`$, with the product of two functions $`f`$ and $`g`$ given by the Moyal formula (or star product):
$$fg(x)=\mathrm{exp}\left[\frac{i}{2}\theta ^{ij}\frac{}{x_1^i}\frac{}{x_2^j}\right]f(x_1)g(x_2)|_{x_1=x_2=x}.$$
A field theory is defined as usual by constructing an action, say in the case of a scalar field theory,
$$(\mathrm{\Phi })=d^dx\left[_i\mathrm{\Phi }_i\mathrm{\Phi }+V(\mathrm{\Phi })+\mathrm{}\right].$$
The symbol $`d^dx`$ is a notation for a trace, $`\mathrm{Tr}`$, on the algebra $`𝒜_\theta `$. When one works on compact noncommutative manifolds (compact manifolds, whose algebra of functions is deformed, e.g. by the techniques of ), for example the noncommutative torus, then the trace is the usual trace, i.e. the linear map $`𝒜_\theta \text{ }\mathrm{C}`$, such that $`\mathrm{Tr}[a,b]=0`$. On $`𝐑^d`$ the notion of the trace is trickier, in particular the trace of the commutator may not vanish, just as the integral of a total derivative may not vanish. We shall encounter such effects in our discussion below, so instead of giving formal definitions at this point we shall treat explicit examples later.
The Lagrangian of a field theory involves derivatives. The derivative $`_i`$ is the infinitesimal automorphism of the algebra (2.1):
$$x^ix^i+\epsilon ^i,$$
where $`\epsilon ^i`$ is a $`c`$-number. For the algebra (2.1) this automorphism is internal:
$$_i\mathrm{\Psi }=i\theta _{ij}[\mathrm{\Psi },x^j],$$
where $`\theta _{ij}`$ is the inverse of $`\theta ^{ij}`$, namely $`\theta _{ij}\theta ^{jk}=\delta _i^k`$. In contrast, on the torus generated by $`U_l=\mathrm{exp}2\pi ix^l`$, it is an outer automorphism. This difference is crucial in the analysis of noncommutative gauge theories.
By an orthogonal change of coordinates we can map the Poisson tensor $`\theta _{ij}`$ onto its canonical form:
$$x^iz_a,\overline{z}_a,a=1,\mathrm{},r;y_b,b=1,\mathrm{},d2r,$$
so that:
$$\begin{array}{cc}& [y_a,y_b]=[y_b,z_a]=[y_b,\overline{z}_a]=0,[z_a,\overline{z}_b]=2\theta _a\delta _{ab},\theta _a>0\hfill \\ & ds^2=dx_i^2+dy_b^2=dz_ad\overline{z}_a+dy_b^2.\hfill \end{array}$$
Since $`z(\overline{z})`$ satisfy (up to a constant) the commutation relations of creation (annihilation) operators we can identify functions $`f(x,y)`$ with operator valued functions of the $`y_a`$ in the Fock space of the $`r`$ creation and annihilation operators (the operators in the Fock space are widely used in the studies of noncommutative theories and matrix models, for their applications to the latter see ):
$$\alpha _a=z_a/\sqrt{2\theta _a},\alpha _a^{}=\overline{z}_a/\sqrt{2\theta _a},[\alpha _a,\alpha _b^{}]=\delta _{ab}.$$
Since we shall be dealing with scale invariant theories in which the only scales will be the $`\theta _a`$ we shall set all $`2\theta _a=1`$. When desired, the $`\theta _a`$’s can be introduced by rescaling the coordinates, $`z_az_a/\sqrt{2\theta _a}`$. Let $`\widehat{n}_a=\alpha _a^{}\alpha _a`$ be the $`a`$’th number operator.
The procedure that maps ordinary commutative functions onto operators in the Fock space acted on by $`z_a,\overline{z}_a`$ is called Weyl ordering and is defined by:
$$f\left(x=(Z_a,\overline{Z}_a)\right)\widehat{f}(z_a,\overline{z}_a)=f(x)\frac{d^{2r}xd^{2r}p}{(2\pi )^{2r}}e^{i\left(\overline{p}_a\left(z_aZ_a\right)+p_a\left(\overline{z}_a\overline{Z}_a\right)\right)}.$$
It is easy to see that
$$\mathrm{if}f\widehat{f},g\widehat{g}\mathrm{then}fg\widehat{f}\widehat{g}.$$
A useful formula is for the matrix elements of $`\widehat{f}`$ in the coherent state basis
$$\xi |\widehat{f}|\eta =f(Z,\overline{Z})\frac{d^rZd^r\overline{Z}}{(2\pi i)^{2r}}e^{\xi \eta 2(\xi \overline{Z})(\eta Z)}$$
where $`\xi |`$ and $`|\eta `$ are coherent states: $`\xi |=\mathrm{𝟎}|\mathrm{exp}\left(\xi _az_a\right),|\eta =\mathrm{exp}\left(\eta _a\overline{z}_a\right)|\mathrm{𝟎}.`$ From (2.1) we can extract the matrix elements of $`\widehat{f}`$ between the standard oscillator states by:
$$𝐤|\widehat{f}|𝐥=\frac{1}{\sqrt{𝐤!𝐥!}}_\xi ^𝐤_\eta ^𝐥|_{\xi =\eta =0}\xi |\widehat{f}|\eta ,$$
where $`𝐤`$, $`𝐥`$ are the vectors of the occupation numbers, e.g. $`𝐤=(k_1,\mathrm{},k_r)`$.
Given the operator $`\widehat{f}`$, or its matrix elements, in the coherent state or occupation number basis, one can easily reconstruct the function $`f(x)`$ to which it corresponds. For example, consider the simplest case where $`r=1,d=3`$, that will be our interest below. Furthermore, consider functions that are axially symmetric. This means that $`f(x)=f(r,x_3)`$, where $`r=\sqrt{x_1^2+x_2^2}`$; or equivalently that $`k|\widehat{f}(x_3)|l=\delta _{kl}f_l(x_3)`$. Then to reconstruct the function $`f(x)`$, from the matrix elements $`f_l(x_3)`$ one uses:
$$f(r,x_3)=2\underset{l=0}{\overset{\mathrm{}}{}}(1)^lf_l(x_3)L_l(4r^2)e^{2r^2},$$
where $`L_l(4r^2)`$ are Laguerre polynomials.
3. Scalar Field Green’s functions
An interesting property of noncommutative field theory is its similarity with lattice field theory, namely the noncommutativity of the coordinates introduces a smearing of space. We shall illustrate this similarity by examining the Green’s functions of the Laplace operator on noncommutative space-time.
3.1. Sources
Consider the noncommutative version of the equation for the Green’s function $`\mathrm{\Delta }_xG(x,x^{})=\delta (xx^{})`$. Recalling (2.1), the Laplace operator can be rewritten as follows:
$$\widehat{\mathrm{\Delta }}=\frac{^2}{y_b^2}4\theta _a^2[[,z_a],\overline{z}_a].$$
Thus the noncommutative equation for the Green’s function is
$$\frac{^2}{y_b^2}\widehat{G}(z,\overline{z},y;z^{},\overline{z}^{},y^{})4[[\widehat{G}(z,\overline{z},y;z^{},\overline{z}^{},y^{}),z_a],\overline{z}_a]=\widehat{\delta }(z,\overline{z};z^{},\overline{z}^{})\delta (yy^{}),$$
where we have introduced two copies of the algebra $`𝒜_\theta `$, generated by $`z_a,\overline{z}_a,z_a^{},\overline{z}_a^{}`$ and $`\widehat{\delta },\widehat{G}`$ are operators in the tensor product of two Fock spaces
$$_{1,2}=_1_2$$
spanned by $`|𝐥_1,𝐥_2=|𝐥_1|𝐥_2`$.
The expression for the delta function, $`\widehat{\delta }(z,\overline{z};z^{},\overline{z}^{})`$, is now easy to obtain directly in the coherent state basis , using the (tensor product form of) (2.1). In terms of $`|\eta _1,\eta _2=|\eta _1|\eta _2`$ we have
$$\begin{array}{cc}\hfill \xi _1,\xi _2|\widehat{\delta }|\eta _1,\eta _2=& e^{\xi _1\eta _1+\xi _2\eta _22(\xi _1\overline{Z})(\eta _1Z)2(\xi _2\overline{Z})(\eta _2Z)}d^rZd^r\overline{Z}\hfill \\ \hfill =& e^{\xi _1\eta _2+\xi _2\eta _1}.\hfill \end{array}$$
The matrix elements of $`\delta `$ in the occupation number basis are:
$$𝐤_{1,2}|\widehat{\delta }|𝐥_{1,2}=\delta _{𝐤_1,𝐥_2}\delta _{𝐤_2,𝐥_1}.$$
Thus $`\widehat{\delta }`$ is a permutation operator $`P:`$, $`P(e_1e_2)=e_2e_1`$, and squares to the identity operator $`P^2=Id`$. It is easy to verify that $`\widehat{\delta }`$ satisfies the defining property of the delta-function, namely
$$\mathrm{Tr}_x^{}\left[\widehat{\delta }(x,x^{})\widehat{f}(x^{})\right]=\widehat{f}(x)$$
What is the noncommutative version of a source localized at the origin, namely $`\delta ^{2r}(x)`$? Using (2.1) we see that
$$\xi |\widehat{\delta }|\eta =\delta ^2(Z)\frac{d^rZd^r\overline{Z}}{(2\pi i)^{2r}}e^{\xi \eta 2(\xi \overline{Z})(\eta Z)}=e^{\xi \eta },$$
or in the occupation number basis:
$$𝐤|\widehat{\delta }|𝐥=\delta _{𝐤,𝐥}(1)^{|𝐤|},\widehat{\delta }=(1)^{\widehat{𝐧}}$$
with $`|𝐤|=_ak_a,\widehat{𝐧}=_a\widehat{n}_a`$. In this way the delta function becomes an operator in the Fock space with the spectrum of the form of the diffraction rings. Note that $`(\widehat{\delta }(x))^2=Id`$, which is the transform of the constant function.
Alternatively we can relate (3.1) to (3.1). by passing to the center-of-mass frame:
$$z^c=\frac{1}{\sqrt{2}}\left(z+z^{}\right),z^r=\frac{1}{\sqrt{2}}\left(zz^{}\right)$$
and similarly for $`\overline{z},\overline{z}^{}`$. The expression (3.1) is written in the number basis for the operators $`z^r,\overline{z}^r`$. The transformation (3.1) is a unitary one:
$$\begin{array}{cc}\hfill Sz^rS^{}=z,& Sz^cS^{}=z^{},\hfill \\ \hfill SzS^{}=z^c,& Sz^{}S^{}=z^r\hfill \\ \hfill S=\mathrm{exp}\frac{\pi }{4}& \left(\overline{z}^{}z\overline{z}z^{}\right)\hfill \end{array}$$
It is easy to check that $`SP=PS^{}`$. Therefore $`SPS^{}=S^2P`$. Now, consider $`U=S^2`$. It acts as follows:
$$UzU^{}=z^{},Uz^{}U^{}=z$$
Let us now apply the $`S`$ transformation to the delta function:
$$SPS^{}|𝐥_1|𝐥_2=U|𝐥_2|𝐥_1=(1)^{|𝐥_1|}|𝐥_1|𝐥_2$$
i.e. we get complete agreement with (3.1).
Thus in the noncommutative case we cannot construct a truly localized source. The transform of $`\delta ^{2r}(x)`$, in which the noncommuting coordinates are all localized at the origin, is spread out over all of space. The most localized source we can construct in the noncommutative case is a Gaussian wave packet $`D(x)=\mathrm{exp}(2Z\overline{Z})`$, whose transform is
$$\widehat{D}=|\mathrm{𝟎}\mathrm{𝟎}|,\xi |\widehat{D}|\eta =1,𝐤|\widehat{D}|𝐥=\delta _{𝐤,𝐥}\delta _{𝐥,\mathrm{𝟎}}$$
One can also develop the similar analysis for finite lattices, in which case one gets the finite matrix versions of the operators (3.1)(3.1)(see ).
3.2. Green’s functions
We now consider the Laplace equation for the Green’s function, $`\widehat{G}`$. Consider a function (an element of $`𝒜_\theta `$) that commutes with all $`N_a`$’s. In the commutative language this means that the functions we wish to look at are invariant under rotations of the $`z_a,\overline{z}_a`$ two-planes. We take $`\widehat{G}`$ to be such a function. On such functions the Laplacian acts as follows:
$$\mathrm{\Delta }G_𝐧=\frac{^2}{y_b^2}G_𝐧+4\underset{a}{}\left((2n_a+1)G_𝐧+(n_a+1)G_{𝐧+𝐞_a}+n_aG_{𝐧𝐞_a}\right)$$
where
$$𝐧=(n_1,\mathrm{},n_r),𝐞_a=(0,0,\mathrm{},1_{\widehat{a}},\mathrm{},0),𝐧|\widehat{G}|𝐧^{}=\delta _{𝐧^{},𝐧}G_𝐧.$$
The formula (3.1) requires the following comment: when evaluating the right hand side of (3.1) the number operators $`n_a`$ must be evaluated first, so that if some of the $`n_a`$’s vanish the whole expression $`n_aG_{𝐧𝐞_a}`$ must be set to zero, no matter how singular the analytic expression for $`G_{𝐧𝐞_a}`$ may look.
One can also rewrite the Laplacian (3.1) using the finite difference operators: $`𝒟_a,𝒩_a`$:
$$𝒟_aG_𝐧=G_𝐧G_{𝐧𝐞_a},𝒩_aG_𝐧=(n_a+1)G_{𝐧+𝐞_a}$$
$$\mathrm{\Delta }=\frac{^2}{y_b^2}+4𝒟_a𝒩_a𝒟_a.$$
The operators $`𝒩,𝒟`$ form a Heisenberg algebra:
$$[𝒟_a,𝒩_b]=\delta _{ab}.$$
Let us compare the expressions (3.1),(3.1) to their commutative analogues. Let $`y_b,Z_a,\overline{Z}_a`$ denote the coordinates on the commutative space-time with the metric
$$ds^2=dy_b^2+\frac{1}{2}dZ_ad\overline{Z}_a,$$
and consider functions that depend only on $`y_b`$ and $`R_a=|Z_a|^2`$. On such a function, say $`𝒢`$, the Laplacian acts as follows:
$$\mathrm{\Delta }𝒢=\frac{^2}{y_b^2}𝒢+4\underset{a}{}\frac{}{R_a}R_a\frac{}{R_a}𝒢.$$
The operators $`R_a`$ and $`_{R_b}`$ form the same Heisenberg algebra (3.1)as $`𝒩,𝒟`$. Consequently, by mapping the representation (3.1) to the standard representation of the Heisenberg algebra acting on functions of $`R_a`$, we can map the Green’s function of the noncommutative Laplacian to that of the commutative one.
Let us note, however, that the algebra (3.1) is represented by functions on the whole space $`𝐑^r`$, whereas the variables $`R_a`$, by definition, must be positive. The same comment applies to the variables $`n_a0`$. This boundnessness of the domain in the definition of these functions produces the source terms in the Laplace (or other) equations that they obey.
Now construct the Laplace transform of the function $`𝒢`$:
$$\stackrel{~}{𝒢}(t)=_0^{\mathrm{}}\underset{a}{}dR_ae^{t_aR_a}𝒢(R_a).$$
At the same time we construct the generating function associated with $`G_𝐧`$:
$$\widehat{G}(t)=\underset{𝐧}{}\underset{a}{}(1t_a)^{n_a}G_𝐧.$$
It is easy to see that the Laplacian operators acting on both $`\stackrel{~}{𝒢}(t)`$ and on $`\widehat{G}(t)`$, are mapped to the same operator:
$$\widehat{\mathrm{\Delta }}=\frac{}{y_b^2}4t_a\frac{}{t_a}t_a,$$
(as long as we assume that the functions don’t grow too fast at infinity or at zero) In this way we define a map from functions of continuous $`R`$ coordinates to functions of discrete $`𝐧`$:
$$G_𝐧=_0^{\mathrm{}}𝒢_R\underset{a}{}\frac{R_a^{n_a}}{n_a!}dR_ae^{R_a}.$$
Note that as $`𝐧\mathrm{}`$ the saddle point approximation gives
$$G_𝐧𝒢_{R_a=n_a}.$$
We can use this to construct the noncommutative version of the Green’s function. For example, take $`d=2r+1`$, then the $`2r+1`$ \- dimensional Green’s function
$$𝒢(y,R)=\frac{1}{\left(y^2+_aR_a\right)^{r+{\scriptscriptstyle \frac{1}{2}}}},$$
transforms into:
$$G_𝐧=_0^{\mathrm{}}\mathrm{d}s\frac{s^{r\frac{3}{2}}}{(1+s)^{|𝐧\mathrm{}+𝐫}}e^{sy^2},|𝐧\mathrm{}=\underset{𝐚}{}𝐧_𝐚.$$
The noncommutative function is everywhere non-singular:
$$G_0(y)\sqrt{\pi }\left(2\frac{(r\frac{3}{2})!}{(r2)!}+|y|+\mathrm{}\right),y.0$$
Thus the map renders functions smoother at the origin.
Another example is when $`d=2r`$. Then we have:
$$𝒢(R)=\frac{1}{\left(_aR_a\right)^{r1}},$$
$$G_𝐧=_0^1d\lambda \lambda ^{r2}(1\lambda )^{|𝐧\mathrm{}}=\frac{(r2)!|𝐧|!}{(|𝐧|+r1)!}.$$
which is also non-singular everywhere (for $`r>1`$). In two dimensions ($`r=1`$) we get:
$$𝒢(R)=\mathrm{log}\left(\mu R\right)G_n=\mathrm{log}\mu +\psi (n+1)=\mathrm{log}\mu C+\underset{k=1}{\overset{N}{}}\frac{1}{k}.$$
The formula (3.1) is not applicable here since it gives a logarithmically divergent integral, of purely infrared origin. However, the divergence is $`n`$ independent, so that it affects $`G_n`$ by an additive constant, i.e. by a zero mode of the Laplacian. The presence of the divergence is reflected in the fact that the cutoff, $`\mathrm{log}\mu `$, in the commutative Green’s function appears in the noncommutative formula (3.1).
The commutative Green’s function $`𝒢(R)`$ solved Laplace’s equation with a delta function source:
$$\mathrm{\Delta }𝒢(R)=\delta ^d(x)$$
Let us see what the source is equal to now. For simplicity let us work in even number of dimensions, $`G_𝐧=G(|𝐧|)`$:
$$\frac{1}{4}\mathrm{\Delta }G(n)=(n+r)G(n+1)+nG(n1)(2n+r)G(n)=_0^1\mathrm{d}\left(\lambda ^r(1\lambda )^n\right)=\delta _{n,0}.$$
So the source got smoothed out:
$$\delta ^d(x)=\delta (R_a)\delta (y_b)\delta _{N_a,0}\delta (y_b).$$
In this way the noncommutativity of the space-time looks similar to the lattice regularization (although in the spherical rather then Cartesian way). However, by the above analysis, the formula (3.1) means that we have ended up with the Gaussian source $`D`$ as in (3.1).
On the other hand, the solution to the equation
$$\mathrm{\Delta }\widehat{G}=\widehat{\delta }$$
with the localized delta function source is also easy to produce: one simply applies the map (2.1) to the ordinary Green’s function $`𝒢_R`$. It is amusing that the result is close to the formula (3.1), namely the Green’s function is again diagonal in the eigenbasis of the occupation number operators $`n_a`$ and depends only on $`n=|𝐧|`$:
$$\widehat{G}_𝐧=\widehat{G}(n)=_0^2d\lambda \lambda ^{r2}(1\lambda )^n.$$
However, (3.1) and (3.1) differ considerably. In four dimensions, $`r=2`$, the difference is striking:
$$G(n)=\frac{1}{n+1},\widehat{G}(n)=\frac{1+(1)^n}{n+1}.$$
What is also striking is the failure of the classical limit for $`r>2`$: one might expect that, for large $`n`$, the Green’s function $`\widehat{G}(n)`$ would go over to its classical counterpart $`𝒢_R\frac{1}{n^{r1}}`$. The integral (3.1) indeed has this property - the integrand is peaked at $`\lambda =0`$ and the saddle point gives precisely the expected asymptotics. But the integral (3.1) has another saddle point at $`\lambda =2`$ which yields the leading asymptotics for $`r>2`$
$$\widehat{G}(n)\frac{(1)^n\mathrm{\hspace{0.17em}2}^{r2}}{n+1},n\mathrm{}.$$
The lesson to be drawn from here is that highly localized distributions (the delta function is such a distribution) become the operators spread out over all the Fock space that they act in, while the operators whose range is comparatively small (such as the Gaussian) correspond in fact to the distributions with finite support of order $`\sqrt{\theta }`$.
Finally, to construct the Green’s function $`\widehat{G}(xx^{})`$ we have to use the formula:
$$\widehat{G}(xx^{})=S^{}\left(\widehat{G}(x)\mathrm{Id}_x^{}\right)S,$$
which gives (for $`d=2r`$):
$$G(xx^{})=_0^2d\lambda \lambda ^{r2}(1\lambda )^{{\scriptscriptstyle \frac{1}{2}}(zz^{})(\overline{z}\overline{z}^{})}.$$
3.3. Gauge theory on noncommutative space
In an ordinary gauge theory with gauge group $`G`$ the gauge fields are connections in some principal $`G`$-bundle. The matter fields are the sections of the vector bundles with the structure group $`G`$. Noncommutative vector bundles are defined as projective modules over the algebra $`𝒜_\theta `$. This definition captures the following two properties of ordinary vector bundles: i) the sections of the bundle can be multiplied by functions on the base manifold - in this way the space of sections is acted on (linearly) by the space of functions — the definition of a module; ii) every vector bundle can be made trivial by the appropriate addition of another vector bundle - this is the definition of the projective module - it becomes free (equals to a direct sum of several copies of the algebra $`𝒜_\theta `$) when we add another module.
Now suppose we are given a module $`M`$ over the algebra $`𝒜_\theta `$. In the noncommutative case there are two types of modules, left and right. The elements $`m_𝐥`$ of the left module are multiplied by the elements $`a`$ of the algebra from the left, while the elements of the right module are multiplied from the right:
$$a:m_𝐥am_𝐥,m_𝐫m_𝐫a.$$
The left module over an algebra $`𝒜`$ is a right module over the algebra $`𝒜^{}`$ which is obtained from $`𝒜`$ by reversing the order of multiplication:
$$a^{}b=ba.$$
The notion of the left/right modules is analogous to the notion of chiral matter fields.
The connection $``$ is the operator
$$:𝐑^d\times MM,_\epsilon (m)M,\epsilon 𝐑^d,mM,$$
where $`𝐑^d`$ denotes the commutative vector space, the Lie algebra of the automorphism group generated by (2.1). The connection is required to obey the Leibnitz rule:
$$\begin{array}{cc}& _\epsilon (am_𝐥)=\epsilon ^i(_ia)m_𝐥+a_\epsilon m_𝐥\hfill \\ & _\epsilon (m_𝐫a)=m_𝐫\epsilon ^i(_ia)+(_\epsilon m_𝐫)a.\hfill \end{array}$$
As usual, one defines the curvature $`F_{ij}=[_i,_j]`$ \- the operator $`\mathrm{\Lambda }^2𝐑^d\times MM`$ which commutes with the multiplication by $`a𝒜_\theta `$. In other words, $`F_{ij}\mathrm{End}_𝒜(M)`$. If the right (left) module $`M`$ is free, i.e. it is a sum of several copies of the algebra $`𝒜`$ itself, then the connection $`_i`$ can be written as
$$_i=_i+A_i$$
where $`A_i`$ is the operator of the left (right) multiplication by the matrix with $`𝒜`$-valued entries. In the same operator sense the curvature obeys the standard identity:
$$F_{ij}=_iA_j_jA_i+A_iA_jA_jA_i.$$
Given the module $`M`$ one can multiply it by a free module $`𝒜^N`$ to make it a module over an algebra $`\mathrm{Mat}_{N\times N}(𝒜)`$ of matrices with elements from $`𝒜`$. In the non-abelian gauge theory over $`𝒜_\theta `$ we are interested in projective modules over $`\mathrm{Mat}_{N\times N}(𝒜_\theta )`$. If the algebra $`𝒜_\theta `$ (or perhaps its subalgebra) has a trace, $`\mathrm{Tr}`$, then the algebra $`\mathrm{Mat}_{N\times N}(𝒜_\theta )`$ has a trace given by the composition of a usual matrix trace and $`\mathrm{Tr}`$.
It is a peculiar property of the noncommutative algebras that the algebras $`𝒜`$ and $`\mathrm{Mat}_{N\times N}(𝒜)`$ have much in common. These algebras are called Morita equivalent and under some additional conditions the gauge theories over $`𝒜`$ and over $`\mathrm{Mat}_N(𝒜)`$ are also equivalent. This phenomenon is responsible for the similarity between the “abelian noncommutative” and “non-abelian commutative” theories.
4. Monopoles and Instantons
4.1. Lagrangian and couplings
After the preparations of the previous section the Lagrangian for the gauge theory is given by:
$$(A)=\frac{1}{4g_{\mathrm{YM}}^2}\underset{i,j}{}\mathrm{Tr}[_i,_j]^2.$$
If additional charged matter fields are present (elements $`\mathrm{\Phi }`$ of a module $`M`$) then the Lagrangian becomes:
$$(A,\mathrm{\Phi })=(A)+\underset{i}{}\mathrm{Tr}_i\mathrm{\Phi }_i\mathrm{\Phi }+\mathrm{}$$
The equations of motion following from the Lagrangian (4.1) are:
$$\underset{i}{}[_i,F_{ij}]=0$$
In four dimensions, $`d=4`$, the Euclidean version of (4.1) can be solved by solving the first order instanton equations:
$$F_{ij}=\pm \frac{1}{2}ϵ_{ijkl}F_{kl},$$
as follows from the Bianci identity, which holds irrespectively of the commutativity:
$$[_i,F_{kl}]+[_k,F_{li}]+[_l,F_{ik}]=0.$$
Introduce the complex coordinates: $`z_1=x_1+ix_2=x_+`$, $`z_2=x_3+ix_4`$. The instanton equations read:
$$\begin{array}{cc}& F_{z_1z_2}=0\hfill \\ & F_{z_1\overline{z}_1}+F_{z_2\overline{z}_2}=0.\hfill \end{array}$$
The first equation in (4.1) can be solved (locally) as follows:
$$A_{\overline{z}_a}=\xi ^1\overline{}_{\overline{z}_a}\xi ,A_{z_a}=_{z_a}\xi \xi ^1.$$
with $`\xi `$ a Hermitian matrix. Then the second equation in (4.1) becomes Yang’s equation:
$$\underset{a=1}{\overset{2}{}}\overline{}_{z_a}\left(_{z_a}\xi ^2\xi ^2\right)=0.$$
This ansatz works in the noncommutative setup as well .
If we look for the solutions to (4.1), that are invariant under translations in the $`4`$’th direction then we will find the monopoles of the gauge theory with an adjoint scalar Higgs field, where the role of the Higgs field is played by the component $`A_4`$ of the gauge field. The equations (4.1) in this case are called the Bogomolny equations, and they can be analyzed in the commutative case via Nahm’s ansatz .
The action (4.1) becomes the energy functional for the coupled gauge-adjoint Higgs system:
$$=\frac{1}{4g_{\mathrm{YM}}^2}d^3x\sqrt{\mathrm{det}G}\mathrm{Tr}\left(G^{ii^{}}G^{jj^{}}F_{ij}F_{i^{}j^{}}+G^{ij}_i\mathrm{\Phi }_j\mathrm{\Phi }\right)$$
where for the sake of generality we have introduced a constant metric $`G_{ij}`$.
$`\underset{¯}{\mathrm{Open}\mathrm{and}\mathrm{closed}\mathrm{string}\mathrm{moduli}.}`$
(4.1) emerges in the decoupling limit of a D3-brane in the Type IIB string theory in a background with a constant Neveu-Schwarz B-field. Let us recall the relation of the parameters of the actions (4.1), (4.1) and the string theory parameters, before taking the Seiberg-Witten limit .
We start with the D3-brane whose worldvolume is occupying the 0123 directions, and turn on a $`B`$-field:
$$\frac{1}{2}Bdx^1dx^2$$
The indices $`i,j`$ below will run from $`1`$ to $`3`$. We assume that the closed string metric $`g_{ij}`$ is flat, and the closed string coupling $`g_s`$ is small. According to the gauge theory on the D3-brane is described by a Lagrangian, which, when restricted to time-independent fields, coincides with (4.1), whose parameters
$$G_{ij},\theta ^{ij},g_{\mathrm{YM}}^2,$$
are related to
$$g_{ij},B_{ij},g_s$$
as follows:
$$\begin{array}{cc}& G_{ij}=g_{ij}(2\pi \alpha ^{})^2\left(Bg^1B\right)_{ij}\hfill \\ & \theta ^{ij}=(2\pi \alpha ^{})^2\left(\frac{1}{g+2\pi \alpha ^{}B}B\frac{1}{g2\pi \alpha ^{}B}\right)^{ij}\hfill \\ & g_{\mathrm{YM}}^2=2\pi g_s\left(\mathrm{det}\left(1+2\pi \alpha ^{}g^1B\right)\right)^{{\scriptscriptstyle \frac{1}{2}}}.\hfill \end{array}$$
Suppose the open string metric is Euclidean: $`G_{ij}=\delta _{ij}`$, then (4.1), (4.1) imply:
$$g=dx_3^2+\frac{(2\pi \alpha ^{})^2}{(2\pi \alpha ^{})^2+\theta ^2}\left(dx_1^2+dx_2^2\right),B=\frac{\theta }{(2\pi \alpha ^{})^2+\theta ^2},$$
and
$$g_s=g_{\mathrm{YM}}^2\frac{\alpha ^{}}{\sqrt{(2\pi \alpha ^{})^2+\theta ^2}}.$$
The Seiberg-Witten limit is achieved by taking $`\alpha ^{}0`$ with $`G,\theta ,g_{\mathrm{YM}}^2`$ kept fixed.
In this limit the effective action of the D3-brane theory will become that of the (super)Yang-Mills theory on a noncommutative space $`𝒜_\theta `$. The relevant part of the energy functional is:
$$=\frac{1}{2g_{\mathrm{YM}}^2}d^3x\mathrm{Tr}\left(B_iB_i+_i\mathrm{\Phi }_i\mathrm{\Phi }\right),$$
where
$$B_i=\frac{i}{2}\epsilon _{ijk}F_{jk}.$$
The fluctuations of the D3-brane in some distinguished transverse direction (which we called $`\mathrm{\Phi }`$) are described by the dynamics of the Higgs field.
As in the ordinary, commutative case, one can rewrite (4.1) as a sum of a total square and a total derivative:
$$=\frac{1}{2g_{\mathrm{YM}}^2}d^3x\mathrm{Tr}\left(_i\mathrm{\Phi }\pm B_i\right)^2_i\mathrm{Tr}\left(B_i\mathrm{\Phi }+\mathrm{\Phi }B_i\right)$$
The total derivative term depends only on the boundary conditions. So, to minimize the energy given boundary conditions we should solve the Bogomolny equations:
$$_i\mathrm{\Phi }=\pm B_i,i=1,2,3.$$
4.2. Nahm’s construction for commutative monopoles
To begin with, we review the techniques which have been used in the commutative case. Specifically, Bogomolny equations take the form:
$$_i\mathrm{\Phi }+B_i=0,i=1,2,3$$
The boundary condition is that at the spatial infinity the Higgs field approaches a constant, corresponding to the Higgs vacuum. In the case of $`SU(2)`$ this means that locally on the two-sphere at infinity:
$$\varphi (x)\mathrm{diag}(\frac{a}{2},\frac{a}{2}).$$
The solutions are classified by the magnetic charge $`k`$. By virtue of the equation (4.1) the monopole charge can be expressed as the winding number which counts how many times the two-sphere $`𝐒_{\mathrm{}}^2`$ at infinity is mapped to the coset space $`SU(2)/U(1)𝐒^2`$ of the abelian subgroups of the gauge group.
Nahm constructs solutions to the monopole equations as follows. Consider the matrix differential operator on the interval $`I`$ with the coordinate $`z`$:
$$i\mathrm{\Delta }=_z+𝒯_i\sigma _i,$$
where
$$𝒯_i=T_i(z)+x_i.$$
$`x_i`$ are the coordinates in the physical space $`𝐑^3`$, and the $`k\times k`$ matrices $`T_i(z)=T_i^{}(z)`$ obey Nahm’s equations:
$$_zT_i=i\epsilon _{ijk}T_jT_k,$$
with certain boundary conditions. We take $`I=(a/2,a/2)`$ where $`a`$ is given in (4.1). At $`zz_0`$, $`z_0=\pm a/2`$ we require that :
$$T_i\frac{t_i}{zz_0}+\mathrm{reg}.,[t_i,t_j]=i\epsilon _{ijk}t_k,$$
i.e. the residues $`t_i`$ must form a $`k`$-dimensional representation of $`SU(2)`$ (irreducible if the solution is to be non-singular).
Then one looks for the fundamental solution to the equation:
$$i\mathrm{\Delta }^{}\mathrm{\Psi }(z)=_z\mathrm{\Psi }𝒯_i\sigma _i\mathrm{\Psi }=0,$$
where
$$\mathrm{\Psi }=\left(\begin{array}{c}\mathrm{\Psi }_+\\ \mathrm{\Psi }_{}\end{array}\right),$$
and $`\mathrm{\Psi }_\pm `$ are $`k\times 2`$ matrices ($`k`$ is the monopole charge, and $`2`$ is for $`SU(2)`$), which must be finite at $`z=\pm a/2`$ and normalized so that:
$$𝑑z\mathrm{\Psi }^{}\mathrm{\Psi }=\mathrm{𝟏}_{2\times 2}.$$
(this $`2\times 2`$ is again for $`SU(2)`$.)
Then:
$$\begin{array}{cc}\hfill A_i& =𝑑z\mathrm{\Psi }^{}_i\mathrm{\Psi },\hfill \\ \hfill \mathrm{\Phi }& =𝑑zz\mathrm{\Psi }^{}\mathrm{\Psi }.\hfill \end{array}$$
Notice that the interval $`I`$ could be $`(a_1,a_2)`$ instead of $`(a/2,+a/2)`$. The only formula that is not invariant under shifts of $`z`$ is the expression (4.1) for $`\varphi `$. By shifting $`\varphi `$ by a scalar $`(a_1+a_2)/2`$ we can make it traceless and map $`I`$ back to the form we used above.
4.3. Nahm’s equations from the D-string point of view
The meaning of the Nahm’s equations becomes clearer in the D-brane realization of gauge theory and the D-string construction of monopoles. The endpoint of a fundamental string touching a D3-brane looks like an electric charge for the $`U(1)`$ gauge field on the brane. By S-duality, a D-string touching a D3-brane creates a magnetic monopole. If one starts with two parallel D3-branes, seperated by distance $`a`$ between them, one is studying the $`U(2)`$ gauge theory, Higgsed down to $`U(1)\times U(1)`$, where the vev of the Higgs field is
$$\mathrm{\Phi }=\left(\begin{array}{cc}a_1& 0\\ 0& a_2\end{array}\right)$$
One can suspend a D-string between these two D3-branes, or a collection of $`k`$ parallel D-strings. These would correspond to a charge $`k`$ magnetic monopole in the Higgsed $`U(2)`$ theory. The BPS configurations of these D-strings are described the corresponding self-duality equations in the 1+1 dimensional $`U(k)`$ gauge theory on the worldsheet of these strings . The equations (4.1) are exactly these BPS equations. The presence of the D3-branes is reflected in the boundary conditions (4.1). The matrices $`T_i`$ correspond to the “matrix” transverse coordinates $`X^i`$, $`i=1,2,3`$ to the D-strings, which lie within D3-branes.
4.4. Charge one monopoles
In the case $`k=1`$ the analysis simplifies: $`T_i=0`$, and
$$\mathrm{\Psi }=\left(\begin{array}{c}\left(_z+x_3\right)v\\ \left(x_1+ix_2\right)v\end{array}\right),_z^2v=r^2v,r^2=\underset{i}{}x_i^2.$$
The condition that $`\mathrm{\Psi }`$ is finite at both ends of the interval allows two solutions for (4.1):
$$v=e^{\pm rz},$$
which after imposing the normalization condition,(4.1), leads to:
$$\mathrm{\Psi }=\frac{1}{\sqrt{2\mathrm{sinh}(ra)}}\left(\begin{array}{cc}\sqrt{r+x_3}e^{rz}& \sqrt{rx_3}e^{rz}\\ \frac{x_+}{\sqrt{r+x_3}}e^{rz}& \frac{x_{}}{\sqrt{rx_3}}e^{rz}\end{array}\right),$$
where we used $`x_\pm =x_1\pm ix_2`$.
In particular,
$$\mathrm{\Phi }=\frac{1}{2}\left(\frac{a}{\mathrm{tanh}(ra)}\frac{1}{r}\right)\sigma _3.$$
4.5. Abelian ordinary monopoles
It is interesting that Nahm’s equations describe Dirac monopoles as well. To achieve this replace the interval $`(a/2,a/2)`$ by the interval $`(\mathrm{},a)`$. Intuitively this is natural, since in the $`U(1)`$ case the Higgs field $`\varphi `$ has only one eigenvalue at infinity.
Then the equation (4.1) becomes simply the condition that the abelian monopole has a magnetic potential $`\varphi `$, which must be harmonic. Let us find this harmonic function. The matrices $`T_i`$ can be taken to have the following form:
$$T_i(z)=\frac{t_i}{z},[t_i,t_j]=i\epsilon _{ijk}t_k,$$
where $`t_i`$ form an irreducible spin $`j`$ representation of $`SU(2)`$. Let $`V\text{ }\mathrm{C}^N,N=2j+1`$, be the space of this representation. The matrices $`\mathrm{\Psi }_\pm `$ are now $`V`$-valued. By an $`SU(2)`$ rotation we can bring the three-vector $`x_i`$ to the form $`(0,0,r)`$, i.e. $`x_1=x_2=0,x_3>0`$. Then one can show that in this basis
$$\mathrm{\Psi }_{}=0,\mathrm{\Psi }_+=\nu _jz^je^{rz}|j,$$
where $`|jV`$ is the highest spin state in $`V`$. The coefficient $`\nu _j`$ is found from the normalization condition:
$$|\nu _j|^2=\frac{r^{2j+1}}{(2j)!}.$$
From this we get the familiar formula for the singular Higgs field
$$\varphi =a\frac{N}{2r}.$$
5. Abelian noncommutative monopoles
In this section we study the solutions to the Bogomolny equations for U(1) gauge theory on a noncommutative three dimensional space. As before, we assume the Poisson structure ($`\theta `$) which deforms the multiplication of the functions to be constant. Then there is essentially a unique choice of coordinate functions $`x_1,x_2,x_3`$ such that the commutation relations between them are as follows:
$$\begin{array}{cc}\hfill [x_1,x_2]=& i\theta ,\theta >0\hfill \\ \hfill [x_1,x_3]=& [x_2,x_3]=0.\hfill \end{array}$$
This algebra, (5.1), defines noncommutative $`𝐑^3`$, which we denote by $`𝒜_\theta `$. Introduce the creation and annihilation operators $`c,c^{}`$:
$$c=\frac{1}{\sqrt{2\theta }}\left(x_1ix_2\right),c^{}=\frac{1}{\sqrt{2\theta }}\left(x_1+ix_2\right),$$
that obey
$$[c,c^{}]=1.$$
5.1. Noncommutative Nahm equations
We start by repeating the procedure that worked in the ADHM instanton case, namely we relax the condition that $`x_i`$’s commute but insist on the equation (4.1) with $`T_i`$ replaced by the relevant matrices $`𝒯_i=T_i+x_i`$. Then the equation (4.1) on $`T_i`$ is modified:
$$_zT_i=i\epsilon _{ijk}T_jT_k+\delta _{i3}\theta .$$
It is obvious that, given a solution $`T_i(z)`$ of the original Nahm equations, it is easy to produce a solution of the noncommutative ones:
$$T_i(z)^{\mathrm{nc}}=T_i(z)+\theta z\delta _{i3}.$$
From this it follows that, unlike the case of instanton moduli space, the monopole moduli space does not change under noncommutative deformation.
This deformation, (5.1) , is exactly what one gets by looking at the D-strings suspended between the D3-branes (or a semi-infinite D-string with one end on a D3-brane) in the presence of a $`B`$-field. One gets the deformation:
$$[X^i,X^j][X^i,X^j]i\theta ^{ij}=[T_i,T_j]\frac{1}{2}\theta \epsilon _{ij3}$$
The reason why $`\theta ^{ij}`$, instead of $`B_{ij}`$, appears on the right hand side of (5.1) is rather simple. By applying T-duality in the directions $`x_1,x_2,x_3`$ we could map the D-string into the D4-brane. The matrices $`X^1,X^2,X^3`$ become the components $`A_{\widehat{1}},A_{\widehat{2}},A_{\widehat{3}}`$ of the gauge field on the D4-brane worldvolume, and the $`B`$-field would couple to these gauge fields via the standard coupling $`F_{\widehat{i}\widehat{j}}\widehat{B}_{\widehat{i}\widehat{j}}`$, where $`\widehat{B}_{\widehat{i}\widehat{j}}`$ is the T-dualized $`B`$-field. It remains to observe that $`\widehat{B}_{\widehat{i}\widehat{j}}=\theta ^{ij}`$, since the T-dualized indices $`\widehat{i}`$ label the coordinates on the space, dual to that of $`x_i`$’s.
5.2. Solving Nahm’s equations
To solve (5.1) we imitate the approach for the $`k=1`$ commutative monopole by taking
$$T_{1,2}=0,T_3=\theta z.$$
To solve (4.1) for $`\mathrm{\Psi }`$ we introduce the operators $`b,b^{}`$:
$$b=\frac{1}{\sqrt{2\theta }}\left(_z+x_3+\theta z\right),b^{}=\frac{1}{\sqrt{2\theta }}\left(_z+x_3+\theta z\right),$$
which obey the oscillator commutation relations:
$$[b,b^{}]=[c,c^{}]=1.$$
We introduce the superpotential
$$W=x_3z+\frac{1}{2}\theta z^2,$$
so that $`b=\frac{1}{\sqrt{2\theta }}e^W_ze^W,b^{}=\frac{1}{\sqrt{2\theta }}e^W_ze^W`$. Then equation (4.1) becomes:
$$\begin{array}{cc}\hfill b^{}\mathrm{\Psi }_++& c\mathrm{\Psi }_{}=0\hfill \\ \hfill c^{}\mathrm{\Psi }_+& b\mathrm{\Psi }_{}=0.\hfill \end{array}$$
The general solution of (4.1) is:
$$\mathrm{\Psi }_+=\left(u_1b+u_2c\right)v,\mathrm{\Psi }_{}=\left(u_1c^{}u_2b^{}\right)v,$$
where $`v`$ satisfies
$$(b^{}b+cc^{})v=0,$$
and $`(u_1,u_2)`$ are two complex numbers defined up to multiplication by a common factor (which can be reabsorbed in the definition of $`v`$). Thus, the generic solution is parameterized by a point $`u=(u_1:u_2)\mathrm{𝐂𝐏}^1`$ on a two-dimensional (twistor) sphere.
To construct the solution we must now solve (5.1) and normalize $`\mathrm{\Psi }^{}\mathrm{\Psi }`$. Recall that $`z`$ is defined on the half line $`(\mathrm{},a)`$, so therefore $`b`$ and $`b^{}`$ are not Hermitian conjugates:
$$\psi _1^{}b\psi _2=\psi _1^{}\psi _2(z=0)+\left(b^{}\psi _1\right)^{}\psi _2.$$
As far as the $`c,c^{}`$ system is concerned we will work in the occupation number basis of the Fock space obtained by quantizing the $`(x_1,x_2)`$ plane. The most general expression for $`v`$ is:
$$v=\underset{n=0}{\overset{\mathrm{}}{}}v_{nm}(z,x_3)|nm|.$$
However it turns out that when $`u=(1:0)`$ or when $`u=(0:1)`$ one can make the following ansatz (which is equivalent to imposing axial symmetry on $`v`$):
$$v=\underset{n=0}{\overset{\mathrm{}}{}}v_n(z,x_3)|nn|,cc^{}|n=(n+1)|n.$$
In this case (5.1) becomes
$$b^{}bv_n(z,x_3)=(n+1)v_n(z,x_3),$$
in the class of functions that lead to a $`\mathrm{\Psi }`$ that is normalizable on the half line $`(\mathrm{},a)`$. It is obvious that if we can solve for $`v_0`$ then $`v_nb^nv_0`$; so the solution of (5.1) is
$$v_n=\nu _nb^n\phi (z),$$
where
$$\phi (z)e^{W(z)}_{\mathrm{}}^ze^{2W(p)}𝑑p$$
and $`\nu _n`$ are normalization constants to be determined below. Will this lead to a normalizable $`\mathrm{\Psi }`$? As $`z\mathrm{}`$ we have the estimate $`\phi (z)<e^{W(z)}`$, so that $`v_0`$, as well as all its descendants $`v_n`$, are good functions. Notice that we can only make it normalizable on a half line, which nicely fits with the intuition that the abelian Higgs field must have values that are bounded. $``$ Notice that by shifting the coordinate $`x_3x_3+\theta a`$ we can always make $`a=0`$ (this is impossible for $`\theta =0`$). From now on we assume $`a=0`$.
$``$ After this shift we see that the only dimensional parameter in the problem is $`\theta `$. Let us choose the length units in which $`2\theta =1`$.
In the case, where $`u`$ is a generic point on the two sphere, we have
$$v=\underset{n,m}{}\nu _{nm}b^n\phi (z)|nm|.$$
In this paper we shall only discuss the case where either $`u=(1:0)`$, or $`u=(0:1)`$.
5.3. The normalization condition
We start by considering the case $`u=(0:1)`$. Accordingly, $`\mathrm{\Psi }_+=cv,\mathrm{\Psi }_{}=b^{}v`$, and
$$\mathrm{\Psi }^{}\mathrm{\Psi }=\underset{n=0}{\overset{\mathrm{}}{}}\left[\left(b^{}v_n\right)^{}\left(b^{}v_n\right)+n|v_n|^2\right]|nn|=\underset{n=0}{\overset{\mathrm{}}{}}_z\left(v_n^{}b^{}v_n\right)|nn|.$$
The noncommutative version of the condition (4.1) is:
$$_{\mathrm{}}^0𝑑z\mathrm{\Psi }^{}\mathrm{\Psi }=1=\underset{n=0}{\overset{\mathrm{}}{}}|nn|,$$
thus
$$\left(v_n^{}b^{}v_n\right)(z=0)=1.$$
which reduces to the sequence of relations:
$$n|\nu _n|^2(^n\mathrm{{\rm Y}}(2x_3))(^{n1}\mathrm{{\rm Y}}(2x_3))=1$$
where
$$\begin{array}{cc}\hfill \mathrm{{\rm Y}}(z)=& _0^{\mathrm{}}e^{\frac{p^2}{2}+zp}𝑑p\hfill \\ & =\sqrt{\frac{\pi }{2}}e^{\frac{z^2}{2}}\left(1+\mathrm{erf}\left(\frac{z}{\sqrt{2}}\right)\right)\hfill \\ & =\underset{n=0}{\overset{\mathrm{}}{}}\frac{\left(\frac{n1}{2}\right)!}{n!}2^{\frac{n1}{2}}z^n\hfill \\ & \sqrt{2\pi }e^{\frac{z^2}{2}},z+\mathrm{}\hfill \\ & \frac{1}{|z|},z\mathrm{}.\hfill \end{array}$$
For $`n=0`$ (5.1) is explicitly given by the analytic continuation of (5.1) to $`n=0`$, namely $`|\nu _0|^2\mathrm{{\rm Y}}(2x_3)=1.`$ The function (5.1) obeys the following differential equation
$$_z\mathrm{{\rm Y}}(z)=z\mathrm{{\rm Y}}(z)+1,\mathrm{{\rm Y}}(0)=\sqrt{\frac{\pi }{2}}.$$
Introduce the expansion coefficients
$$\begin{array}{cc}\hfill \mathrm{{\rm Y}}(2x_3+y)=& \underset{n=0}{}\zeta _n\frac{y^n}{n!},\hfill \\ \hfill \zeta _n=& _0^{\mathrm{}}p^ne^{\frac{p^2}{2}+2px_3}𝑑p,\hfill \end{array}$$
which obey the following equations:
$$\begin{array}{cc}\hfill \zeta _{n+1}=& 2x_3\zeta _n+n\zeta _{n1}\hfill \\ \hfill _3\zeta _n=& 2\zeta _{n+1}\hfill \\ \hfill \zeta _n(x_3=0)=& 2^{\frac{n1}{2}}\left(\frac{n1}{2}\right)!.\hfill \end{array}$$
The recursion relation in (5.1) for $`n=0`$ is to be understood by analytic continuation as $`n0`$. In this limit we have $`n\zeta _{n1}1`$, as $`n0`$. Thus $`\zeta _1=2x_3\zeta _0+1`$, as can also be checked directly from (5.1).
To find the normalization constants we substitute (5.1) into (5.1) to deduce:
$$|\nu _n|^2=\frac{1}{n\zeta _n\zeta _{n1}}.$$
Again, for $`n=0`$ the last equality is understood as $`|\nu _0|^2=1/\zeta _0`$. This completes the solution. We have explicitly constructed $`v`$ and thus $`\mathrm{\Psi }_\pm `$, from which we can determine, using (4.1), the Higgs and gauge fields. To do this we shall need to evaluate the overlap integrals:
$$_{\mathrm{}}^0(b^n\phi )(b^{n+1}\phi )=\zeta _{n+2}\zeta _n\zeta _{n+1}^2=(n+1)\zeta _n^2n\zeta _{n+1}\zeta _{n1}.$$
Again, for $`n=0`$ the last equality is understood with $`n\zeta _{n1}=1`$ for $`n=0`$.
Let us also introduce the functions $`\xi ,\stackrel{~}{\xi }`$ and $`\eta =\stackrel{~}{\xi }^2`$ :
$$\stackrel{~}{\xi }(n)=\sqrt{\frac{\zeta _n}{\zeta _{n+1}}},\eta (n)=\frac{\zeta _n}{\zeta _{n+1}},\xi (n)=\sqrt{\frac{n\zeta _{n1}}{\zeta _n}}.$$
We will need the asymptotics of these functions for large $`x_3`$. Let $`r_n^2=x_3^2+n`$. For $`r_n+x_3\mathrm{}`$ we can estimate the integral in (5.1) by the saddle point method. The saddle point and the approximate values of $`\zeta _n`$ and $`\eta _n`$ are:
$$\begin{array}{cc}\hfill \overline{p}=& x_3+r_n\hfill \\ \hfill \zeta _n\sqrt{\frac{\pi }{r_n}}& \left(x_3+r_n\right)^{n+\frac{1}{2}}e^{\frac{1}{2}\left(x_3+r_n\right)\left(3x_3r_n\right)}\hfill \\ \hfill \eta _n\frac{1}{x_3+r_{n+1}}& \left(1+\frac{1}{4r_n^2}+\mathrm{}\right).\hfill \end{array}$$
5.4. The explicit solution for the gauge and Higgs fields
$`\underset{¯}{\mathrm{The}\mathrm{Higgs}\mathrm{Field}.}`$
The Higgs field is given by (4.1):
$$\mathrm{\Phi }=𝑑zz\mathrm{\Psi }^{}\mathrm{\Psi }\underset{n=0}{\overset{\mathrm{}}{}}\mathrm{\Phi }_n(x_3)|nn|,$$
it has axial symmetry, that is commutes with the the number operator $`c^{}c`$. Explicitly:
$$\begin{array}{cc}\hfill \mathrm{\Phi }_n=& v_n^{}b^{}v_n𝑑z=\hfill \\ & =\frac{\zeta _n}{\zeta _{n1}}\frac{\zeta _{n+1}}{\zeta _n}=_3\mathrm{log}\xi _n\hfill \\ & =(n1)\eta _{n2}n\eta _{n1},n>0\hfill \\ & =\frac{\zeta _1}{\zeta _0}=2x_3\frac{1}{\zeta _0},n=0.\hfill \end{array}$$
To arrive at the third line we used the fact that
$$\frac{1}{\eta _n}\frac{1}{\eta _{n+1}}=n\eta _{n1}(n+1)\eta _n,$$
which follows immediately from the recursion relation for the $`\zeta ^{}`$s in (5.1). These fields are finite at $`x_3=0`$. Indeed as $`x_30`$,
$$\mathrm{\Phi }_n(x_3=0)=\sqrt{2}(\frac{\left(\frac{n1}{2}\right)!}{\left(\frac{n2}{2}\right)!}\frac{\left(\frac{n}{2}\right)!}{\left(\frac{n1}{2}\right)!}).$$
At the origin:
$$\mathrm{\Phi }_0(x_3=0)=\sqrt{\frac{2}{\pi }}.$$
$`\underset{¯}{\mathrm{The}\mathrm{Gauge}\mathrm{Field}.}`$
Using (4.1) it is easy to see that the component $`A_3`$ vanishes
$$\begin{array}{cc}\hfill A_3=& (b^{}v_n)^{}_3(b^{}v_n)+nv_n^{}_3v_n=\hfill \\ & (b^{}v_n)^{}_3v_n(z=0)+(b^{}v_n)^{}v_n=\hfill \\ & _3\mathrm{log}\xi (n)+\frac{\zeta _{n+1}\zeta _{n1}\zeta _n^2}{\zeta _{n1}\zeta _n}=0.\hfill \end{array}$$
In the same gauge the components $`A_1,A_2`$ (which we consider to be anti-hermitian) are given by:
$$\begin{array}{cc}\hfill A_c=\frac{1}{2}\left(A_1+iA_2\right),& A_c^{}=\frac{1}{2}\left(A_1iA_2\right)=A_c^{}\hfill \\ \hfill A_c=& \mathrm{\Psi }^{}[\mathrm{\Psi },c^{}]\hfill \\ \hfill =& \xi ^1[\xi ,c^{}]=c^{}\left(1\frac{\xi (n)}{\xi (n+1)}\right).\hfill \end{array}$$
Again we see that the matrix elements of $`A_c`$ are all finite and non singular.
$`\underset{¯}{\mathrm{The}\mathrm{Field}\mathrm{strength}.}`$
From (5.1) we deduce:
$$F_{12}=2i\left(_cA_c^{}_c^{}A_c+[A_c,A_c^{}]\right)=$$
$$\begin{array}{cc}& 2\left([\frac{\xi (n)}{\xi (n+1)}c,c^{}\frac{\xi (n)}{\xi (n+1)}]1\right)=\hfill \\ & =2\underset{n>0}{}\left(1+(n+1)\left(\frac{\xi (n)}{\xi (n+1)}\right)^2n\left(\frac{\xi (n1)}{\xi (n)}\right)^2\right)|nn|+\hfill \\ & +2\left(1+\left(\frac{\xi (0)}{\xi (1)}\right)^2\right)|00|,\hfill \end{array}$$
from which it follows, that:
$$\begin{array}{cc}\hfill B_3(n)=& 2\left(1n\frac{\eta _{n1}}{\eta _n}+\left(n1\right)\frac{\eta _{n2}}{\eta _{n1}}\right)\hfill \\ \hfill B_c=& \frac{1}{2}\left(B_1+iB_2\right)=c^{}\frac{\xi (n)}{\xi (n+1)}\left(\mathrm{\Phi }(n)\mathrm{\Phi }(n+1)\right).\hfill \end{array}$$
with the understanding that at $`n=0`$:
$$B_3(0)=2\left(1\frac{\zeta _1}{\zeta _0^2}\right).$$
5.5. Checking the Bogomolny equations.
With our conventions it is relatively easy to check that our solution satisfies the Bogomolny equations everywhere:
$$\begin{array}{cc}& _3\mathrm{\Phi }=_3\mathrm{\Phi }=B_3\hfill \\ & _c\mathrm{\Phi }=\xi ^1_c\mathrm{\Phi }\xi =B_c.\hfill \end{array}$$
For example, consider the equation $`_3\mathrm{\Phi }=B_3`$. We have:
$$_3\mathrm{\Phi }(n)=_3[(n1)\eta _{n2}n\eta _{n1}].$$
Then we use the fact that
$$_3\eta _n=2\left(1\frac{\eta _n}{\eta _{n+1}}\right),$$
to see that
$$\begin{array}{cc}& _3\mathrm{\Phi }(n)=2(n1)\left(1\frac{\eta _{n2}}{\eta _{n1}}\right)2n\left(1\frac{\eta _{n1}}{\eta _n}\right)=\hfill \\ & 2\left(1n\frac{\eta _{n1}}{\eta _n}+(n1)\frac{\eta _{n2}}{\eta _{n1}}\right)=B_3(n).\hfill \end{array}$$
$`\underset{¯}{\mathrm{Other}\mathrm{solutions}\mathrm{and}\mathrm{Seiberg}\mathrm{Witten}\mathrm{map}.}`$
It is plausible that the solutions corresponding to the other values of $`u=(u_1:u_2)`$ also have a physical meaning. In fact, the solutions of the Dirac-Born-Infeld theory that describe a a D-string touching a D3-brane (or D-string suspended between two D3-branes) in the presence of the $`B`$-field suggest that $`\mathrm{\Phi }`$ is multi-valued. Moreover, the solution for $`\mathrm{\Phi }`$ is implicit, whereas our solution is explicit. On the other hand the Seiberg-Witten map from the noncommutative gauge fields to the commutative ones must map our explicit solution for $`(A,\mathrm{\Phi })`$ into the solution of the DBI theory. It could mean that our solution is just one branch of the full solution, somehow incorporating other choices of $`u`$. However, we have found that the solution corresponding to the choice $`u=(1:0)`$ does not quite satisfy the BPS equations everywhere. Instead, it has a source, localized along a semi-infinite string pointing in the $`x_3\mathrm{}`$ direction. Nevertheless, it is clear that the other “branches”, corresponding to the generic $`u`$, seem worth investigating further. It is also plausible that in order to have a better understanding of the matching of the solutions to the DBI theory and the noncommutative gauge theory one would need to incorporate the $`\alpha ^{}`$ corrections.
$`\underset{¯}{\mathrm{A}\mathrm{remark}\mathrm{concerning}\mathrm{instantons}.}`$
As in the ordinary gauge theory case the monopoles are the solutions of the instanton equations in four dimensions, that are invariant under translations in the fourth direction $`x_4`$. We observe that the solution presented above ((5.1),(5.1) ), can also be cast in the Yang form: Take $`\xi =\xi (x_3,n)`$ as in (5.1). Then $`_3\xi `$ commutes with $`\xi `$ and we can write $`_3\xi \xi ^1=_3\mathrm{log}\xi `$. The formulae (4.1) yield exactly (5.1) and (5.1) with $`\mathrm{\Phi }=iA_4`$. Indeed, the equation (4.1) is nothing but the first equation in (5.1).
5.6. Toda lattice
At this point it is worth mentioning the relation of the noncommutative Bogomolny equations with the Polyakov’s non-abelian Toda system (see for the recent studies of this system) on the semi-infinite one-dimensional lattice. Let us try to solve the equations (4.1) using the Yang ansatz and imposing the axial symmetry: we assume that $`\xi (x_1,x_2,x_3)=\xi (n,x_3)`$, $`n=c^{}c`$. Then the equation (4.1) for the $`x_4`$-independent fields reduces to the system:
$$_t(_tg_ng_n^1)g_ng_{n+1}^1+g_{n1}g_n^1=0$$
where
$$g_n(t)=\frac{e^{\frac{t^2}{2}}}{n!}\xi ^2(n,\frac{t}{2}),$$
(notice that $`g_n(t)`$ are ordinary matrices). In the $`U(1)`$ case we can write
$$g_n(t)=e^{\alpha _n(t)},$$
and rewrite (5.1) in a more familiar form:
$$_t^2\alpha _n+e^{\alpha _{n1}\alpha _n}e^{\alpha _n\alpha _{n+1}}=0$$
For $`n=0`$ these equations also formally hold if we set $`g_1=0`$ (this boundary condition follows both from the Bogomolny equations and the same condition is imposed on the Toda variables on the lattice with the end-points).
Our Higgs field $`\mathrm{\Phi }_n`$ has a simple relation to the $`\alpha `$’s:
$$\mathrm{\Phi }(x_3,n)=2x_3+\alpha _n^{}(2x_3).$$
Our solution to (5.1) is:
$$\alpha _n=\frac{1}{2}t^2+\mathrm{log}\left(\frac{n\zeta _{n1}(t/2)}{\zeta _n(t/2)}\right)\mathrm{log}(n!).$$
It is amusing that Polyakov’s motivation for studying the system (5.1) was the structure of loop equations for lattice gauge theory. Here we encountered these equations in the study of the continuous, but noncommutative, gauge theory, thus giving more evidence for their similarity.
We should note in passing that in the integrable non-abelian Toda system one usually has two ‘times’ $`t,\overline{t}`$, so that the equation (5.1) has actually the form :
$$_t(_{\overline{t}}g_ng_n^1)g_ng_{n+1}^1+g_{n1}g_n^1=0.$$
It is obvious that these equations describe four-dimensional axial symmetric instantons on the noncommutative space with the coordinates $`t,\overline{t},c,c^{}`$ of which only half is noncommuting.
5.7. The mass of the monopole
In this section we restore our original units, so that $`2\theta `$ has dimensions of (length)<sup>2</sup>. From the formulae (5.1) we can derive the following estimates:
$$\mathrm{\Phi }(n)\frac{1}{2r_n}=\frac{1}{2\sqrt{x_3^2+2\theta n}}n0,r\mathrm{}.$$
Instead, for $`n=0`$ we have:
$$\begin{array}{cc}& \mathrm{\Phi }(0)\frac{x_3}{\theta },x_3+\mathrm{}\hfill \\ & \mathrm{\Phi }(0)\frac{1}{2|x_3|},x_3\mathrm{}.\hfill \end{array}$$
The asymptotics of the magnetic field is clear from the Bogomolny equations and the behaviour of $`\mathrm{\Phi }`$. Thus, for example,
$$B_3(n)=_3\mathrm{\Phi }(n)=\frac{x_3}{2r_n^3},n0,$$
and similarly for the other components of $`B`$. This is easily translated into ordinary position space, as in the discussion following equation(3.1), since, for large $`n`$, $`B_i(n,x_3)B_i(x_1^2+x_2^2n,x_3)`$. Therefore the magnetic field for large values of $`x_3`$ and $`n`$, or equivalently large $`x_i`$ is that of a point-like magnetic charge at the origin. However the $`n=0`$ component of $`B_3`$ behaves differently for large positive $`x_3`$:
$$B_3(n=0)=_3\mathrm{\Phi }(0)=\frac{1}{\theta }.$$
Notice, that this is exactly the value of the $`B`$-field. Thus, in addition to the magnetic charge at the origin we have a flux tube, localized in a Gaussian packet in the $`(x_1,x_2)`$ plane, of the size $`\theta `$, along the positive $`x_3`$ axis. The monopole solution is indeed a smeared version of the Dirac monopole, wherein the Dirac string (the D-string!) is physical.
To calculate the energy of the monopole we use the Bogomolny equations to reduce the total energy to a boundary term:
$$\begin{array}{cc}\hfill =& \frac{1}{2g_{\mathrm{YM}}^2}d^3x\left(\stackrel{}{B}\stackrel{}{B}+\stackrel{}{}\mathrm{\Phi }\stackrel{}{}\mathrm{\Phi }\right)=\hfill \\ & \frac{1}{2g_{\mathrm{YM}}^2}d^3x\left(\stackrel{}{B}+\stackrel{}{}\mathrm{\Phi }\right)^2\frac{1}{2g_{\mathrm{YM}}^2}d^3x\stackrel{}{}\left(\stackrel{}{B}\mathrm{\Phi }+\mathrm{\Phi }\stackrel{}{B}\right)=\hfill \\ & \frac{2\pi \theta }{2g_{\mathrm{YM}}^2}𝑑x_3\underset{n}{}n|_3^2\mathrm{\Phi }^2+4_c\left(\xi ^2\left(_c^{}\mathrm{\Phi }^2\right)\xi ^2\right)|n,\hfill \end{array}$$
where in the last line we switched back to the Fock space, by using the relation:
$$𝑑x_1𝑑x_2f(x_1,x_2)=2\pi \theta \mathrm{Tr}_{}\widehat{f}$$
Thus, the energy is given by the boundary term. To evaluate this expression we need to figure out what the boundary term is in the noncommutative, Fock space setup?
Consider the derivative terms in (5.1) involving $`x_1,x_2`$. They can be expressed as the commutators with $`c`$ or $`c^{}`$. In computing the trace
$$\mathrm{Tr}_{}[c,𝒳]=\underset{n}{}n|[c,𝒳]|n,$$
where we denote by $`𝒳`$ the terms $`\mathrm{\Phi }B_c+B_c\mathrm{\Phi }`$ in (5.1) , we get naively get that the trace of a commutator is zero. But we should be careful, since the matrices are infinite and the trace is an infinite sum. If we regulate it by restricting the sum to $`nN`$, then the matrix element $`N|c|N+1N+1|𝒳|N`$ is not cancelled, so that the regularized trace is
$$\mathrm{Tr}__N[c,𝒳]=\sqrt{N+1}N+1|𝒳|N$$
and similarly for $`c^{}`$. Let us choose as the infrared regulator box the “region” where $`|x_3|L,0nN`$, $`L\sqrt{2\theta N}1`$. Then the total integral in (5.1) reduces to the sum of two terms (up to the factor $`\frac{\pi \theta }{g_{\mathrm{YM}}^2}`$):
$$\begin{array}{cc}& 4N_L^L𝑑x_3\frac{\eta _{N1}}{\eta _N}\left(\mathrm{\Phi }_N^2\mathrm{\Phi }_{N+1}^2\right)\hfill \\ & +\underset{n=0}{\overset{N}{}}_3\mathrm{\Phi }_n^2|_{x_3=L}^{x_3=+L}.\hfill \end{array}$$
The first line in (5.1) is easy to evaluate. Since $`x_3+r_nL+\sqrt{L^2+2\theta N}\sqrt{\theta N}\mathrm{}`$ we can use the asymptotic expressions (5.1) to make an estimate:
$$4N_L^L𝑑x_3\frac{\zeta _{N+1}\zeta _{N1}}{\zeta _N^2}\left(\mathrm{\Phi }_N^2\mathrm{\Phi }_{N+1}^2\right)_L^L𝑑x\frac{2\theta N}{(x^2+2\theta N)^2}$$
$$\frac{L}{L^2+2\theta N}0$$
The second line in (5.1) contains derivatives of the Higgs field evaluated at $`x_3=L0`$ and at $`x_3=L0`$. The former is estimated using the $`z0`$ asymptotics in (5.1) or (5.1), and produces:
$$\underset{n=0}{\overset{N}{}}_3\mathrm{\Phi }_n^2(x_3=L)\frac{2\theta (N1)}{L^3}+2\frac{L}{\theta ^2}$$
The diverging with $`L`$ piece comes solely from the $`n=0`$ term. Finally, the $`x_3=L`$ case is treated via $`z0`$ asymptotics in (5.1) yielding the estimate $`\theta N/L^3`$ vanishing in the limit of large $`L,N`$.
Hence the total energy is given by
$$\frac{2\pi \theta \times 2L}{2g_{\mathrm{YM}}^2\theta ^2}=\frac{2\pi L}{g_{\mathrm{YM}}^2\theta },$$
which is the mass of a string of length $`L`$ whose tension is
$$T=\frac{2\pi }{g_{\mathrm{YM}}^2\theta }.$$
5.8. Magnetic charge
It is instructive to see what is the magnetic charge of our solution. On the one hand, it is clearly zero:
$$Q_{(\mathrm{space})}\stackrel{}{B}𝑑\stackrel{}{S}=d^3x\stackrel{}{}\stackrel{}{B}=0$$
since the gauge field is everywhere non-singular. On the other hand, we were performing a $`\theta `$-deformation of the Dirac monopole, which definitely had magnetic charge. To see what has happened let us look at (5.1) more carefully. We again introduce the box and evaluate the boundary integral (5.1) as in (5.1):
$$\frac{Q}{2\pi }=\underset{n=0}{\overset{N}{}}\left[B_3(x_3=L,n)B_3(x_3=L,n)\right]+4N_L^L𝑑x_3\frac{\eta _{N1}}{\eta _N}\left(\mathrm{\Phi }_N\mathrm{\Phi }_{N+1}\right)$$
It is easy to compute the sums
$$\begin{array}{cc}& \underset{n=0}{\overset{N}{}}B_3(x,n)=_3\frac{\zeta _{N+1}}{\zeta _N}\hfill \\ & 4N_L^L𝑑x_3\frac{\eta _{N1}}{\eta _N}\left(\mathrm{\Phi }_N\mathrm{\Phi }_{N+1}\right)=4(N+1)\frac{\xi _N^2}{\xi _{N+1}^2}𝑑\mathrm{log}\frac{\xi _N}{\xi _{N+1}}=\hfill \\ & =2(N+1)\left(\frac{\xi _N}{\xi _{N+1}}\right)^2|_{x_3=L}^{x_3=+L}=2N\frac{\zeta _{N1}\zeta _{N+1}}{\zeta _N^2}|_{x_3=L}^{x_3=+L},\hfill \end{array}$$
and the total charge vanishes as:
$$Q=\left[2N\frac{\zeta _{N1}\zeta _{N+1}}{\zeta _N^2}+_3\left(\frac{\zeta _{N+1}}{\zeta _N}\right)\right]_{x_3=L}^{x_3=+L}2(N+1)|_{x_3=+L}2(N+1)|_{x_3=L}.$$
We can better understand the distribution of the magnetic field by looking separately at the fluxes through the “lids” $`x_3=\pm L`$ of our box and through the “walls” $`n=N`$.
The walls contribute
$$\left[2N\frac{\zeta _{N1}\zeta _{N+1}}{\zeta _N^2}\right]_{x_3=L}^{x_3=+L}\frac{L}{\sqrt{L^2+N}}1,$$
while the lids contribute $`+1`$. Let us isolate the term $`B_3(+L,n=0)+2`$ (recall (5.1)). It contributes to the flux through the upper lid. The rest of the flux through the lids is therefore $`1`$. Hence the flux through the rest of the “sphere at infinity” is $`2`$ and it is roughly uniformly distributed ($`1`$ contribute the walls and $`1`$ the lids). So we get a picture of a spherical magnetic field of a monopole together with a flux tube pointing in one direction.
This spherical flux becomes observable in the naive $`\theta 0`$ limit, in which the string becomes localized at the point $`x_3=0`$, $`n=0`$ (since the slope of the linearly growing $`\mathrm{\Phi }_0\frac{x_3}{\theta }`$ becomes infinite). In the $`\theta =0`$ limit we throw out this point and all of the string.
6. Discussion
In this paper we found an explicit analytic expression for a soliton in the U(1) gauge theory on a noncommutative space. The solution describes a magnetic monopole attached to a finite tension string, that runs off to infinity tranverse to the noncommutative plane. This soliton has a clear reflection in type IIB string theory. If the gauge theory is realized as the $`\alpha ^{}0`$ limit of the theory on a D3-brane in the IIB string theory in the presence of a background NS B-field, then the monopole with the string attached is nothing but the D1-string ending on the D3-brane. What is unusual about the solution that we found is that it describes this string as a non-singular field configuration.
Whether this string is a dynamical object in the gauge theory, with full stringy degrees of freedom, remains to be determined. To this end we should analyze the spectrum of the fluctuations of this string. Several remarks are in order:
$``$ First, the “location” of the string is not very well defined., In noncommutative gauge theory the local energy density, as all local operators, is not gauge invariant. However the energy of a line element of a string, as a function of $`x_3`$ is a well-defined gauge-invariant notion:
$$t(x_3)=\frac{1}{2g_{\mathrm{YM}}^2}𝑑x_1𝑑x_2\left(\stackrel{}{B}^2+\left(\stackrel{}{}\mathrm{\Phi }\right)^2\right).$$
For our solution this “tension” turned out to be exponentially small for $`x_3<0`$ and essentially a constant
$$t=\frac{2\pi }{g_{\mathrm{YM}}^2\theta },$$
for $`x_3>0`$.
$``$ We worked in the gauge where $`A_3=0`$, which still allows for $`x_3`$-independent gauge transformations. This gauge freedom is broken down to a global $`U(1)`$ rotation by demanding that for $`x_3+\mathrm{}`$
$$\mathrm{\Phi }_0(x_3)\frac{x_3}{\theta }.$$
If we impose this asymptotic behaviour on $`\mathrm{\Phi }`$, then the problem of finding the spectrum of the fluctuations of the string becomes well-posed.
$``$ Our solution breaks translational invariance. One would expect the derivatives $`_\mu (\mathrm{\Phi },A_c,A_c^{})`$ to show up as zero modes. However, the derivatives in the $`x_1,x_2`$ directions are infinitesimal gauge transformations, while the derivatives in the $`x_3`$ direction are not normalizable:
$$_c\mathrm{\Phi }=[\mathrm{\Phi },c^{}],_cA_\mu =D_\mu c^{}\delta _{\mu ,c^{}}$$
(the shift of $`A_\mu `$ by a constant is a symmetry of the theory).
The next subject which we plan to elaborate further on is the extension of our analytic solution to the case of $`U(2)`$ noncommutative gauge theory. In this case one expects to find strings of finite extent, according to the brane picture .
What is the relation between the string we have found and the electric flux strings found recently in . These authors also study the coupled gauge field - Higgs field system, with the Higgs field in the adjoint representation. Their Higgs field $`t`$ arises from the open string tachyon, and has a non-trivial potential $`V(t)`$. In the limit of large noncommutativity $`\theta `$ the kinetic term can be neglected, according to , and the soliton can be found as a Gaussian wave-packet localized at the origin of the transverse plane to the would-be-string space, with the values of the tachyon field at the origin and far away given by at two different critical points of the potential $`V(t)`$. In our case we have no potential for $`\mathrm{\Phi }`$, nor did we assume $`\theta `$ to be large. However, our solitonic string also had an effective thickness of the order of $`\theta `$, and also disappears when $`\theta =0`$. It would be interesting to see, whether S-duality will map our magnetic strings to the electric strings of .
As a step in this direction we would like to compare the tension of our string with that of D-string (the authors of claim to have a complete agreement of the tension of their soliton with the tension of the fundamental string). As already mentioned, a D-string ending on a D3-brane in the presence of the constant $`B`$-field bends. To analyze this bending one could use the exact solution of the DBI theory , the B-deformed spike solutions of . However, for our qualitative analysis, it is sufficient to look at the linearized equations. If we replace the DBI Lagrangian by its Maxwell approximation, then the BPS equations in the presence of the $`B`$-field will have the form:
$$B_{ij}+F_{ij}+\sqrt{\mathrm{det}g}\epsilon _{ijk}g^{kl}_l\mathrm{\Phi }=0,$$
where we should use the closed string metric (4.1). The solution of (6.1) is:
$$\mathrm{\Phi }=B\left(1+\left(\frac{\theta }{2\pi \alpha ^{}}\right)^2\right)x_3\frac{1}{2r},r^2=x_3^2+\frac{1}{\left(1+\left(\frac{\theta }{2\pi \alpha ^{}}\right)^2\right)}\left(x_1^2+x_2^2\right).$$
The linearly growing piece in $`\mathrm{\Phi }`$ should be interpreted as a global rotation of the D3-brane, by an angle $`\psi `$, $`\mathrm{tan}\psi =\frac{\theta }{(2\pi \alpha ^{})}`$. This conclusion remains correct even after the full non-linear BPS equation is solved (see . Notice however that we fix $`G_{ij}=\delta _{ij}`$ instead of $`g_{ij}=\delta _{ij}`$ as in ). The singular part of $`\mathrm{\Phi }`$, the spike, represents the D-string. If we rotate the brane, then the spike forms an angle $`\frac{\pi }{2}\psi `$ with the brane. If we project this spike on the brane, then the energy, carried by its shadow per unit length, is related to the tension of the D-string via:
$$\frac{T_{D1}}{\mathrm{sin}\psi }=\frac{1}{2\pi \alpha ^{}g_s}\frac{\sqrt{(2\pi \alpha ^{})^2+\theta ^2}}{\theta }=\frac{(2\pi \alpha ^{})^2+\theta ^2}{2\pi g_{\mathrm{YM}}^2(\alpha ^{})^2\theta }.$$
However, this is not the full story. The endpoint of the D-string is a magnetic charge, which experiences a constant force, induced by the background magnetic field. If we had introduced a box of the extent $`2L`$ in the $`x_3`$-direction, then in order to bring a tilted D-string into our system from outside of the box we would have had to spend an energy equal to $`\frac{T_{D1}}{\mathrm{sin}\psi }L`$, but we would have been helped by the magnetic force, which would decrease the work done by
$$\frac{2\pi }{g_{\mathrm{YM}}^2}B_3=\frac{2\pi }{g_{\mathrm{YM}}^2}B_{12}g^{11}g^{22}\sqrt{g}\frac{1}{(\alpha ^{}g_{\mathrm{YM}})^2}\theta .$$
In sum, the energy of the semi-infinite D-string in the box will be given by<sup>1</sup> We thank K. Hashimoto for bringing a very helpful argument from the second reference in to our attention:
$$\frac{2\pi }{g_{\mathrm{YM}}^2\theta }.$$
This expression coincides with our tension (6.1). On dimensional grounds, non-commutative gauge theory cannot produce any other dependence of the tension on $`\theta `$ but that given in (6.1).
Finally, the large $`\theta `$ limit of the noncommutative gauge theory may provide an exciting opportunity to learn more about the large $`N`$ non-abelian commutative gauge theories, for both theories become essentially planar in this limit. If we keep $`g_{\mathrm{YM}}^2`$ small and take $`\theta \mathrm{}`$ then our magnetic strings become tensionless. Whether this could lead to condensation of the magnetic charges and a mechanism for confinement remains to be seen.
References
relax H. S. Snyder, “Quantized Space-Time”, Phys. Rev. 71 (1947) 38; “The Electromagnetic Field in Quantized Space-Time”, Phys. Rev. 72 (1947) 68 relax A. Connes, “Noncommutative geometry”, Academic Press (1994) relax A. Connes, M. Douglas, A. Schwarz, JHEP 9802(1998) 003 relax V. Schomerus, “D-Branes and Deformation Quantization”, JHEP 9906(1999) 030 relax N. Seiberg, E. Witten, hep-th/9908142, JHEP 9909(1999) 032 relax T. Filk, “Divergencies in a Field Theory on Quantum Space”, Phys. Lett. 376B (1996) 53 relax S. Minwala, M. van Raamsdonk, N. Seiberg, “Noncommutative Perturbative Dynamics”, hep-th/9912072 relax E. Witten, Nucl. Phys. B268 (1986) 253 relax M. van Raamsdonk, N. Seiberg, “Comments of Noncommutative Perturbative Dynamics”, hep-th/0002186, JHEP 0003(2000) 035 relax N. Nekrasov, A. S. Schwarz, hep-th/9802068, Comm. Math. Phys. 198 (1998) 689 relax H. Braden, N. Nekrasov, hep-th/9912019; K. Furuuchi, hep-th/9912047 relax K. Hashimoto, H. Hata, S. Moriyama, hep-th/9910196, JHEP 9912(1999) 021; A. Hashimoto, K. Hashimoto, hep-th/9909202, JHEP 9911(1999) 005; K. Hashimoto, T. Hirayama, hep-th/0002090 relax L. Jiang, “Dirac Monopole in Non-Commutative Space”, hep-th/0001073 relax R. Gopakumar, S. Minwala, J. Maldacena, A. Strominger, hep-th/0005048; O. Ganor, G. Rajesh, S. Sethi, hep-th/00050046 relax J. Harvey, P. Kraus, F. Larsen, E. Martinec, hep-th/0005031 relax S. Gukov, I. Klebanov, A. Polyakov, hep-th/9711112, Phys. Lett. 423B (1998) 64-70 relax N. Seiberg, L. Susskind, N. Toumbas, hep-th/0005040 relax M. Kontsevich, “Deformation quantization of Poisson manifolds”, q-alg/9709040 relax A. Dhar, G. Mandal and S. R. Wadia, Mod. Phys. Lett. A7 (1992) 3129-3146; A. Dhar, G. Mandal and S. R. Wadia, Int. J. Mod. Phys. A8 (1993) 3811-3828; A. Dhar, G. Mandal and S. R. Wadia, Mod. Phys. Lett. A8 (1993) 3557-3568; A. Dhar, G. Mandal and S. R. Wadia, Phys. Lett. 329B (1994) 15-26 relax I. Bars, D. Minic, “Non-Commutative Geometry on a Discrete Periodic Lattice and Gauge Theory”, hep-th/9910091 relax W. Nahm, Phys. Lett. 90B (1980) 413; W. Nahm, “The Construction of All Self-Dual Multimonopoles by the ADHM Method”, in “Monopoles in quantum field theory”, Craigie et al., Eds., World Scientific, Singapore (1982) ; N.J. Hitchin, Comm. Math. Phys. 89 (1983) 145 relax D.-E. Diaconescu, Nucl. Phys. B503 (1997) 220-238, hep-th/9608163 relax D. Bak, Phys. Lett. 471B (1999) 149-154, hep-th/9910135 relax S. Moriyama, hep-th/0003231 relax D. Mateos, “Noncommutative vs. commutative descriptions of D-brane BIons”, hep-th/0002020 relax P. Etingof, I. Gelfand, V. Retakh, “Factorization of differential operators, quasideterminants, and nonabelian Toda field equations” q-alg/9701008 relax R. Gopakumar, S. Minwala, A. Strominger, hep-th/0003160, JHEP 0005(2000) 020 relax C. G. Callan, Jr., J. M. Maldacena, Nucl. Phys. B513 (1998) 198-212, hep-th/9708147
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# Quantum correlations from local amplitudes and the resolution of the Einstein-Podolsky-Rosen nonlocality puzzle
## 1 Introduction
Quantum nonlocality as manifested in the EPR correlations has been, without doubt, the most important unresolved problem in the foundational aspects of physics. All experiments on quantum correlations and test of Bell’s inequality, and their present interpretations based on the multiparticle wavefunction in quantum mechanics suggest that there is nonlocality. Yet, we understand neither the nature of this nonlocality nor the physical mechanism that establishes the nonlocal correlations. The situation is somewhat akin to that in the earlier part of the last century when the ether was thought to be a necessary concept for the propagation of electromagnetic waves, yet something not apparent or detectable.
Sixty five years ago, Einstein, Podolsky and Rosen (EPR) addressed the question whether the wave-function represented a complete description of reality in quantum mechanics, and argued that it didn’t. The crucial and essential assumption in their argument was strict nonlocality, in the spirit of special relativity. Also, they had considered and included the concept of objective reality in the analysis. If the value of an observable was predictable with certainty without a measurement, the observable had a physical reality according to EPR. Their assertions lead to attempts at constructing a hidden variable theory that was more complete. Bell’s analysis of the EPR problem in the early sixties established the Bell’s inequalities obeyed by any local hidden variable theory for the correlations of entangled particles . Quantum mechanical correlations calculated using the entangled wave-function and spin operators violate these inequalities. Various experiments have established beyond doubt that there cannot be a viable local realistic hidden variable description of quantum mechanics . Further, these results also have been interpreted as evidence for nonlocal influences in quantum measurements involving entangled particles. Since no instruction set carried by the particles from their source of origin (possibly with the addition of several local hidden variables) can manage to create the correct correlations observed in experiments, the only way out seems to be that measurement of an observable on one of the particles in an entangled pair seems to convey the result of this measurement instantaneously to the other particle resulting in the correct behaviour of the other particle during a measurement on the second particle. In the quantum mechanical terminology, the measurement of an observable on one of the particles collapses the entire wave-function instantaneously and nonlocally and the second particle acquires a definite value for the same observable, consistent with the relevant conservation law. The no signalling theorems in this context prohibit any faster than light signalling using this feature, and therefore signal locality is not violated. But the stronger requirement of Einstein locality is violated. We seem to be stuck with the puzzling nonlocality which is probably the deepest mystery in the behaviour of entangled systems.
Accepting the concept of nonlocality without being able to understand its nature is already a disturbing feature. There is also serious conflict with the spirit of relativity. If one measurement precede the other in one frame, one can always find a moving frame in which the converse it true, the second measurement preceding the first . Therefore, one cannot attribute cause and effect relationships for measurements causing nonlocal collapse.
It turns out that the long-standing problem of EPR nonlocality is resolved by a simple quantum step that is physically well motivated . The crucial new idea is to incorporate the fact that quantum systems have ‘wave aspects’ in their behaviour and all calculations should take into account the phase relationships (coherence) that might be there in the multiparticle system. If locality is assumed at the level of probability amplitudes, and if the correlation is calculated directly from these amplitudes the correct quantum correlations emerge. This means that the EPR definition of objective reality was too restrictive. There is objective reality, but that is at the level of quantum phases and not at the level of eigenvalues. The quantum correlation is encoded in the relative phase appropriate for the problem. In the local hidden variable theories the correlations are calculated from eigenvalues and this procedure does not preserve the phase information. The situation has some analogy to the description of interference in quantum mechanics. Any attempt to reproduce the interference pattern using locality and the information on ‘which-path’ will fail since the phase information is lost or modified in such an attempt.
## 2 The solution of the EPR nonlocality puzzle
Consider the breaking up of a correlated state as in the standard Bohm version of the EPR problem . The two-particle state is described by the wave function
$$\mathrm{\Psi }_S=\frac{1}{\sqrt{2}}\{|1,1|1,1\}$$
(1)
where the state $`|1,1`$ is short form for $`|1_1|1_2`$, and represents an eigenvalue of $`+1`$ for the first particle and $`1`$ for the second particle if measured in any particular direction. $`\mathrm{\Psi }_S`$ is inherently nonlocal, describing both particles together, even when they are far apart in space-like separated regions.
Two observers make measurements on these particles individually at space like separated regions with time stamps such that these results can be correlated later through a classical channel. We assume that strict locality is valid at the level of probability amplitudes . A measurement changes probability amplitudes only locally. Measurements performed in one region do not change the magnitude or phase of the complex amplitude for the companion particle in a space-like separated region. The local setting of the polarizers, analyzers, Stern-Gerlach analyzers etc. (collectively denoted as analyzer) is represented by $`𝐚`$ and $`𝐛`$ for the two distant apparatus. These could be the directions of the analyzers, for example. Since we need to deal with correlated particles which may have a definite phase relationship at source (when the particles are produced together, for example) we introduce internal variables associated with each particle. We denote these variable as $`\varphi _1`$ and $`\varphi _2`$. Their values are unaltered once the particles are separated. Measurement on one particle does not change the value of this internal variable for the other particle. The assumption of locality is that the amplitudes (as opposed to eigenvalues) are functions of only these local variables.
We now state the assumptions (1 and 2) and some related comments (3, 4,and 5) :
1. The local amplitude for the first particle $`C_1`$ that decide the passage of the particle through an analyzer depends only on the local variables $`𝐚`$ and $`\varphi _1.`$ Similarly $`C_2`$ depends only on $`b`$ and $`\varphi _2`$. If we denote the passage as $`+`$ and the alternate outcome as $`,`$ then the statement of locality is for the relevant amplitudes is
$$C_{1\pm }=C_{1\pm }(𝐚,\varphi _1),C_{2\pm }=C_{2\pm }(𝐛,\varphi _2)$$
(2)
2. The correlations of the particles are encoded in the difference of the internal variables $`\varphi _1`$ and $`\varphi _2.`$ If the particles have perfect correlations at source then all the pairs in the ensemble have the same value for the difference $`\left|\varphi _1\varphi _2\right|=\varphi _0.`$
3. We do not make any assumption on determinism. Given the initial values of the internal variables $`\varphi _1`$ and $`\varphi _2`$, we do not attempt to make any prediction of the eigenvalues that would be measured in each run of the experiment.
4. We do not assume any hidden variable in the problem. The variables $`\varphi _1`$ and $`\varphi _2`$ are associated with the particles and they could be considered as hidden variables in a formal sense, though these values are not measurable since only relative phases are measurable.
5. We will also state the locality at the level of the eigenvalues, though we do not use this in the calculation. For observables $`A`$ and $`B`$,
$$A(𝐚,\varphi _1)=\pm 1,B(𝐛,\varphi _2)=\pm 1$$
(3)
This is the same locality assumption as in local realistic theories . But, this has a meaning different from its meaning in standard local realistic theories. Here, this means that the outcomes, when measured, depend only on the local setting and the local internal variable. There is no objective reality to $`A`$ and $`B`$ before a measurement. There is objective reality to $`\varphi _1`$ and $`\varphi _2,`$ but there is no way to observe these absolute phases.
Note that $`\varphi `$ is not a dynamical phase evolving as the particle propagates. It is an internal variable whose difference (possibly zero) remains constant for the particles of the correlated pair. The value of $`\varphi `$ can vary from particle to particle, but the relative phase $`\varphi _0`$ between the two particles in all correlated pairs is constant. Consider $`\varphi `$ as a reference for the particles to determine the angle of a polarizer or analyzer encountered on their way, locally.
The first particle encounters analyzer1 kept at an angle $`\theta _1`$ with respect to some global direction. We denote this angle of the analyzer with reference to $`\varphi `$ as $`\theta .`$ Similarly, the second particle which has the internal phase angle $`\varphi +\varphi _0,`$ where $`\varphi _0`$ is a constant, encounters the second analyzer oriented at angle $`\theta _2`$ at another space-like separated point. Let the orientation of this analyzer with respect to the internal phase angle of the second particle is $`\theta ^{}.`$ We have $`\theta \theta ^{}=\theta _1\theta _2+\varphi _0.`$
An experiment in which each particle is analyzed by orienting the analyzers at various angles $`\theta _1`$ and $`\theta _2`$ is considered next. At each location the result is two-valued denoted by ($`+1`$) for transmission and ($`1`$) for absorption of each particle, for any angle of orientation. The classical correlation function, which is also the experimenter’s correlation function, $`P(𝐚,𝐛)=\frac{1}{N}(A_iB_i)`$ satisfies $`1P(𝐚,𝐛)1.`$ Here $`(𝐚,𝐛)`$ denotes the two directions along which the analyzers are oriented and $`A_i`$ and $`B_i`$ are the two valued results. We note that $`P(𝐚,𝐛)`$ denotes the average of the quantity (number of detections in coincidence $``$ number of detections in anticoincidence), where ‘coincidence’ denotes both particles showing same value for the measurement and ‘anticoincidence’ denotes those with opposite values. The defect in the local realistic theories is that they try to calculate this correlation essentially by averaging over the products of eigenvalues. Obviously the phase information is thrown away in this procedure and there is no way, conceptually, such an attempt would have reproduced quantum mechanical results. We calculate the experimenter’s correlations starting from local amplitudes.
Now we state the expressions for the amplitudes and the amplitude correlation function. This is new physical input .
1. The local amplitudes for transmission associated with the first particle is $`C_{1+}=\frac{1}{\sqrt{2}}\mathrm{exp}(i\theta s)`$ for measurements at analyzer1, and for the second particle it is $`C_{2+}=\frac{1}{\sqrt{2}}\mathrm{exp}(i\theta ^{}s)`$ at analyzer2. The amplitudes for the orthogonal outcome are $`C_1`$ and $`C_2`$ and these are rotated by $`\pi /2`$ in the complex plane from the amplitudes for transmission. The square of the amplitudes give the corresponding probabilities ($`C_{1+}C_{1+}^{}`$ gives the probability for transmission, for example). $`s`$ is the spin of the particle (1 for photons and $`\frac{1}{2}`$ for the spin-$`\frac{1}{2}`$ singlet state- see below)
2. The amplitude correlation function is a normalized inner product of the amplitudes. This is of the form
$$U(𝐚,𝐛)=Real(NC_iC_j^{})$$
(4)
where $`N`$ is the normalization constant. The square of the amplitude correlation function gives the joint probabilities for events of the form $`(++),`$ $`(),`$ $`(+),`$ and $`(+).`$ All probabilities are guaranteed to be positive definite in our formalism since the amplitude correlation function is real.
The crucial difference from local realistic theories is that the correlation is calculated from quantities which preserve the relative phases.
Let us consider the maximally entangled singlet system described by Eq. 1, the most widely discussed example in the context of nonlocality. We prescribe the local amplitudes as $`C_{1+}=\frac{1}{\sqrt{2}}\mathrm{exp}\{is(\theta _1\varphi _1)\}`$ for the first particle at the first polarizer and $`C_{2+}=`$ $`\frac{1}{\sqrt{2}}\mathrm{exp}\{is(\theta _2\varphi _2)\}`$ for the second particle at the second polarizer. There are corresponding amplitudes, $`C_1`$ and $`C_2`$ for the events denoted by $`,`$ and they differ only in the phase for the maximally entangled state.
The explicit dependence of the amplitude on the spin of the particle is motivated by the fact that we are dealing with systems with phases and the phase associated with the spin rotations (a geometric phase) is a necessary input in this description . The correlation at source is encoded in $`\varphi _0`$. The locality assumption is strictly enforced since the two amplitudes depend only on local variables and on an internal variable generated at the source and then individually carried by the particles without any subsequent interaction of any sort. The individual measurements at each end separately will now give the correct result for transmission for any angle of orientation. These probabilities are
$$C_1C_1^{}=C_2C_2^{}=\frac{1}{2}$$
(5)
Events of both types ($`++`$) and ($``$) contribute to a “coincidence”. The correlation function for an outcome of either $`(++)`$ or $`()`$ of two maximally entangled particles is
$$U(\theta _1,\theta _2,\varphi _o)=2Re(C_1C_2^{})=\mathrm{cos}\{s(\theta _1\theta _2)+s\varphi _o\}.$$
(6)
It is normalized such that its square will give the conditional joint probabilities of the type ‘outcome $`+`$ for the second particle, given that the outcome for the first particle is $`+,`$ etc. All references to the individual values of the internal variable $`\varphi `$ has dropped out.
We now derive the relation between this correlation function and the experimenter’s correlation function $`P(𝐚,𝐛)=\frac{1}{N}(A_iB_i)`$. Since $`U_{++}^2=U_{}^2`$ for the maximally entangled state, $`U^2(\theta _1,\theta _2,\varphi _o)`$ is the probability for a coincidence detection ($`++`$ or $``$), and $`(1U^2(\theta _1,\theta _2,\varphi _o))`$ is the probability for an anticoincidence (events of the type $`+`$ and $`+`$). Since the average of the quantity (number of coincidences $``$ number of anticoincidences) =
$$U^2(\theta _1,\theta _2,\varphi _o)(1U^2(\theta _1,\theta _2,\varphi _o))=2U^2(\theta _1,\theta _2,\varphi _o)1,$$
(7)
the correspondence between $`P(𝐚,𝐛)`$ and $`U(\theta _1,\theta _2,\varphi _o)`$ is given by the expression,
$`P(𝐚,𝐛)`$ $`=`$ $`2U^2(\theta _1,\theta _2,\varphi _o)1`$ (8)
$`=`$ $`2\mathrm{cos}^2\{s(\theta _1\theta _2)+s\varphi _o\}1`$
This completes the back-bone of our formalism and we are ready to discuss some specific examples.
## 3 Spin-$`\frac{1}{2}`$ particles and Photons
Consider the singlet state breaking up into two spin-$`\frac{1}{2}`$ particles propagating in opposite directions to spatially separated regions. Since orthogonality of the two particles in any basis implies a relative angle of $`\pi `$ for spinors, we set $`\varphi _o=\pi `$ . Then the correlation function and $`P(𝐚,𝐛)`$ calculated from this function are
$`U(\theta _1,\theta _2,\varphi _o)`$ $`=`$ $`\mathrm{cos}\{s(\theta _1\theta _2)+s\varphi _o\}`$ (9)
$`=`$ $`\mathrm{cos}\{{\displaystyle \frac{1}{2}}(\theta _1\theta _2)+\pi /2\}`$
$`=`$ $`\mathrm{sin}{\displaystyle \frac{1}{2}}(\theta _1\theta _2)`$
$`P(𝐚,𝐛)`$ $`=`$ $`2\mathrm{sin}^2({\displaystyle \frac{1}{2}}(\theta _1\theta _2))1`$ (10)
$`=`$ $`\mathrm{cos}(\theta _1\theta _2)=𝐚𝐛`$
This is identical to the quantum mechanical predictions obtained from the singlet entangled state and Pauli spin operators. We have reproduced the correct correlation function using local amplitudes.
For the case of photons entangled in orthogonal polarization states we get, by setting $`s=1`$ and $`\varphi _o=\pi /2`$ to represent orthogonal polarization,
$`U(\theta _1,\theta _2,\varphi _o)`$ $`=`$ $`\mathrm{cos}\{(\theta _1\theta _2)+\pi /2\}`$ (11)
$`=`$ $`\mathrm{sin}(\theta _1\theta _2)`$
$$P(𝐚,𝐛)=2\mathrm{sin}^2(\theta _1\theta _2)1=\mathrm{cos}(2((\theta _1\theta _2))$$
(12)
which is the correct quantum mechanical correlation.
The same analysis works for particles entangled in other sets of variables like momentum and coordinate, and energy and time. These cases of two particle entanglement can be mapped on to the spin-$`\frac{1}{2}`$ singlet problem with two-valued outcomes. Starting from the local amplitudes $`C_1=\frac{1}{\sqrt{2}}\mathrm{exp}(i\alpha k(x_1x_o)/2),`$ and $`C_2=\frac{1}{\sqrt{2}}\mathrm{exp}(i\alpha k(x_2x_o)/2)`$ we can derive the probability for coincidence detection as
$$P(x_1,x_2)=\mathrm{cos}^2(\alpha k(x_1,x_2)/2)=\frac{1}{2}(1+\mathrm{cos}k\alpha (x_1x_2))$$
(13)
This is the two photon correlation pattern with 100% visibility, obtained without nonlocality. $`x_1`$ and $`x_2`$ are the coordinates of the two detectors separated by a space-like interval. $`k`$ is the wave vector and $`\alpha `$ is a scaling factor for the angle subtended by the two slits at the detectors, source etc. The factor $`2`$ dividing the angular variable comes from the mapping with the spin-$`\frac{1}{2}`$ problem.
## 4 Three-particle GHZ correlations
The three particle G-H-Z state is defined as
$$|\mathrm{\Psi }_{GHZ}=\frac{1}{\sqrt{2}}(|1,1,1|1,1,1)$$
(14)
where the eigenvalues in the kets are with respect to the $`z`$-axis basis.
The prediction from quantum mechanics for the measurement represented by the operator $`\sigma _x^1\sigma _x^2\sigma _x^3`$ is given by
$$\sigma _x^1\sigma _x^2\sigma _x^3|\mathrm{\Psi }_{GHZ}=|\mathrm{\Psi }_{GHZ}$$
(15)
Equivalently the joint probabilities for various outcomes in the $`x`$ direction are
$`P(+,+,+)`$ $`=`$ $`P(,,+)=P(+,,)=P(,+,)=0`$ (16)
$`P(,,)`$ $`=`$ $`P(+,+,)=P(+,,+)=P(,+,+)=1`$ (17)
Local realistic theories predict that the product of the outcomes in the $`x`$ direction for the three particles should be $`+1,`$ i.e.e,
$$P(+,+,+)=P(,,+)=P(+,,)=P(,+,)=1$$
This contradicts Eqs. 15-17 and highlights the conflict between a local realistic theory and quantum mechanics.
The solution using local amplitudes is simple and physically revealing . We define the local amplitudes for the outcomes $`+`$ and $``$ at the analyzer (with respect to the $`x`$ basis) for the first particle as $`C_{1+}=\frac{1}{\sqrt{2}}\mathrm{exp}(i\theta _1)`$, and $`C_1=\frac{1}{\sqrt{2}}\mathrm{exp}(i(\theta _1+\pi /2)).`$ The amplitude $`C_1`$ contains the added angle $`\pi /2`$ because this amplitude is orthogonal to $`C_{1+}.`$ Similarly, we have $`C_{2+}=\frac{1}{\sqrt{2}}\mathrm{exp}(i\theta _2)`$, and $`C_2=\frac{1}{\sqrt{2}}\mathrm{exp}(i(\theta _2+\pi /2))`$ for the second particle and $`C_{3+}=\frac{1}{\sqrt{2}}\mathrm{exp}(i\theta _3)`$, and $`C_3=\frac{1}{\sqrt{2}}\mathrm{exp}(i(\theta _3+\pi /2))`$ for the third particle.
Correlation function is obtained from $`N`$Real($`C_1C_2^{}C_3^{}),`$ where $`N`$ is a normalization constant, and its square is the relevant joint probability. (There is no unique definition of the amplitude correlation function. The final results are independent of the particular definition we use). Since we want $`N`$Re$`(C_1C_2^{}C_3^{})=\pm 1,`$ we choose $`C_1C_2^{}C_3^{}`$ to be pure real. This gives
$$\frac{N}{2\sqrt{2}}Real(\mathrm{exp}i(\theta _1\theta _2\theta _3\pi /2))=\pm 1$$
$$\theta _1\theta _2\theta _3\pi /2=0\mathrm{or}\pm \pi $$
We can choose the relevant relative phases to satisfy this condition. Then we get
$$P(,,)=1$$
Rest of the joint probabilities given in Eq. 6 automatically follow, since flipping sign once rotates the complex number $`C_1C_2^{}C_3^{}`$ through $`\pi /2.`$ The square of $`N`$Real($`C_1C_2^{}C_3^{})`$ is then $`1`$ for an odd number of $`()`$ outcomes and $`0`$ for even number of $`()`$ outcomes.
Similar construction also applies to four- particle maximally entangled state and general multiparticle maximally entangled states.
## 5 Concluding remarks
We have also constructed local amplitudes for the Hardy experiment in which quantum mechanics predicts three particular zero joint probabilities are one nonzero joint probability (the other possible joint probabilities in the problem can be nonzero and are not relevant for the demonstration of nonlocality). Local complex amplitudes that reproduce the four relevant joint probabilities can be constructed easily. It is impossible to achieve this if local realism at the level of eigenvalues are assumed.
We note that there is a simple way to physically understand the fact that the quantum correlations are typically larger than the corresponding classical correlations. The overlap (inner product) between a normalized random vector with $`N`$ elements and any basis vector is $`1/N`$ for a classical vector, and $`1/\sqrt{N}`$ for a quantum vector (amplitude). Therefore quantum correlations are typically stronger than classical correlations. In fact, it is this same physical fact that forms the basis of quantum search algorithms , where the initial overlap between a random vector and the desired basis vector is $`\sqrt{N}`$ times larger in the quantum case, making the search faster by $`\sqrt{N}.`$
The following table summarizes the locality and reality properties in various approaches to quantum correlations:
| Theory/ Formalism | Basic quantity | Locality | Reality | Determinism | Predictions |
| --- | --- | --- | --- | --- | --- |
| Quantum mechanics | Multiparticle wavefunction | NO | NO | NO | Correct |
| Local Realistic theories with hidden variables | Eigen values | YES | YES | YES | Incorrect |
| Present formalism | Amplitudes | YES | Yes (for phase) | NO | Correct |
Quantum entanglement swapping is understood within this frame work by noting that Bell state measurements choose subensembles of particle pairs that show a particular joint outcome. Particles entangled independently with the pair of particles that are subjected to the Bell state measurement will show a joint outcome consistent with swapped entanglement due to the correlation encoded in the internal variable. But the Bell state measurement does not collapse the distant particle into a definite state. Yet all correlations are correctly reproduced. This has important implication to the interpretation of quantum teleportation. The present nonlocal interpretation of quantum teleportation is not correct.
In summary, the long standing puzzle of nonlocality in the EPR correlations is resolved. There is no nonlocal influence between correlated particles separated into space-like regions. The solution has new physical and philosophical implications regarding the nature of reality, measurement and state reduction in quantum systems.
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# Curvature and Smooth Topology in Dimension Four
## 1 Four-Dimensional Geometry
Let $`M`$ be a smooth compact oriented 4-manifold. If $`g`$ is any Riemannian metric on $`M`$, the middle-dimensional Hodge star operator
$$:\mathrm{\Lambda }^2\mathrm{\Lambda }^2$$
has eigenvalues $`\pm 1`$. This gives rise to a natural decomposition
$$\mathrm{\Lambda }^2=\mathrm{\Lambda }^+\mathrm{\Lambda }^{}$$
(1)
of the rank-6 bundle of 2-forms into two rank-3 bundles, where
$$\psi \mathrm{\Lambda }^\pm \psi =\pm \psi .$$
Sections of $`\mathrm{\Lambda }^+`$ are called self-dual 2-forms, while sections of $`\mathrm{\Lambda }^{}`$ are called anti-self-dual 2-forms. The middle-dimensional Hodge star operator is unchanged if $`g`$ is multiplied by a smooth positive function, so the decomposition (1) really only depends on the conformal class $`\gamma =[g]`$ rather than on the Riemannian metric itself.
The decomposition (1) has important ramifications for Riemannian geometry. In particular, since the curvature tensor $``$ may be thought of as a linear map $`\mathrm{\Lambda }^2\mathrm{\Lambda }^2`$, there is an induced decomposition
=(
+W+s12rr+W-s12 )
+W+s12rr+W-s12 {\mathcal{R}}=\left(\mbox{
\begin{tabular}[]{c|c}&\\
$W_{+}+\frac{s}{12}$&$\stackrel{{\scriptstyle\circ}}{{r}}$\\
&\\
\cline{1-2}\cr&\\
$\stackrel{{\scriptstyle\circ}}{{r}}$&$W_{-}+\frac{s}{12}$\\
&\\
\end{tabular}
}\right)
into simpler curvature tensors. Here the self-dual and anti-self-dual Weyl curvatures $`W_\pm `$ are the trace-free pieces of the appropriate blocks. The scalar curvature $`s`$ is understood to act by scalar multiplication, and $`\stackrel{}{r}`$ can be identified with the trace-free part $`r\frac{s}{4}g`$ of the Ricci curvature.
Now the $`L^2`$-norm of each of these curvatures are all scale invariant, so one may define sensible diffeomorphism invariants of a 4-manifold, such as
$$(infs)(M)=\underset{g}{inf}\left(_Ms_g^2𝑑\mu _g\right)^{1/2}$$
by considering the infima of the $`L^2`$ norms of these curvatures over all metrics $`g`$ on $`M`$. Notice that these invariants might a priori depend on the differentiable structure of $`M`$. In any case, there is no obvious way to read off these invariants from homotopy invariants of the $`M`$.
Nonetheless, homotopy invariants do impose some important relations between these invariants. One such relation arises from the intersection form
$`:H^2(M,)\times H^2(M,)`$ $``$ $``$
$`([\varphi ],[\psi ])`$ $``$ $`{\displaystyle _M}\varphi \psi `$
which may be diagonalized as
$$\left[\begin{array}{cc}\hfill \underset{b_+(M)}{\underset{}{\begin{array}{ccc}1& & \\ & \mathrm{}& \\ & & 1\end{array}}}& \\ \hfill b_{}(M)\{\begin{array}{c}\\ \\ \end{array}& \begin{array}{ccc}1& & \\ & \mathrm{}& \\ & & 1\end{array}\hfill \end{array}\right]$$
by choosing a suitable basis for the de Rham cohomology $`H^2(M,)`$. The numbers $`b_\pm (M)`$ are independent of the choice of basis, and so are oriented homotopy invariants of $`M`$. Their difference
$$\tau (M)=b_+(M)b_{}(M),$$
is called the signature of $`M`$. The Hirzebruch signature theorem asserts that this invariant is expressible as a curvature integral:
$$\tau (M)=\frac{1}{12\pi ^2}_M\left(|W^+|^2|W^{}|^2\right)𝑑\mu .$$
(2)
Here the curvatures, norms $`||`$, and volume form $`d\mu `$ are, of course, those of the Riemannian metric $`g`$, but the entire point is that the answer is independent of which metric we use. Thus $`(infW_+)^2=(infW_{})^2+12\pi ^2\tau (M)`$.
Another relationship between the curvature L<sup>2</sup>-norms under consideration is given by the $`4`$-dimensional case of the generalized Gauss-Bonnet theorem; this asserts that the Euler characteristic
$$\chi (M)=22b_1(M)+b_2(M)$$
is given by
$$\chi (M)=\frac{1}{8\pi ^2}_M\left(|W^+|^2+|W^{}|^2+\frac{s^2}{24}\frac{|\stackrel{}{r}|^2}{2}\right)𝑑\mu .$$
(3)
This then gives rise to inequalities concerning $`infW_\pm `$, $`infs`$, and $`inf\stackrel{}{r}`$ which are imposed by the homotopy type of $`M`$.
Not long ago, the problem of calculating invariants such as $`infs`$ would have simply been considered intractable. However, there has been a remarkable amount of recent progress on these issues. The key event in this regard was Witten’s introduction of the so-called Seiberg-Witten invariants, which display an unexpected relationship between Donaldson’s polynomial invariants and Riemannian geometry on $`M`$. In particular, $`infs`$ turns out not to be a homeomorphism invariant, but is nonetheless exactly calculable for a huge class of 4-manifolds, including all complex algebraic surfaces . In the next section, we will give a brief introduction to Seiberg-Witten theory, and then show how it allows one to estimate certain linear combinations of $`s`$ and $`W_+`$.
## 2 Seiberg-Witten Theory
Let $`(M,g)`$ be a compact oriented Riemannian $`4`$-manifold. On any contractible open subset $`UM`$, one can define Hermitian vector bundles
$$\begin{array}{cc}\hfill ^2& 𝕊_\pm |_U\hfill \\ & \hfill \\ & UM\hfill \end{array}$$
called spin bundles, characterized by the fact that their determinant line bundles $`^2𝕊_\pm `$ are canonically trivial and that their projectivizations
$$\begin{array}{cc}\hfill _1& (𝕊_\pm )\\ & \\ & M\end{array}$$
are exactly the unit 2-sphere bundles $`S(\mathrm{\Lambda }^\pm )`$. As one passes between open subset $`U`$ and $`U^{}`$, however, the corresponding locally-defined spin bundles are not quite canonically isomorphic; instead, there are two equally ‘canonical’ isomorphisms, differing by a sign. Because of this, one cannot generally define the bundles $`𝕊_\pm `$ globally on $`M`$; manifolds on which this can be done are called spin, and are characterized by the vanishing of the Stiefel-Whitney class $`w_2=w_2(TM)H^2(M,_2)`$. However, one can always find Hermitian complex line bundles $`LM`$ with first Chern class $`c_1=c_1(L)`$ satisfying
$$c_1w_2mod2.$$
(4)
Given such a line bundle, one can then construct Hermitian vector bundles $`𝕍_\pm `$ with
$$(𝕍_\pm )=S(\mathrm{\Lambda }^\pm )$$
by formally setting
$$𝕍_\pm =𝕊_\pm L^{1/2},$$
because the sign problems encountered in consistently defining the transition functions of $`𝕊_\pm `$ are exactly canceled by those associated with trying to find consistent square-roots of the transition functions of $`L`$.
The isomorphism class $`𝔠`$ of such a choice of $`𝕍_\pm `$ is called a spin<sup>c</sup> structure on $`M`$. The cohomology group $`H^2(M,)`$ acts freely and transitively on the spin<sup>c</sup> structures by tensoring $`𝕍_\pm `$ with complex line bundles. Each spin<sup>c</sup> structure has a first Chern class $`c_1:=c_1(L)=c_1(𝕍_\pm )H^2(M,)`$ satisfying (4), and the $`H^2(M,)`$-action on spin<sup>c</sup> structures induces the action
$$c_1c_1+2\alpha ,$$
$`\alpha H^2(M,)`$, on first Chern classes. Thus, if $`H^2(M,)`$ has trivial 2-torsion — as can always be arranged by replacing $`M`$ with a finite cover — the spin<sup>c</sup> structures are precisely in one-to-one correspondence with the set of cohomology classes $`c_1H^2(M,)`$ satisfying (4).
To make this discussion more concrete, suppose that $`M`$ admits an almost-complex structure. Any given almost-complex structure can be deformed to an almost complex structures $`J`$ which is compatible with $`g`$ in the sense that $`J^{}g=g`$. Choose such a $`J`$, and consider the rank-2 complex vector bundles
$`𝕍_+`$ $`=`$ $`\mathrm{\Lambda }^{0,0}\mathrm{\Lambda }^{0,2}`$ (5)
$`𝕍_{}`$ $`=`$ $`\mathrm{\Lambda }^{0,1}.`$
These are precisely the twisted spinor bundles of the spin<sup>c</sup> structure obtained by taking $`L`$ to be the anti-canonical line bundle $`\mathrm{\Lambda }^{0,2}`$ of the almost-complex structure. A spin<sup>c</sup> structure $`𝔠`$ arising in this way will be said to be of almost-complex type. These are exactly the spin<sup>c</sup> structures for which
$$c_1^2=(2\chi +3\tau )(M).$$
On a spin manifold, the spin bundles $`𝕊_\pm `$ carry natural connections induced by the Levi-Civita connection of the given Riemannian metric $`g`$. On a spin<sup>c</sup> manifold, however, there is not a natural unique choice of connections on $`𝕍_\pm `$. Nonetheless, since we formally have $`𝕍_\pm =𝕊_\pm L^{1/2}`$, every Hermitian connection $`A`$ on $`L`$ induces associated Hermitian connections $`_A`$ on $`𝕍_\pm `$.
On the other hand, there is a canonical isomorphism $`\mathrm{\Lambda }^1=\text{Hom }(𝕊_+,𝕊_{})`$, so that $`\mathrm{\Lambda }^1\text{Hom }(𝕍_+,𝕍_{})`$ for any spin<sup>c</sup> structure, and this induces a canonical homomorphism
$$:\mathrm{\Lambda }^1𝕍_+𝕍_{}$$
called Clifford multiplication. Composing these operations allows us to define a so-called twisted Dirac operator
$$D_A:\mathrm{\Gamma }(𝕍_+)\mathrm{\Gamma }(𝕍_{})$$
by $`D_A\mathrm{\Phi }=_A\mathrm{\Phi }`$.
For any spin<sup>c</sup> structure, we have already noted that there is a canonical diffeomorphism $`(𝕍_+)\stackrel{}{}S(\mathrm{\Lambda }^+)`$. In polar coordinates, we now use this to define the angular part of a unique continuous map
$$\sigma :𝕍_+\mathrm{\Lambda }^+$$
with
$$|\sigma (\mathrm{\Phi })|=\frac{1}{2\sqrt{2}}|\mathrm{\Phi }|^2.$$
This map is actually real-quadratic on each fiber of $`𝕍_+`$; indeed, assuming our spin<sup>c</sup> structure is induced by a complex structure $`J`$, then, in terms of (5), $`\sigma `$ is explicitly given by
$$\sigma (f,\varphi )=(|f|^2|\varphi |^2)\frac{\omega }{4}+\mathrm{}m(\overline{f}\varphi ),$$
where $`f\mathrm{\Lambda }^{0,0}`$, $`\varphi \mathrm{\Lambda }^{0,2}`$, and where $`\omega (,)=g(J,)`$ is the associated 2-form of $`(M,g,J)`$. On a deeper level, $`\sigma `$ directly arises from the fact that $`𝕍_+=𝕊_+L^{1/2}`$, while $`\mathrm{\Lambda }^+=^2𝕊_+`$. For this reason, $`\sigma `$ is is invariant under parallel transport.
We are now in a position to introduce the Seiberg-Witten equations
$`D_A\mathrm{\Phi }`$ $`=`$ $`0`$ (6)
$`F_A^+`$ $`=`$ $`i\sigma (\mathrm{\Phi }),`$ (7)
where the unknowns are a Hermitian connection $`A`$ on $`L`$ and a section $`\mathrm{\Phi }`$ of $`𝕍_+`$. Here $`F_A^+`$ is the self-dual part of the curvature of $`A`$, and so is a purely imaginary 2-form.
For many $`4`$-manifolds, it turns out that there is a solution of the Seiberg-Witten equations for each metric. Let us introduce some convenient terminology to describe this situation.
###### Definition 2.1
Let $`M`$ be a smooth compact oriented $`4`$-manifold with $`b_+2`$, and suppose that $`M`$ carries a spin<sup>c</sup> structure $`𝔠`$ for which the Seiberg-Witten equations (67) have a solution for every Riemannian metric $`g`$ on $`M`$. Then the first Chern class $`c_1H^2(M,)`$ of $`𝔠`$ will be called a monopole class.
This definition is useful in practice primarily because there are topological arguments which lead to the existence of solutions the Seiberg-Witten equations. For example , if $`𝔠`$ is a spin<sup>c</sup> structure of almost-complex type, then the Seiberg-Witten invariant $`𝒮𝒲_𝔠(M)`$ can be defined as the number of solutions, modulo gauge transformations and counted with orientations, of a generic perturbation
$`D_A\mathrm{\Phi }`$ $`=`$ $`0`$
$`iF_A^++\sigma (\mathrm{\Phi })`$ $`=`$ $`\varphi `$
of (67), where $`\varphi `$ is a smooth self-dual 2-form. If $`b_+(M)2`$, this integer is independent of the metric $`g`$; and if it is non-zero, the first Chern class $`c_1`$ of $`𝔠`$ is then a monopole class. Similar things are also true when $`b_+(M)=1`$, although the story becomes rather more complicated.
Now, via the Hodge theorem, every Riemannian metric $`g`$ on $`M`$ determines a direct sum decomposition
$$H^2(M,)=_g^+_g^{},$$
where $`_g^+`$ (respectively, $`_g^{}`$) consists of those cohomology classes for which the harmonic representative is self-dual (respectively, anti-self-dual). Because the restriction of the intersection form to $`_g^+`$ (respectively, $`_g^{}`$) is positive (respectively, negative) definite, and because these subspaces are mutually orthogonal with respect to the intersection pairing, the dimensions of these spaces are exactly the invariants $`b_\pm `$ defined in §1. If the first Chern class $`c_1`$ of the spin<sup>c</sup> structure $`𝔠`$ is now decomposed as
$$c_1=c_1^++c_1^{},$$
where $`c_1^\pm _g^\pm `$, we get the important inequality
$$_M|\mathrm{\Phi }|^4𝑑\mu 32\pi ^2(c_1^+)^2$$
(8)
because (7) tells us that $`2\pi c_1^+`$ is the harmonic part of $`\sigma (\mathrm{\Phi })`$.
Many of the most remarkable consequences of Seiberg-Witten theory stem from the fact that the equations (67) imply the Weitzenböck formula
$$0=4^{}\mathrm{\Phi }+s\mathrm{\Phi }+|\mathrm{\Phi }|^2\mathrm{\Phi },$$
(9)
where $`s`$ denotes the scalar curvature of $`g`$, and where we have introduced the abbreviation $`_A=`$. Taking the inner product with $`\mathrm{\Phi }`$, it follows that
$$0=2\mathrm{\Delta }|\mathrm{\Phi }|^2+4|\mathrm{\Phi }|^2+s|\mathrm{\Phi }|^2+|\mathrm{\Phi }|^4.$$
(10)
If we multiply (10) by $`|\mathrm{\Phi }|^2`$ and integrate, we have
$$0=_M\left[2\left|d|\mathrm{\Phi }|^2\right|^2+4|\mathrm{\Phi }|^2|\mathrm{\Phi }|^2+s|\mathrm{\Phi }|^4+|\mathrm{\Phi }|^6\right]𝑑\mu _g,$$
so that
$$(s)|\mathrm{\Phi }|^4𝑑\mu 4|\mathrm{\Phi }|^2|\mathrm{\Phi }|^2𝑑\mu +|\mathrm{\Phi }|^6𝑑\mu .$$
(11)
This leads to the following curvature estimate:
###### Theorem 2.2
Let $`M`$ be a smooth compact oriented 4-manifold with monopole class $`c_1`$. Then every Riemannian metric $`g`$ on $`M`$ satisfies
$$_M\left(\frac{2}{3}s2\sqrt{\frac{2}{3}}|W_+|\right)^2𝑑\mu 32\pi ^2(c_1^+)^2,$$
(12)
where $`c_1^+`$ is the self-dual part of $`c_1`$ with respect to $`g`$.
Proof. The first step is to prove the inequality
$$V^{1/3}\left(_M\left|\frac{2}{3}s_g2\sqrt{\frac{2}{3}}|W_+|\right|^3𝑑\mu \right)^{2/3}32\pi ^2(c_1^+)^2,$$
(13)
where $`V=\text{Vol}(M,g)=_M𝑑\mu _g`$ is the total volume of $`(M,g)`$.
Any self-dual 2-form $`\psi `$ on any oriented 4-manifold satisfies the Weitzenböck formula
$$(d+d^{})^2\psi =^{}\psi 2W_+(\psi ,)+\frac{s}{3}\psi .$$
It follows that
$$_M(2W_+)(\psi ,\psi )𝑑\mu _M(\frac{s}{3})|\psi |^2𝑑\mu _M|\psi |^2𝑑\mu .$$
However,
$$|W_+|_g|\psi |^2\sqrt{\frac{3}{2}}W_+(\psi ,\psi )$$
simply because $`W_+`$ is trace-free. Thus
$$_M2\sqrt{\frac{2}{3}}|W_+||\psi |^2𝑑\mu _M(\frac{s}{3})|\psi |^2𝑑\mu _M|\psi |^2𝑑\mu ,$$
and hence
$$_M(\frac{2}{3}s2\sqrt{\frac{2}{3}}|W_+|)|\psi |^2𝑑\mu _M(s)|\psi |^2𝑑\mu _M|\psi |^2𝑑\mu .$$
On the other hand, the particular self-dual 2-form $`\phi =\sigma (\mathrm{\Phi })=iF_A^+`$ satisfies
$`|\phi |^2`$ $`=`$ $`{\displaystyle \frac{1}{8}}|\mathrm{\Phi }|^4,`$
$`|\phi |^2`$ $``$ $`{\displaystyle \frac{1}{2}}|\mathrm{\Phi }|^2|\mathrm{\Phi }|^2.`$
Setting $`\psi =\phi `$, we thus have
$$_M(\frac{2}{3}s2\sqrt{\frac{2}{3}}|W_+|)|\mathrm{\Phi }|^4𝑑\mu _M(s)|\mathrm{\Phi }|^4𝑑\mu 4_M|\mathrm{\Phi }|^2|\mathrm{\Phi }|^2𝑑\mu .$$
But (11) tells us that
$$_M(s)|\mathrm{\Phi }|^4𝑑\mu 4_M|\mathrm{\Phi }|^2|\mathrm{\Phi }|^2𝑑\mu _M|\mathrm{\Phi }|^6𝑑\mu ,$$
so we obtain
$$_M(\frac{2}{3}s2\sqrt{\frac{2}{3}}|W_+|)|\mathrm{\Phi }|^4𝑑\mu _M|\mathrm{\Phi }|^6𝑑\mu .$$
(14)
By the Hölder inequality, we thus have
$$\left(\left|\frac{2}{3}s2\sqrt{\frac{2}{3}}|W_+|\right|^3𝑑\mu \right)^{1/3}\left(|\mathrm{\Phi }|^6𝑑\mu \right)^{2/3}|\mathrm{\Phi }|^6𝑑\mu ,$$
Since the Hölder inequality also tells us that
$$|\mathrm{\Phi }|^6𝑑\mu V^{1/2}\left(|\mathrm{\Phi }|^4𝑑\mu \right)^{3/2},$$
we thus have
$$V^{1/3}\left(_M\left|\frac{2}{3}s2\sqrt{\frac{2}{3}}|W_+|\right|^3𝑑\mu \right)^{2/3}|\mathrm{\Phi }|^4𝑑\mu 32\pi ^2(c_1^+)^2,$$
where the last inequality is exactly (8). This completes the first part of the proof.
Next, we observe that any smooth conformal $`\gamma `$ class on any oriented $`4`$-manifold contains a $`C^2`$ metric such that $`s\sqrt{6}|W_+|`$ is constant. Indeed, as observed by Gursky , this readily follows from the standard proof of the Yamabe problem. The main point is that the curvature expression
$$𝔖_g=s_g\sqrt{6}|W_+|_g$$
transforms under conformal changes $`g\widehat{g}=u^2g`$ by the rule
$$𝔖_{\widehat{g}}=u^3\left(6\mathrm{\Delta }_g+𝔖_g\right)u,$$
just like the ordinary scalar curvature $`s`$. We will actually use this only in the negative case, where the proof is technically the simplest, and simply repeats<sup>1</sup><sup>1</sup>1However, since $`|W_+|`$ is generally only Lipschitz continuous, the minimizer generally only has regularity $`C^{2,\alpha }`$ in the vicinity of a zero of $`W_+`$. the arguments of Trudinger .
The conformal class $`\gamma `$ of a given metric $`g`$ thus always contains a metric $`g_\gamma `$ for which $`\frac{2}{3}s2\sqrt{\frac{2}{3}}|W_+|`$ is constant. But since the existence of solutions of the Seiberg-Witten equations precludes the possibility that we might have $`s_{g_\gamma }>0`$, this constant is necessarily non-positive. We thus have
$$_M(\frac{2}{3}s_{g_\gamma }2\sqrt{\frac{2}{3}}|W_+|_{g_\gamma })^2d\mu _{g_\gamma }=V_{g_\gamma }^{1/3}\left(_M\right|(\frac{2}{3}s_{g_\gamma }2\sqrt{\frac{2}{3}}|W_+|_{g_\gamma }|^3d\mu _{g_\gamma })^{2/3},$$
so that
$$_M\left(\frac{2}{3}s_{g_\gamma }2\sqrt{\frac{2}{3}}|W_+|_{g_\gamma }\right)^2𝑑\mu _{g_\gamma }32\pi ^2(c_1^+)^2.$$
Thus we at least have the desired $`L^2`$ estimate for a specific metric $`g_\gamma `$ which is conformally related to the given metric $`g`$.
Let us now compare the left-hand side with analogous expression for the given metric $`g`$. To do so, we express $`g`$ in the form $`g=u^2g_\gamma `$, where $`u`$ is a positive $`C^2`$ function, and observe that
$`{\displaystyle _M}\left({\displaystyle \frac{2}{3}}s_g2\sqrt{{\displaystyle \frac{2}{3}}}|W_+|_g\right)u^2𝑑\mu _{g_\gamma }`$ $`=`$ $`{\displaystyle \frac{2}{3}}{\displaystyle 𝔖_gu^2𝑑\mu _{g_\gamma }}`$
$`=`$ $`{\displaystyle \frac{2}{3}}{\displaystyle u^3\left(6\mathrm{\Delta }_{g_\gamma }u+𝔖_{g_\gamma }u\right)u^2𝑑\mu _{g_\gamma }}`$
$`=`$ $`{\displaystyle \frac{2}{3}}{\displaystyle \left(6u^2|du|_{g_\gamma }^2+𝔖_{g_\gamma }\right)𝑑\mu _{g_\gamma }}`$
$``$ $`{\displaystyle \frac{2}{3}}{\displaystyle 𝔖_{g_\gamma }𝑑\mu _{g_\gamma }}`$
$`=`$ $`{\displaystyle _M}\left({\displaystyle \frac{2}{3}}s_{g_\gamma }2\sqrt{{\displaystyle \frac{2}{3}}}|W_+|_{g_\gamma }\right)𝑑\mu _{g_\gamma }.`$
Applying Cauchy-Schwarz, we thus have
$`V_{g_\gamma }^{1/2}\left({\displaystyle \left(\frac{2}{3}s_g2\sqrt{\frac{2}{3}}|W_+|_g\right)^2𝑑\mu _g}\right)^{1/2}`$ $``$ $`{\displaystyle _M}\left({\displaystyle \frac{2}{3}}s_g2\sqrt{{\displaystyle \frac{2}{3}}}|W_+|_g\right)u^2𝑑\mu _{g_\gamma }`$
$``$ $`{\displaystyle _M}\left({\displaystyle \frac{2}{3}}s_{g_\gamma }2\sqrt{{\displaystyle \frac{2}{3}}}|W_+|_{g_\gamma }\right)𝑑\mu _{g_\gamma }`$
$`=`$ $`V_{g_\gamma }^{1/2}\left({\displaystyle \left(\frac{2}{3}s_{g_\gamma }2\sqrt{\frac{2}{3}}|W_+|_{g_\gamma }\right)^2𝑑\mu _{g_\gamma }}\right)^{1/2},`$
and hence
$$_M\left(\frac{2}{3}s_g2\sqrt{\frac{2}{3}}|W_+|_g\right)^2𝑑\mu _g_M\left(\frac{2}{3}s_{g_\gamma }2\sqrt{\frac{2}{3}}|W_+|_{g_\gamma }\right)^2𝑑\mu _{g_\gamma },$$
exactly as claimed.
Notice that we can rewrite the inequality (12) as
$$\frac{2}{3}s2\sqrt{\frac{2}{3}}|W_+|4\sqrt{2}\pi |c_1^+|,$$
where $``$ denotes the $`L^2`$ norm with respect to $`g`$. Dividing by $`\sqrt{24}`$ and applying the triangle inequality, we thus have
###### Corollary 2.3
Let $`M`$ be a smooth compact oriented 4-manifold with monopole class $`c_1`$. Then every Riemannian metric $`g`$ on $`M`$ satisfies
$$\frac{2}{3}\frac{s}{\sqrt{24}}+\frac{1}{3}W_+\frac{2\pi }{\sqrt{3}}|c_1^+|.$$
(15)
Inequality (12) actually belongs to a family of related estimates:
###### Theorem 2.4
Let $`M`$ be a smooth compact oriented 4-manifold with monopole class $`c_1`$, and let $`\delta [0,\frac{1}{3}]`$ be a constant. Then every Riemannian metric $`g`$ on $`M`$ satisfies
$$_M\left[(1\delta )s\delta \sqrt{24}|W_+|\right]^2𝑑\mu 32\pi ^2(c_1^+)^2,$$
(16)
Proof. Inequality (11) implies
$$(s)|\mathrm{\Phi }|^4𝑑\mu |\mathrm{\Phi }|^6𝑑\mu .$$
(17)
On the other hand, inequality (14) asserts that
$$_M(\frac{2}{3}s2\sqrt{\frac{2}{3}}|W_+|)|\mathrm{\Phi }|^4𝑑\mu _M|\mathrm{\Phi }|^6𝑑\mu .$$
Now multiply (17) by $`13\delta `$, multiply (14) by $`3\delta `$, and add. The result is
$$\left[(1\delta )s\delta \sqrt{24}|W_+|\right]|\mathrm{\Phi }|^4𝑑\mu |\mathrm{\Phi }|^6𝑑\mu .$$
(18)
Applying the same Hölder inequalities as before, we now obtain
$$V^{1/3}\left(_M\left|(1\delta )s\delta \sqrt{24}|W_+|\right|^3𝑑\mu \right)^{2/3}|\mathrm{\Phi }|^4𝑑\mu 32\pi ^2(c_1^+)^2.$$
Passage from this $`L^3`$ estimate to the desired $`L^2`$ estimate is then accomplished by the same means as before: every conformal class contains a metric for which $`(1\delta )s\delta \sqrt{24}|W_+|`$ is constant, and this metric minimizes
$$_M\left[(1\delta )s\delta \sqrt{24}|W_+|\right]^2𝑑\mu $$
among metrics in its conformal class.
Rewriting (16) as
$$(1\delta )s\delta \sqrt{24}|W_+|4\sqrt{2}\pi |c_1^+|,$$
dividing by $`\sqrt{24}`$, and applying the triangle inequality, we thus have
###### Corollary 2.5
Let $`M`$ be a smooth compact oriented 4-manifold with monopole class $`c_1`$. Then every Riemannian metric $`g`$ on $`M`$ satisfies
$$(1\delta )\frac{s}{\sqrt{24}}+\delta W_+\frac{2\pi }{\sqrt{3}}|c_1^+|$$
(19)
for every $`\delta [0,\frac{1}{3}]`$.
The $`\delta =0`$ version of (16) is implicit in the work of Witten ; it was later made explicit in , where it was also shown that equality holds for $`\delta =0`$ iff $`g`$ is a Kähler metric of constant, non-positive scalar curvature. But indeed, since $`\sqrt{24}|W_+||s|`$ for any Kähler manifold of real dimension $`4`$, metrics of this kind saturate (16) for each value of $`\delta `$. Conversely:
###### Proposition 2.6
Let $`\delta [0,\frac{1}{3})`$ be a fixed constant. If $`g`$ is a metric such that equality holds in (16), then $`g`$ is Kähler, and has constant scalar curvature.
Proof. Equality in (16) implies equality in (18). However, $`(13\delta )`$ times inequality (11) plus $`3\delta `$ times inequality (14) reads
$$\left[(1\delta )s\delta \sqrt{24}|W_+|\right]|\mathrm{\Phi }|^4𝑑\mu |\mathrm{\Phi }|^6𝑑\mu +4(13\delta )|\mathrm{\Phi }|^2|\mathrm{\Phi }|^2𝑑\mu .$$
Equality in (16) therefore implies that
$$0=\frac{1}{2}|\mathrm{\Phi }|^2|\mathrm{\Phi }|^2𝑑\mu |\phi |^2𝑑\mu ,$$
forcing the $`2`$-form $`\phi `$ to be parallel. If $`\phi 0`$, we conclude that the metric is Kähler, and the constancy of $`s`$ then follows from the Yamabe portion of the argument.
On the other hand, since $`b_+(M)2`$ and $`c_1`$ is a monopole class, $`M`$ does not admit any metrics of positive scalar curvature. If $`\phi 0`$ and (16) is saturated, one can therefore show that $`(M,g)`$ is $`K3`$ or $`T^4`$ with a Ricci-flat Kähler metric. The details are left as an exercise for the interested reader.
When $`\delta =\frac{1}{3}`$, the above argument breaks down. However, a metric $`g`$ can saturate (12) only if equality holds in (8), and this forces the self-dual $`2`$-form $`\phi =\sigma (\mathrm{\Phi })`$ to be harmonic. Moreover, the relevant Hölder inequalities would also have to be saturated, forcing $`\phi `$ to have constant length. This forces $`g`$ to be almost-Kähler, in the sense that there is an orientation-compatible orthogonal almost-complex structure for which the associated $`2`$-form is closed. For details, see .
It is reasonable to ask whether the inequalities (16) and (19) continue to hold when $`\delta >1/3`$. This issue will be addressed in §4.
## 3 Einstein Metrics
Recall that a smooth Riemannian metric $`g`$ is said to be Einstein if its Ricci curvature $`r`$ is a constant multiple of the metric:
$$r=\lambda g.$$
Not every 4-manifold admits such metrics. A necessary condition for the existence of an Einstein metric on a compact oriented 4-manifold is that the Hitchin-Thorpe inequality $`2\chi (M)3|\tau (M)|`$ must hold . Indeed, (2) and (3) tell us that
$$(2\chi \pm 3\tau )(M)=\frac{1}{4\pi ^2}_M\left(\frac{s^2}{24}+2|W_\pm |^2\frac{|\stackrel{}{r}|^2}{2}\right)𝑑\mu .$$
The Hitchin-Thorpe inequality follows, since the integrand is non-negative when $`\stackrel{}{r}=0`$. This argument, however, treats the scalar and Weyl contributions as ‘junk’ terms, about which one knows nothing except that they are non-negative. We now remedy this by invoking the estimates of §2.
###### Proposition 3.1
Let $`M`$ be a smooth compact oriented 4-manifold with monopole class $`c_1`$. Then every metric $`g`$ on $`M`$ satisfies
$$\frac{1}{4\pi ^2}_M\left(\frac{s_g^2}{24}+2|W_+|_g^2\right)𝑑\mu _g\frac{2}{3}(c_1^+)^2.$$
If $`c_1^+0`$, moreover, equality can only hold if $`g`$ is almost-Kähler, with almost-Kähler class proportional to $`c_1^+`$.
Proof. We begin begin with inequality (15)
$$\frac{2}{3}\frac{s}{\sqrt{24}}+\frac{1}{3}W_+\frac{2\pi }{\sqrt{3}}|c_1^+|,$$
and elect to interpret the left-hand side as the dot product
$$(\frac{2}{3},\frac{1}{3\sqrt{2}})(\frac{s}{\sqrt{24}},\sqrt{2}W_+)$$
in $`^2`$. Applying Cauchy-Schwarz, we thus have
$$\left((\frac{2}{3})^2+(\frac{1}{3\sqrt{2}})^2\right)^{1/2}\left(_M(\frac{s^2}{24}+2|W_+|^2)𝑑\mu \right)^{1/2}\frac{2}{3}\frac{s}{\sqrt{24}}+\frac{1}{3}W_+.$$
Thus
$$\frac{1}{2}_M(\frac{s^2}{24}+2|W_+|^2)𝑑\mu \frac{4\pi ^2}{3}(c_1^+)^2,$$
and hence
$$\frac{1}{4\pi ^2}_M\left(\frac{s_g^2}{24}+2|W_+|_g^2\right)𝑑\mu _g\frac{2}{3}(c_1^+)^2,$$
as claimed.
In the equality case, $`\phi `$ would be a closed self-dual form of constant norm, so $`g`$ would be almost-Kähler unless $`\phi 0`$.
To give some concrete applications, we now focus on the case of complex surfaces.
###### Proposition 3.2
Let $`(X,J_X)`$ be a compact complex surface with $`b_+>1`$, and let $`(M,J_X)`$ be the complex surface obtained from $`X`$ by blowing up $`k>0`$ points. Then any Riemannian metric $`g`$ on the $`4`$-manifold
$$M=X\mathrm{\#}k\overline{}_2$$
satisfies
$$\frac{1}{4\pi ^2}_M\left(\frac{s_g^2}{24}+2|W_+|_g^2\right)𝑑\mu _g>\frac{2}{3}(2\chi +3\tau )(X).$$
Proof. Let $`c_1(X)`$ denote the first Chern class of the given complex structure $`J_X`$, and, by a standard abuse of notation, let $`c_1(X)`$ also denote the pull-back class of this class to $`M`$. If $`E_1,\mathrm{},E_k`$ are the Poincaré duals of the exceptional divisors in $`M`$ introduced by blowing up, the complex structure $`J_M`$ has Chern class
$$c_1(M)=c_1(X)\underset{j=1}{\overset{k}{}}E_j.$$
By a result of Witten , this is a monopole class of $`M`$. However, there are self-diffeomorphisms of $`M`$ which act on $`H^2(M)`$ in a manner such that
$`c_1(X)`$ $``$ $`c_1(X)`$
$`E_j`$ $``$ $`\pm E_j`$
for any choice of signs we like. Thus
$$c_1=c_1(X)+\underset{j=1}{\overset{k}{}}(\pm E_j)$$
is a monopole class on $`M`$ for each choice of signs. We now fix our choice of signs so that
$$[c_1(X)]^+(\pm E_j)0,$$
for each $`j`$, with respect to the decomposition induced by the given metric $`g`$. We then have
$`(c_1^+)^2`$ $`=`$ $`\left([c_1(X)]^++{\displaystyle \underset{j=1}{\overset{k}{}}}(\pm E_j^+)\right)^2`$
$`=`$ $`([c_1(X)]^+)^2+2{\displaystyle \underset{j=1}{\overset{k}{}}}[c_1(X)]^+(\pm E_j)+({\displaystyle \underset{j=1}{\overset{k}{}}}(\pm E_j^+))^2`$
$``$ $`([c_1(X)]^+)^2`$
$``$ $`(2\chi +3\tau )(X).`$
This shows that
$$\frac{1}{4\pi ^2}_M\left(\frac{s_g^2}{24}+2|W_+|_g^2\right)𝑑\mu _g\frac{2}{3}(2\chi +3\tau )(X).$$
If equality held, $`g`$ would be almost-Kähler, with almost-Kähler class $`[\omega ]`$ proportional to $`c_1^+`$. On the other hand, we would also have $`[c_1(X)]^+E_j=0`$, so it would then follow that $`[\omega ]E_j=0`$ for all $`j`$. However, the Seiberg-Witten invariant would be non-trivial for a spin<sup>c</sup> structure with $`c_1(\stackrel{~}{L})=c_1(L)2(\pm E_1)`$, and a celebrated theorem of Taubes would then force the homology class $`E_j`$ to be represented by a pseudo-holomorphic $`2`$-sphere in the symplectic manifold $`(M,\omega )`$. But the (positive!) area of this sphere with respect to $`g`$ would then be exactly $`[\omega ]E_j`$, contradicting the observation that $`[\omega ]E_j=0`$.
###### Theorem 3.3
Let $`(X,J_X)`$ be a compact complex surface with $`b_+>1`$, and let $`(M,J_M)`$ be obtained from $`X`$ by blowing up $`k`$ points. Then the smooth compact $`4`$-manifold $`M`$ does not admit any Einstein metrics if $`k\frac{1}{3}c_1^2(X)`$.
Proof. We may assume that $`(2\chi +3\tau )(X)>0`$, since otherwise the result follows from the Hitchin-Thorpe inequality.
Now
$$(2\chi +3\tau )(M)=\frac{1}{4\pi ^2}_M\left(\frac{s_g^2}{24}+2|W_+|_g^2\frac{|\stackrel{}{r}|^2}{2}\right)𝑑\mu _g$$
for any metric on $`g`$ on $`M`$. If $`g`$ is an Einstein metric, the trace-free part $`\stackrel{}{r}`$ of the Ricci curvature vanishes, and we then have
$`(2\chi +3\tau )(X)k`$ $`=`$ $`(2\chi +3\tau )(M)`$
$`=`$ $`{\displaystyle \frac{1}{4\pi ^2}}{\displaystyle _M}\left({\displaystyle \frac{s_g^2}{24}}+2|W_+|_g^2\right)𝑑\mu _g`$
$`>`$ $`{\displaystyle \frac{2}{3}}(2\chi +3\tau )(X)`$
by Proposition 3.2. If $`M`$ carries an Einstein metric, it therefore follows that
$$\frac{1}{3}(2\chi +3\tau )(X)>k.$$
The claim thus follows by contraposition.
Example Let $`X_4`$ be the intersection of two cubic hypersurfaces in general position. Since the canonical class on $`X`$ is exactly the hyperplane class, $`c_1^2(X)=1^233=9`$. Theorem 3.3 therefore tells us that if we blow up $`X`$ at $`3`$ points, the resulting $`4`$-manifold
$$M=X\mathrm{\#}3\overline{}_2$$
does not admit Einstein metrics.
But now consider the Horikawa surface $`N`$ obtained as a ramified double cover of the blown-up projective plane $`_2\mathrm{\#}\overline{}_2`$ branched over the (smooth) proper transform $`\widehat{C}`$ of the singular curve $`C`$ given by
$$x^{10}+y^{10}+z^6(x^4+y^4)=0$$
in the complex projective plane, where the singular point $`[0:0:1]`$ of $`C`$ is the point at which we blow up $`_2`$. By the Freedman classification of $`4`$-manifolds , both of these complex surfaces are homeomorphic to
$$11_2\mathrm{\#}53\overline{}_2.$$
However, $`N`$ has $`c_1<0`$, and so admits a Kähler-Einstein metric by the Aubin/Yau theorem . Thus, although $`M`$ and $`N`$ are homeomorphic, one admits Einstein metrics, while the other doesn’t. $`\mathrm{}`$
Example Let $`X_3`$ be a hypersurface of degree $`6`$. Since the canonical class on $`X`$ is twice the hyperplane class, $`c_1^2(X)=2^26=24`$. Theorem 3.3 therefore tells us that if we blow up $`X`$ at $`8`$ points, the resulting $`4`$-manifold
$$M=X\mathrm{\#}8\overline{}_2$$
does not admit Einstein metrics.
However, the Freedman classification can be used to show that $`M`$ is homeomorphic to the Horikawa surface $`N`$ obtained as a ramified double cover of $`_1\times _1`$ branched at a generic curve of bidegree $`(6,12)`$; indeed, both of these complex surfaces are homeomorphic to
$$21_2\mathrm{\#}93\overline{}_2.$$
However, this $`N`$ also admits a Kähler-Einstein metric, even though the existence of Einstein metric is obstructed on $`M`$. $`\mathrm{}`$
Example Let $`X_3`$ be a hypersurface of degree $`10`$. Since the canonical class on $`X`$ is six times the hyperplane class, $`c_1^2(X)=6^210=360`$. Theorem 3.3 therefore tells us that if we blow up $`X`$ at $`120`$ or more points, the resulting $`4`$-manifold does not admit Einstein metrics. In particular, this assertion applies to
$$M=X\mathrm{\#}144\overline{}_2.$$
Now let $`N`$ be obtained from $`_1\times _1`$ as a ramified double cover branched at a generic curve of bidegree $`(8,58)`$. Both $`M`$ and $`N`$ are then simply connected, and have $`c_1^2=216`$ and $`p_g=84`$; and both are therefore homeomorphic to
$$129_2\mathrm{\#}633\overline{}_2.$$
But again, $`N`$ has $`c_1<0`$, and so admits a Kähler-Einstein metric, even though $`M`$ does not admit an Einstein metric of any kind whatsoever.
In most respects, this example is much like the previous examples. However, this choice of $`N`$ is not a Horikawa surface, but instead sits well away from the Noether line of complex-surface geography. $`\mathrm{}`$
Infinitely many such examples can be constructed using the above techniques, and the interested reader might wish to explore their geography.
It should be noted that Theorem 3.3 is the direct descendant of an analogous result in , where scalar curvature estimates alone were used to obtain an obstruction when $`k\frac{2}{3}c_1^2(X)`$. It was later pointed out by Kotschick that this suffices to imply the existence of homeomorphic pairs consisting of an Einstein manifold and a $`4`$-manifold which does not admit Einstein metrics. An intermediate step between and Theorem 3.3 may be found in , where cruder Seiberg-Witten estimates of Weyl curvature were used to obtain an obstruction for $`k\frac{25}{57}c_1^2(X)`$.
## 4 How Sharp are the Estimates?
The estimates we have described in §2 are optimal in the sense that equality is achieved for Kähler metrics of constant negative scalar curvature, with the standard orientation and spin<sup>c</sup> structure. In this section, we will attempt to probe the limits of these estimates by considering metrics of precisely this type, but with non-standard choices of orientation and spin<sup>c</sup> structure.
One interesting class of $`4`$-manifolds which admit constant-scalar-curvature Kähler metrics are the complex surfaces with ample canonical line bundle. In terms of complex-surface classification , these are precisely those minimal surfaces of general type which do not contain $`_1`$’s of self-intersection $`2`$. The ampleness of the canonical line bundle is often written as $`c_1<0`$, meaning that $`c_1`$ is a Kähler class. A celebrated result of Aubin/Yau guarantees that there is a unique Kähler-Einstein metric on $`M`$, compatible with the given complex structure, and with Kähler class $`[\omega ]=c_1=H^{1,1}(M,)`$. The scalar curvature of such a metric is, of course, a negative constant; indeed, $`s=dim_{}M=4`$.
Now if $`M`$ is a compact complex manifold without holomorphic vector fields, the set of Kähler classes which are representable by metrics of constant scalar curvature is open in $`H^{1,1}(M,)`$. On the other hand, a manifold with $`c_1<0`$ never carries a non-zero holomorphic vector field, so it follows that a complex surface with ample canonical line bundle will carry lots of constant-scalar-curvature Kähler metrics which are non-Einstein if $`b_{}=h^{1,1}1`$ is non-zero. However, one might actually hope to find such metrics even in those Kähler classes which are far from the anti-canonical class. This expectation may be codified as follows:
###### Conjecture 4.1
Let $`M`$ be any compact complex surface with $`c_1<0`$. Then every Kähler class $`[\omega ]H^{1,1}(M,)`$ contains a unique Kähler metric of constant scalar curvature.
The uniqueness clause was recently proved by X.-X. Chen , using ideas due to Donaldson and Semmes. A direct continuity-method attack on conjecture has also been explored by S.-R. Simanca.
Let us now narrow our discussion to a very special class of complex surfaces.
###### Definition 4.2
A Kodaira fibration is a holomorphic submersion $`\varpi :MB`$ from a compact complex surface to a compact complex curve, such that the base $`B`$ and fiber $`F_z=\varpi ^1(z)`$ both have genus $`2`$. If $`M`$ admits such a fibration $`\varpi `$, we will say that is a Kodaira-fibered surface.
The underlying $`4`$-manifold $`M`$ of a Kodaira-fibered surface is a fiber bundle over $`B`$, with fiber $`F`$. We thus have a long exact sequence
$$\mathrm{}\pi _k(F)\pi _k(M)\pi _k(B)\pi _{k1}(F)\mathrm{}$$
of homotopy groups, so that $`M`$ is a $`K(\pi ,1)`$. Thus, any $`2`$-sphere in $`M`$ is homologically trivial, and so has self-intersection $`0`$; in particular, the complex surface $`M`$ cannot contain any $`_1`$’s of self-intersection $`1`$ or $`2`$. On the other hand, $`M`$ is of general type, so the above implies that $`c_1(M)<0`$. Kodaira-fibered surfaces thus provide us with an interesting testing-ground for Conjecture 4.1.
Now the product $`B\times F`$ of two complex curves of genus $`2`$ is certainly Kodaira fibered, but such a product also admits orientation-reversing diffeomorphisms, and so has signature $`\tau =0`$. However, as was first observed by Kodaira , one can construct examples with $`\tau >0`$ by taking branched covers of products; cf. . For example, let $`B`$ be a curve of genus $`3`$ with a holomorphic involution $`\iota :BB`$ without fixed points; one may visualize such an involution as a $`180^{}`$ rotation of a $`3`$-holed doughnut about an axis which passes though the middle hole, without meeting the doughnut. Let $`f:CB`$ be the unique $`64`$-fold unbranched cover with $`f_{}[\pi _1(C)]=\mathrm{ker}[\pi _1(B)H_1(B,_2)]`$; thus $`C`$ is a complex curve of genus $`129`$. Let $`\mathrm{\Sigma }C\times B`$ be the union of the graphs of $`f`$ and $`\iota f`$. Then the homology class of $`\mathrm{\Sigma }`$ is divisible by $`2`$. We may therefore construct a ramified double cover $`MB\times C`$ branched over $`\mathrm{\Sigma }`$. The projection $`MB`$ is then a Kodaira fibration, with fiber $`F`$ of genus $`321`$. The projection $`MC`$ is also a Kodaira fibration, with fiber of genus $`6`$. The signature of this example is $`\tau (M)=256`$, and so coincidentally equals one-tenth of its Euler characteristic $`\chi (M)=2560`$.
Now, more generally, let $`M`$ be any Kodaira-fibered surface with $`\tau >0`$, and let $`\varpi :MB`$ be a Kodaira fibration. Let $`p`$ denote the the genus of $`B`$, and let $`q`$ denote the genus of a fiber $`F`$ of $`\varpi `$. Indulging in a standard notational abuse, let us also use $`F`$ to denote the Poincaré dual of the homology class of the fiber. Since $`F`$ can be represented in de Rham cohomology by the pull-back of an area form on $`B`$, this $`(1,1)`$-class is positive semi-definite. On the other hand, $`c_1`$ is a Kähler class on $`M`$, and so it follows that
$$[\omega _ϵ]=2(p1)Fϵc_1$$
is a Kähler class on $`M`$ for any $`ϵ>0`$. If Conjecture 4.1 is true, there must therefore exist a Kähler metric $`g_ϵ`$ on $`M`$ of constant scalar curvature with Kähler class $`[\omega _ϵ]`$. Let us explore the global geometric invariants of this putative metric.
The metric in question, being Kähler, would have total scalar curvature
$$s_{g_ϵ}𝑑\mu _{g_ϵ}=4\pi c_1[\omega _ϵ]=4\pi (\chi +ϵc_1^2)(M)$$
and total volume
$$𝑑\mu _{g_ϵ}=\frac{[\omega _ϵ]^2}{2}=\frac{ϵ}{2}(2\chi +ϵc_1^2)(M).$$
The assumption that $`s_{g_ϵ}=\text{const}`$ would thus imply that
$`s^2={\displaystyle s_{g_ϵ}^2𝑑\mu _{g_ϵ}}`$ $`=`$ $`{\displaystyle \frac{32\pi ^2}{ϵ}}{\displaystyle \frac{(\chi +ϵc_1^2)^2}{2\chi +ϵc_1^2}}`$
$`=`$ $`16\pi ^2{\displaystyle \frac{\chi }{ϵ}}\left[1+(3+{\displaystyle \frac{9}{2}}\varrho )ϵ+O(ϵ^2)\right],`$
where we have set
$$\varrho =\frac{\tau (M)}{\chi (M)}.$$
Since a Kähler metric on a complex surface satisfies $`|W_+|^2s^2/24`$, we would also consequently have
$`{\displaystyle |W_+|_{g_ϵ}^2𝑑\mu _{g_ϵ}}`$ $`=`$ $`{\displaystyle \frac{1}{24}}{\displaystyle s_{g_ϵ}^2𝑑\mu _{g_ϵ}}`$
$`=`$ $`{\displaystyle \frac{2}{3}}\pi ^2{\displaystyle \frac{\chi }{ϵ}}\left[1+(3+{\displaystyle \frac{9}{2}}\varrho )ϵ+O(ϵ^2)\right].`$
It would thus follow that
$`W_{}^2={\displaystyle |W_{}|_{g_ϵ}^2𝑑\mu _{g_ϵ}}`$ $`=`$ $`12\pi ^2\tau (M)+{\displaystyle |W_+|_{g_ϵ}^2𝑑\mu _{g_ϵ}}`$
$`=`$ $`{\displaystyle \frac{2}{3}}\pi ^2{\displaystyle \frac{\chi }{ϵ}}\left[1+(3{\displaystyle \frac{27}{2}}\varrho )ϵ+O(ϵ^2)\right].`$
On the other hand, there are symplectic forms on $`M`$ which are compatible with the non-standard orientation of $`M`$; for example, the cohomology class $`F+ϵc_1`$ is represented by such forms if $`ϵ`$ is sufficiently small. A celebrated theorem of Taubes therefore tells us that the reverse-oriented version $`\overline{M}`$ of $`M`$ has a non-trivial Seiberg-Witten invariant . The relevant spin<sup>c</sup> structure on $`\overline{M}`$ is of almost-complex type, and its first Chern class, which we will denote by $`\overline{c}_1`$, is given by
$$\overline{c}_1=c_1+4(p1)F.$$
Of course, the conjugate almost-complex structure, with first Chern class $`\overline{c}_1`$,is also a monopole class of $`\overline{M}`$, and $`\overline{M}`$ will have yet other monopole classes if, for example, $`M`$ admits more than one Kodaira fibration and $`\tau (M)0`$.
Now recall that (19) asserts that
$$(1\delta )\frac{s}{\sqrt{24}}+\delta W_+\frac{2\pi }{\sqrt{3}}|c_1^+|$$
for all $`\delta [0,\frac{1}{3}]`$. One would like to know whether this inequality might also hold, quite generally, for some value of $`\delta >\frac{1}{3}`$. In order to find out, we apply this inequality to $`\overline{M}`$ with the above monopole class. Rewriting the inequality with respect to the complex orientation of $`M`$, we then get
$$(1\delta )\frac{s}{\sqrt{24}}+\delta W_{}\frac{2\pi }{\sqrt{3}}|\overline{c}_{1}^{}{}_{}{}^{}|,$$
(20)
and it is this inequality we shall now use to probe the limits of the theory.
Relative to any Kähler metric with Kähler class $`[\omega _ϵ]`$, one has
$`\overline{c}_1^+`$ $`=`$ $`{\displaystyle \frac{\overline{c}_1[\omega _ϵ]}{[\omega _ϵ]^2}}[\omega _ϵ]`$
$`=`$ $`{\displaystyle \frac{[c_1+4(p1)F][2(p1)Fϵc_1]}{[\omega _ϵ]^2}}[\omega _ϵ]`$
$`=`$ $`{\displaystyle \frac{(\chi +3ϵ\tau )}{[\omega _ϵ]^2}}[\omega _ϵ],`$
so that
$`|\overline{c}_1^+|^2`$ $`=`$ $`{\displaystyle \frac{(\chi +3ϵ\tau )^2}{[\omega _ϵ]^2}}`$
$`=`$ $`{\displaystyle \frac{1}{ϵ}}{\displaystyle \frac{(\chi +3ϵ\tau )^2}{2\chi +ϵc_1^2}}`$
$`=`$ $`{\displaystyle \frac{\chi }{2ϵ}}\left[1(1{\displaystyle \frac{9}{2}}\varrho )ϵ+O(ϵ^2)\right].`$
Now since $`\overline{c}_1`$ is an almost-complex structure on $`\overline{M}`$, we have
$$|\overline{c}_{1}^{}{}_{}{}^{}|^2|\overline{c}_1^+|^2=2\chi 3\tau ,$$
so
$`|\overline{c}_{1}^{}{}_{}{}^{}|^2`$ $`=`$ $`(2\chi 3\tau )+{\displaystyle \frac{\chi }{2ϵ}}\left[1(1{\displaystyle \frac{9}{2}}\varrho )ϵ+O(ϵ^2)\right]`$
$`=`$ $`{\displaystyle \frac{\chi }{2}}(46\varrho )+{\displaystyle \frac{\chi }{2ϵ}}\left[1(1{\displaystyle \frac{9}{2}}\varrho )ϵ+O(ϵ^2)\right]`$
$`=`$ $`{\displaystyle \frac{\chi }{2ϵ}}\left[1+(3{\displaystyle \frac{3}{2}}\varrho )ϵ+O(ϵ^2)\right].`$
After dividing by $`\pi \sqrt{2\chi /3ϵ}`$, the inequality (20) would thus read
$$(1\delta )\sqrt{1+(3+\frac{9}{2}\varrho )ϵ+O(ϵ^2)}+\delta \sqrt{1+(3\frac{27}{2}\varrho )ϵ+O(ϵ^2)}\sqrt{1+(3\frac{3}{2}\varrho )ϵ+O(ϵ^2)}.$$
Dropping the terms of order $`ϵ^2`$, we would thus have
$$(1\delta )\left[1+(\frac{3}{2}+\frac{9}{4}\varrho )ϵ\right]+\delta \left[1+(\frac{3}{2}\frac{27}{4}\varrho )ϵ\right]1+(\frac{3}{2}\frac{3}{4}\varrho )ϵ,$$
so that, upon collecting terms, we would obtain
$$3\varrho ϵ9\varrho ϵ\delta .$$
Taking $`\varrho =\tau /\chi `$ to be positive, and noting that $`ϵ`$ is positive by construction, this shows that Conjecture 4.1 would imply that
$$\frac{1}{3}\delta ,$$
or in other words that (15) is optimal. We have thus proved the following result:
###### Theorem 4.3
Either
* the estimate (15) is optimal; or else
* Conjecture 4.1 is false.
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# A Determination of the Hubble Constant From Cepheid Distances and a Model of the Local Peculiar Velocity Field
## 1 Introduction
A long-standing goal of observational cosmology is the measurement of the expansion rate of the universe, parameterized by the Hubble constant, $`H_0.`$ Knowledge of $`H_0`$ enables us to assign galaxies absolute distances $`d`$ from their redshifts $`cz`$,using $`d=cz/H_0,`$ for $`z1.`$ More fundamentally, the Hubble constant measures the time since the Big Bang, or the “expansion age” of the universe: $`t_0=f(\mathrm{\Omega }_\mathrm{m},\mathrm{\Omega }_\mathrm{\Lambda })H_0^1,`$ where $`\mathrm{\Omega }_\mathrm{m}`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ are the density parameters for mass and the cosmological constant (or “dark energy”), respectively. The function $`f(\mathrm{\Omega }_\mathrm{m},\mathrm{\Omega }_\mathrm{\Lambda })`$ has well known limiting values $`f=2/3`$ for $`\mathrm{\Omega }_\mathrm{m}=1`$ and $`f=1`$ for $`\mathrm{\Omega }_\mathrm{m}=0,`$ if $`\mathrm{\Omega }_\mathrm{\Lambda }=0.`$ In the flat ($`\mathrm{\Omega }_\mathrm{m}+\mathrm{\Omega }_\mathrm{\Lambda }=1`$) models now favored by CMB anisotropy measurements (e.g., Tegmark & Zaldarriaga 2000; Lange et al. 2000), $`f`$ is larger, at given $`\mathrm{\Omega }_\mathrm{m},`$ than in the $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ case. However, unless $`\mathrm{\Omega }_\mathrm{m}\stackrel{<}{}0.25,`$ which is disfavored by a variety of data (Primack 2000), $`f1`$, even in a flat universe. It follows that for currently acceptable values of the density parameters, the expansion age of the universe is $`\stackrel{<}{}H_0^1.`$
By convention, extragalactic distances are measured in Mpc, redshifts in $`\mathrm{km}\mathrm{s}^1,`$ and $`H_0`$ in $`\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1.`$ From the mid-1960s through the mid-1980s, two groups dominated the debate over $`H_0.`$ One, associated mainly with Sandage and collaborators, argued that $`H_0=50\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ with relatively small ($`10\%`$) uncertainty. A second, led by de Vaucouleurs, advocated $`H_0=100\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ with similarly small error. The corresponding values of the expansion timescale are $`H_0^1=9.8h^1`$ Gyr, where $`hH_0/100\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1.`$ Thus, the large Hubble constant favored by de Vaucouleurs leads to a “young” ($`t_0\stackrel{<}{}10`$ Gyr) universe, while the small $`H_0`$ favored by Sandage corresponds to an “old” ($`t_0\stackrel{>}{}15`$ Gyr) universe. In recent years, the debate has shifted, with many groups finding $`H_0`$ to be intermediate between the Sandage and de Vaucouleurs values. Especially important in this regard has been the work of the Hubble Space Telescope (HST) Key Project on the Extragalactic Distance Scale ($`H_0`$KP), which finds $`H_0=71\pm 6\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ (Mould et al. 2000). We discuss their methods and results further below.
An independent constraint on the age of the universe comes from the age of the oldest stars $`t_{}`$. This can be measured from the turnoff point in the Hertzsprung-Russell diagrams of old globular clusters. The best current estimates (Krauss 1999; see also Caretta et al. 1999) suggest that $`t_{}=12.8\pm 1.0`$ Gyr ($`1\sigma `$ error), and that $`10t_{}17`$ Gyr at 95% confidence. If one furthermore assumes that the globular clusters did not form until about $`\mathrm{\Delta }t1`$ Gyr after the Big Bang, then the age of the universe as indicated by the oldest stars is $`t_{}+\mathrm{\Delta }t14\pm 1`$ Gyr. With the above estimates, we thus require that $`t_0`$ be strictly larger than 10 Gyr, and prefer that it be $`\stackrel{>}{}13`$ Gyr, to ensure consistency of the Big Bang model with stellar ages.
From this perspective, the de Vaucouleurs value of $`H_0`$ yields far too small an expansion time, while the Sandage value produces one that is comfortably large. The current modern value ($`H_0`$KP) gives a marginally consistent $`t_0=13.3\pm 1.3`$ Gyr if we assume an $`\mathrm{\Omega }_\mathrm{m}=0.3,`$ $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ universe as preferred by a variety of present data. A Hubble constant only 20% larger, however, would give an expansion age of 10.6 Gyr, and thus conflict with the best estimates of the globular cluster ages.
The determination of the Hubble constant clearly remains a crucial part of the cosmological puzzle. Recent efforts by the $`H_0`$KP and other groups have greatly reduced the allowed range for $`H_0,`$ but have not unequivocally demonstrated consistency between the timescales of Big Bang cosmology and stellar evolution. The main purpose of this paper is to underscore the importance of ongoing work on the problem, by presenting an alternative approach, using existing data, to measuring $`H_0.`$ The outline of this paper is as follows: In §2 we discuss the effects of peculiar velocities on the determination of $`H_0,`$ and strategies for overcoming these effects. In §3 we present a derivation of the Cepheid PL relation using the Optical Gravitational Lensing Experiment (OGLE; Udalski et al. 1999a,b) database of LMC Cepheids, and in §4 we apply the new PL relation to the $`H_0`$KP Cepheid database to obtain distances for the $`H_0`$KP galaxies. In §5 we constrain the local peculiar velocity field by applying the maximum likelihood VELMOD method to a sample of galaxies with accurate relative distances from surface brightness fluctuations. In §6 we apply the resulting velocity models to the $`H_0`$KP Cepheid galaxies, and thereby obtain a value of $`H_0.`$ In §7 we further discuss and summarize our results.
## 2 Peculiar velocities and the strategies for measuring $`H_0`$
Measuring $`H_0`$ requires addressing the effects of peculiar velocities—those deviations from the underlying Hubble flow. The observed redshift $`cz`$ of a galaxy at a distance $`d`$ is given by $`cz=H_0d+u,`$ where $`u`$ is the radial component of its peculiar velocity. In the (unfortunately) hypothetical case of pure Hubble flow ($`u0`$), a few good distance and redshift measurements of very nearby galaxies would cleanly yield $`H_0.`$
Over the last two decades, however, it has become clear that galactic peculiar velocities are both substantial ($``$ several hundred $`\mathrm{km}\mathrm{s}^1`$) and systematic, i.e., coherent over volumes of diameter $`10`$$`20`$ Mpc (see Willick 2000 for a recent review). Neglecting peculiar velocities can produce an error $`\delta H_0/H_0u/cz`$ for a single galaxy, while observing $`N`$ galaxies does not necessarily lead to a $`\sqrt{N}`$ reduction due to the coherence of the velocity field. $`H_0`$ measurement errors due to uncorrected peculiar velocities can approach $`30\%`$ at redshifts as large as $`1500\mathrm{km}\mathrm{s}^1,`$ roughly the distance to the Virgo cluster, even if moderately large ($`N10`$) samples are used.
At least two straightforward paradigms exist to handle the impact on $`H_0`$ from peculiar velocities: In the first approach (Method I ), one measures $`H_0`$ in the “far field” of the Hubble flow—peculiar velocities here are washed out in comparison with the large expansion velocities $`H_0d.`$ In the second approach (Method II ), one measures $`H_0`$ in the “near field” of the Hubble flow after using an accurate model of the local pecuilar velocity field $`u(𝐫)`$ to prune the observed redshifts $`cz`$ of their peculiar velocity contributions.
Method I requires no knowledge of peculiar velocities, provided they are fractionally small. At redshifts beyond $`5000\mathrm{km}\mathrm{s}^1,`$ for example, the fractional uncertainty in the Hubble constant, $`u/cz,`$ will generally be $`<10\%`$ for a single galaxy. Method I does, however, require data at distances where it is difficult or impossible to directly employ primary distance indicators such as Cepheids. Instead, secondary distance indicators, such as the Tully-Fisher relation or the Surface Brightness Fluctuation method (SBF), must be used. Not only are secondary DI’s less intrinsically accurate than Cepheids, but they also lack a priori absolute calibrations. This motivates a first look at how sources of error in the calibration process and distance scale relatively impact Method I and Method II analyses.
Primarily, uncertainty in the absolute calibration of the Cepheid PL relation, (itself due to uncertainty in the distance to the LMC), floods the systematic uncertainty in $`H_0`$. As Gibson (1999) has emphasized, values as small as $`\mu _{\mathrm{LMC}}=18.2`$ mag and as large as $`\mu _{\mathrm{LMC}}=18.7`$ mag have appeared recently in the literature, though $`H_0`$KP adopts $`\mu _{\mathrm{LMC}}=18.50`$ mag. $`H_0`$ estimates could be revised upwards by as much as $`15\%`$ if the smallest LMC distances prove correct, or downward by as much as $`10\%`$ should the largest LMC distance hold.
Also, a recalibration of the Cepheid P-L relation is desired because the current P-L relation used by the $`H_0`$KP is based on only 32 Cepheids in the LMC (Ferrarese et al. 2000b) while the OGLE has found many more. The $`H_0`$KP is also in the process of redetermining Cepheid galaxy distances using a calibration based on the OGLE data (W. Freedman, private communication; Madore & Freedman 2000, in preparation). Our own work with the OGLE data (see §4) indicates that the actual Cepheid distances are systematically $`5\%`$ shorter than the distances originally used by the HOKP.
These Cepheid uncertainties will propagate through any Method I or II approach that uses the Cepheid distance scale. The extent of this propagation motivates our choice of method. Method II approaches cut the propagation chain and use these Cepheid distances directly. Method I approaches use this Cepheid distance scale to calibrate a secondary distance scale, as mentioned above, which then provides the distances used to measure $`H_0`$. We focus on those Method I approaches that use Surface Brightness Fluctuations (SBF; Tonry & Schneider 1988) as the secondary distance indicator to highlight the difficulties in the calibration process and offer a first motivation toward Method II approaches. Ferrarese et al. (2000a) give a more exhaustive review of the calibration process in other secondary distance indicators.
Table 1 shows the tabulated SBF calibration data from Tonry et al. (2000; hereafter TBAD00). Column 1 lists the 6 galaxy names for which both color-adjusted SBF apparent bulge magnitudes ($`\overline{m_I}`$; Column 2) and Cepheid-based distance moduli ($`\mu _{ceph}`$; Column 3) are available.
Table 1 reveals a spread in $`\overline{M_I}`$ of $`.8`$ mag, with half of the measured values within the estimated systematic error of $`.16`$ mag of the $`H_0`$KP SBF zero point, $`1.79`$ mag (Ferrarese et al. 2000a), and the Tonry et al. (TBAD00) zero-point of $`1.74`$ mag. Both values agree with the theoretically predicted value, $`1.81`$ mag (Worthey 1993); the $`H_0`$KP value leads to a distance scale that is in good agreement with other secondary distance scales (Ferrarese et al. 2000a).
The SBF and Cepheid Virgo distances indicate that the preferred calibration of Tonry et al. (TBAD00), as well as the similar $`H_0`$KP calibration (Ferrarese et al. 2000a), is inconsistent with our revised Cepheid distance scale. As noted in (§4) the mean Cepheid distance to the five Virgo galaxies considered in this paper is 14.7 Mpc. If only the three galaxies within the canonical 6 degrees of the Virgo center are considered, this rises to 15 Mpc. By comparison, the mean SBF distance, using the Tonry et al. calibration above, of 27 Virgo galaxies in the SBF sample is 16.5 Mpc. If the SBF zero point, formerly in agreement with theory and other distance indicators, were shifted to bring the SBF Virgo galaxies to the 15 Mpc value suggested by the Cepheids, the value for $`H_0`$ derived from the analyses of TBAD00 and Blakeslee et al. (1999; see §5) would rise by $`7\%`$, while the $`H_0`$KP value would rise by $`10\%`$, close to the upper edge of the (Cepheid + SBF) systematic error envelope. This new result would still depend critically on the 6 calibrating galaxies of Table 1. As illustrated, the inevitable errors in the Cepheid calibration propagate directly into the secondary distance calibration, and the derived $`H_0`$.
The main advantage of Method II over Method I is that no such intermediate calibration step is required. One uses only the primary distances from Cepheid variables, which are more reliable within $`2000\mathrm{km}\mathrm{s}^1`$, but still subject to uncertainties. The disadvantage of Method II is that it requires an accurate model of the local peculiar velocity field. Turner, Cen & Ostriker (1992) performed theoretical analyses of the impact a lack of an accurate velocity field would have on local measures of $`H_0`$; constructing a satisfactory model has proved difficult but possible, as we discuss below. However, any systematic errors in the model result in errors in the corrected redshifts, and thus in the derived Hubble constant. The gains of removing the systematics from secondary calibration are countered by the velocity field systematics.
In the past decade Method I has taken precedence in Hubble constant determinations. The $`H_0`$KP, in particular, has adopted the philosophy of Method I, focusing its efforts on determining Cepheid distances to galaxies that can serve as suitable calibrators for secondary distance indicators. As the foregoing discussion indicates, however, the two approaches emphasize different systematic effects, and should each be employed as mutual checks. We attempt in this paper to redress the imbalance by applying Method II. Our approach has been made possible by the advent of data sets that allow the local velocity field to be more accurately modeled than previously, as we discuss in detail in §5.
## 3 Recalibration of the Cepheid Period-Luminosity Relation
In §4 we will derive distances to the $`H_0`$KP sample. First, however, we perform the recalibration of the Cepheid PL relation discussed in §2, based on the OGLE data.
The OGLE project has monitored fields in the LMC nearly every clear night for the last three years. Its primary goal was to detect gravitational lensing events; however, in the course of this monitoring, the OGLE discovered more than 1300 Cepheid variable stars in the LMC.<sup>2</sup><sup>2</sup>2The OGLE catalog of Cepheid variables in the LMC is publicly available at http://astro.princeton.edu/ogle/ogle2/var\_stars/lmc/cep/catalog. Of these, a majority are fundamental mode pulsators of the sort observed by the $`H_0`$KP in more distant galaxies. We use a subsample of the OGLE Cepheids judged by Udalski et al. (1999a) to be fundamental mode pulsators, and that have periods greater than 2.5 days, to calibrate the I and V band PL relations. Cepheids with shorter periods may not be fundamental mode pulsators, and moreover are not present in the $`H_0`$KP data set which contains only luminous Cepheids. There are 729 such stars in the OGLE catalog. Of these, 331 also have B band data, and we use this subset to calibrate the B band PL relation. This latter calibration is for reference only, however, as the $`H_0`$KP data consists only of V and I band Cepheid data.
We correct the observed Cepheid mean magnitudes for total extinction, Galactic plus LMC, as determined by the reddenings for each field given by Udalski et al. (1999b). We adopt the ratios of extinction to reddening, $`R_XA_X/E(BV)`$ (where $`X=B,`$$`V,`$$`I`$) given by Schlegel, Finkbeiner, & Davis (1998; SFD). These values are given in Table 2. We then fit a linear apparent magnitude versus log period relation in each of the three bandpasses:
$$m_X=a_X+b_X\mathrm{log}P.$$
(1)
The extinction corrected magnitudes are plotted versus $`\mathrm{log}P`$ in Figure 1. The solid lines show the best-fit PL relation for each bandpass. The parameters $`a_X`$ and $`b_X,`$ along with related information, are given in Table 2. Outliers were eliminated by iterating the fit several times and excluding objects which deviated by more than $`2.5\sigma .`$ The final fit parameters are those obtained after this iterative procedure converged. The number of objects used in the final fit is given in Table 2 as $`N_{fit},`$ as compared with the quantity $`N_{tot},`$ the total number of data points in the given bandpass meeting the initial cut on $`\mathrm{log}P.`$ The exclusion of outliers is justified because of the likelihood of contamination of the sample by first-overtone Cepheids (see Udalski et al. 1999a for further details).
### 3.1 Absolute Calibration of the Cepheid PL Relation
The apparent magnitude-log period relations given in Table 2 need to be converted into absolute PL relations in order to use them as distance indicators for other galaxies. To do this, we must adopt a distance modulus for the LMC, $`\mu _{\mathrm{LMC}}.`$ The LMC distance has been a contentious issue and remains, perhaps, the greatest source of systematic uncertainty in the extraglactic distance scale (Ferrarese et al. 2000b). Recent determinations of $`\mu _{\mathrm{LMC}}`$ (see Gibson 1999 for a review) have ranged from $`\mu _{\mathrm{LMC}}=18.20\pm 0.05`$ (Popowski & Gould 1998; Stanek, Zaritsky, & Harris 1998; Udalski et al. 1998a,b) to $`\mu _{\mathrm{LMC}}=18.55\pm 0.05`$ (Cioni et al. 2000; Hoyle, Shanks, & Tanvir 2000). We will follow the $`H_0`$KP and adopt $`\mu _{\mathrm{LMC}}=18.50`$. Although this may be somewhat higher than the mean of recent measurements, we believe it is the conservative choice for now. In any case, our ultimate Hubble constant scales in a simple way with $`\mu _{\mathrm{LMC}},`$ and can be adjusted accordingly should the LMC distance become better determined in the future.
Using $`\mu _{\mathrm{LMC}}18.50,`$ the absolute Cepheid PL relations follow directly from the parameters given in Table 2. We write our final PL relations in the form
$$M_X=A_X+b_X(\mathrm{log}P1).$$
(2)
We thus set the zero point at $`P=10`$ d, corresponding to a “typical” fundamental mode Cepheid. Comparison of Eqs. (1) and (2) shows that $`A_X=a_X+b_X18.5`$ for our adopted value of $`\mu _{\mathrm{LMC}}.`$ The PL slopes are of course the same. We thus obtain the final parameters for the Cepheid PL relations shown in columns (2) and (3) of Table 3. The slope errors are the same as those given in Table 2; we do not tabulate zero point errors because they are completely dominated by the $`0.2`$ mag systematic uncertainty in the LMC distance modulus.
It is useful to compare these PL parameters to those obtained from a completely separate sample: Cepheids from the Hipparcos database with parallax distances. Obviously, the latter are unaffected by the distance to the LMC, although they have other problems, such as incompleteness and potential systematic parallax errors. Lanoix, Paturel, & Garnier (1999b) have calibrated the PL relation for 174 Cepheids with Hipparcos parallaxes and have obtained the zero point and slope given in columns (3) and (4) of Table 3. As can be seen, there is excellent agreement to within the Hipparcos-based errors of the PL zero points. The V band slopes are in similarly excellent agreement. The I band slopes agree less well, but it is difficult to gauge the significance of the disagreement because Lanoix et al. (1999b) do not give slope errors. In any case, because $`\mathrm{log}P=1`$ is a typical extragalactic period, the slope difference will not translate into a large predicted absolute magnitude difference. Hence we can say with confidence that the OGLE-based PL relations from the LMC are in good agreement with the PL relations derived from Galactic Cepheids with Hipparcos parallaxes. This argues against a large ($`\stackrel{>}{}0.2`$ mag) error in our adopted LMC distance modulus.
## 4 Distances to Cepheid Galaxies
We apply the OGLE PL relations given in Table 2 to a sample of thirty-four galaxies for which V and I band Cepheid data are available from the $`H_0`$KP team and two other groups. It is supplemented by a small number of Cepheids with ground-based V and I band data. All the data were acquired from the electronic archive maintained by P. Lanoix at the URL http://www-obs.univ-lyon1.fr/$``$planoix/ECD; See Lanoix et al. (1999a) for more details.
Basic data for the thirty-four galaxies, listed in order of increasing heliocentric redshift, are given in Table 4. The names, in Column 1, are those preferred by Lanoix. Galactic longitude ($`\mathrm{}`$, Column 2) and latitude ($`b,`$ Column 3) and heliocentric redshift ($`cz_{},`$ Column 4) were all obtained from the NASA Extragalactic Database (NED; http://nedwww.ipac.caltech.edu). For reference we also list the Local Group (LG) frame and CMB frame redshifts, $`cz_{\mathrm{LG}}`$ and $`cz_{\mathrm{CMB}},`$ in Columns 5 and 6. Galactic reddenings $`E(BV)`$ determined by SFD (also obtained from NED) are listed in Column 7.<sup>3</sup><sup>3</sup>3The SFD reddenings represent the effects of Galactic dust only, and therefore are not suitable for correcting the Cepheid magnitudes for extinction, which is often dominated by extinction within the host galaxy. We do not use the SFD reddenings in any case, but present them here for completeness. Column 8 lists the number of Cepheids in each galaxy for which $`V`$ and $`I`$ band PL data are available.
Further references for each galaxy are given in column 9 of Table 4. Seven galaxies are denoted “LG,” indicating that they are members of the LG according to the compilation of Mateo (1998).<sup>4</sup><sup>4</sup>4For all seven we obtain Cepheid distances (see Table 5) smaller than 1.5 Mpc, and for four of them we obtain distances less than 0.8 Mpc. Our distances agree with those given by Mateo (1998) to within the errors. The $`V`$ and $`I`$ band data compiled by Lanoix for these objects are ground-based. The LG galaxies are included here for completeness, but are not used in the $`H_0`$ determination presented in §6. Galaxies in the LG have no leverage on $`H_0`$ because their Hubble velocities are negligible in comparison with random peculiar motions.
Of the 27 remaining objects, all of which are used in the $`H_0`$ analysis in §6, 26 have HST data. The one exception is NGC 300, which has ground-based data from Freedman et al. (1992; F92). PL data for this galaxy in the Lanoix database that are not from F92 are not used in our analysis.
For the twenty-six galaxies with HST data, eighteen were observed originally by the $`H_0`$KP team. These objects are denoted “$`H_0`$KP” in Table 4. The complete Cepheid database for these galaxies, as well as a detailed compendium of publications by the $`H_0`$KP, is given at the Key Project website.<sup>5</sup><sup>5</sup>5 The Key project web site can be found at www.ipac.caltech.edu/$`H_0`$KP. The Lanoix database includes this database as a subset; a crosscheck confirms that the periods and magnitudes in the two are identical. A further seven galaxies were originally observed by other HST investigators: six by the Sandage/Saha (Saha et al. 1999) group, and one by Tanvir et al. (1995), denoted “SS/KP” and “T/KP” respectively, in Table 4. The data for all seven were subsequently reanalyzed by the $`H_0`$KP team (Gibson et al. 2000). We use only the reanalyzed data from these seven galaxies in order to ensure uniformity with the rest of the $`H_0`$KP analyzed data that we use. Finally, one galaxy, NGC 4258, is denoted “Maoz/KP” in Table 4. It is not formally part of the $`H_0`$KP sample, but was observed with the HST by Maoz et al. (1999), in collaboration with members of the $`H_0`$KP team, in order to compare Cepheid distances with the highly accurate maser distance of Herrnstein et al. (1999; we discuss this special case further in §7). The Cepheid data for NGC 4258 are as yet unpublished, but were kindly provided to us in advance of publication by J. Newman. They are also available in Lanoix’s database.
### 4.1 Calculation of Distances
We determine the distance to each galaxy by minimizing a $`\chi ^2`$ statistic that measures deviations from the $`V`$ and $`I`$ band PL relations. The statistic is minimized with respect to variations in the galaxy distance modulus $`\mu `$ and the total reddening, Galactic plus internal, along the line of sight toward each Cepheid in the galaxy. We also assumed a single reddening value for all stars in the galaxy. This led to an identical value of the distance modulus but a much higher value of the PL scatter, indicating that the reddening is variable across the face of each galaxy.
Suppose we have PL data for $`i=1,\mathrm{},N_{\mathrm{Ceph}}`$ Cepheids in the galaxy in question, with one $`V`$ band and one $`I`$ band mean magnitude, and one period, for each star. Let $`m_{ij},`$ where $`j=V,I,`$ denote the magnitudes, and $`X_i=\mathrm{log}P_i,`$ where $`P_i`$ is the pulsation period in days, of the $`i`$th star. The appropriate $`\chi ^2`$ statistic is then
$$\chi ^2=\underset{i=1}{\overset{N_{\mathrm{Ceph}}}{}}\underset{j=V,I}{}\left[m_{ij}\left(A_j+b_j(X_j1)+\mu +R_jE(BV)_i\right)\right]^2/\sigma ^2.$$
(3)
Minimization of $`\chi ^2`$ yields the distance modulus $`\mu `$ and the $`N_{\mathrm{Ceph}}`$ reddenings $`E(BV)_i.`$ The PL parameters for the two bandpasses are hardwired to the values given in Table 3.
When we minimize this $`\chi ^2`$ statistic for each of the 34 galaxies listed in Table 4, we obtain the distance moduli, and the corresponding distances in Mpc, given in Table 5. We discuss the calculation of uncertainties below. We do not tabulate the reddenings for each star here, although this information can be made available electronically to interested readers. We do, however, list the mean reddenings, $`E(BV),`$ for each galaxy in Table 5. We note that, to within the errors, these reddenings are consistent with being larger than the SFD Galactic reddenings given in Table 4. When the reddening errors are sufficiently small for an estimate to be made—typically, when $`N_{\mathrm{Ceph}}\stackrel{>}{}30`$—we find that the mean internal reddening, $`E(BV)_{int}=E(BV)E(BV)_{SFD98}`$, is of order $`0.1`$ mag.
Once a galaxy distance modulus and the $`N_{\mathrm{Ceph}}`$ reddenings are determined, one may calculate $`V`$ and $`I`$ band absolute magnitudes for each star as follows:
$$M_{i,j}=m_{ij}R_jE(BV)_i\mu ,$$
(4)
where $`i=1,\mathrm{},N_{\mathrm{Ceph}}`$ and $`j=V,I.`$ The PL relation for the entire sample may then be exhibited as a plot of $`M_{i,j}`$ versus $`\mathrm{log}P_i`$ for all stars in all galaxies, as shown in Figure 2. A total of 1021 stars is plotted in the Figure. The OGLE PL relations given in Table 3 are plotted through the points as dotted lines. Note that the $`V`$ band PL relation has smaller scatter than the $`I`$ band relation. In fact, the scatter visible in the plots is smaller than the true scatter, because the fits have one degree of freedom, the reddening, for each star (plus one for each galaxy).
When these degrees of freedom are properly taken into account, the PL scatters for the 1021 Cepheids are
$$\begin{array}{cccc}\hfill \sigma _V& =& 0.111\hfill & \mathrm{mag}\hfill \\ \hfill \sigma _I& =& 0.165\hfill & \mathrm{mag}.\hfill \end{array}$$
(5)
These scatters are similar to those derived from the LMC fit (Table 2), although there it was the $`I`$ band scatter that was smaller. This most likely reflects the fact that in the LMC fit the reddenings were given a priori, and reddening errors increase the $`V`$ band PL residuals more than $`I`$ band residuals. Taken together, the LMC and 34-galaxy sample results suggest that the $`V`$ and $`I`$ band PL scatters are about the same in the absence of extinction, and that $`\sigma _V\sigma _I\sigma _{\mathrm{Ceph}}=0.15`$ mag. We assume that this approximation holds for the remainder of the paper. Our method for calculating random distance errors for the Cepheid galaxies is given in Appendix A.
## 5 Fitting Velocity Field Models
In order to use the Cepheid galaxies to estimate the Hubble constant, we must, as explained in § 2, use an accurate model of the local velocity field. We employ two quite different models in this paper:
1. An IRAS model, in which peculiar velocities are based on the distribution of galaxies in the nearby universe ($`cz\stackrel{<}{}10,000\mathrm{km}\mathrm{s}^1`$) as determined by the 1.2 Jy IRAS redshift survey (Strauss et al. 1992; Fisher et al. 1995). The velocity field is then a function of the parameter $`\beta \mathrm{\Omega }_\mathrm{m}^{0.6}/b,`$ where $`b`$ is the biasing factor for IRAS galaxies. Before applying the model the IRAS model to the Cepheids, then, we must first determine the appropriate value(s) of $`\beta `$ to use.
2. A phenomenological model, in which the local velocity field is dominated by a few simple components. We adopt a model of the form recently used by TBAD00, which includes a dipole, a quadrupole, and two “attractors,” one centered on the Virgo Cluster and one on the Great Attractor (GA). This model has a large number of free parameters, whose values must be determined before the model can be applied to the Cepheids.
Each of the two models has strong and weak points. The IRAS model is more realistic and is better motivated physically. All mass fluctuations in the local universe, not only prominent attractors, affect the velocity field, as they must if structure grows from gravitational instability. However, the IRAS model suffers from undercounting early-type galaxies in clusters, and thus may well underestimate the importance of massive concentrations such as Virgo and the GA. In contrast, the Tonry (TBAD00) model allows one to ascribe as much influence on the velocity field to Virgo and the GA as the data warrant. However, because it includes only these attractors, the Tonry model may attribute greater or lesser importance to these than they would have in reality to compensate for missing mass concentrations and voids.
The fact that each model is imperfect suggests that the prudent approach is to try both. In what follows, we first describe (§5.1) our method for fitting velocity models. We next describe (§5.2) the $`281`$ galaxy Surface Brightness Fluctuation (SBF) data set to which we fit each of the two models and comment on a key difference between our treatment of this data set and previous treatments. After that we present (§5.3) the constraints that the SBF data set allows us to place on the IRAS model (i.e., the allowed values of $`\beta `$ and several ancillary parameters to be described), and (§5.4) the best-fitting parameters of the phenomenological Tonry model.
### 5.1 The VELMOD Method
For both optimizing our peculiar velocity models with the SBF data, and for determining $`H_0`$ with the Cepheid data (in §6), we use the VELMOD maximum likelihood approach. The method was described in detail in two papers (Willick et al. 1997b, hereafter WSDK; Willick & Strauss 1998, hereafter WS), so we limit ourselves to a brief overview here. When we fit a peculiar velocity model to redshift-distance data, two sources of variance enter in: distance measurement error and small-scale velocity noise. The latter is the part of the the peculiar velocity field that cannot be predicted by any model. It results from close gravitational encounters with other galaxies in groups, rather than from the systematic pull of the large-scale gravity field. When velocity models are fit to relatively distant ($`\stackrel{>}{}30h^1\mathrm{Mpc}`$) galaxies, the distance errors ($`\stackrel{<}{}10\%`$ for SBF or Cepheid distances, $`20\%`$ for Tully-Fisher distances) dominate. However, at the much smaller distances of interest here ($`\stackrel{<}{}20h^1\mathrm{Mpc}`$), the velocity noise is equal to or larger than distance errors. Thus, it is essential that we calibrate it properly.
VELMOD was designed with precisely this aim. It is based on the explicit expression for the probability that a galaxy along a given line of sight has a measured distance $`d`$ and redshift $`cz:`$
$$P(\mathrm{ln}d,cz;𝐩)_0^{\mathrm{}}r^2n(r)P(\mathrm{ln}d|r)P(cz|r)𝑑r,$$
(6)
where $`r`$ is the true (as opposed to measured) distance along the line of sight, and $`𝐩`$ is a vector of parameters that determine the model peculiar velocity field. Here and in what follows, we assume that $`r`$ and $`d`$ are measured in units of $`\mathrm{km}\mathrm{s}^1,`$ i.e., we take $`H_01.`$ We will drop this assumption only in the last step of the analysis, when we apply VELMOD to the Cepheid sample to determine $`H_0.`$ In Eq. (6), $`n(r)`$ is the number density of galaxies at distance $`r`$ along the line of sight,
$$P(\mathrm{ln}d|r)=\frac{1}{\sqrt{2\pi }\mathrm{\Delta }}\mathrm{exp}\left\{\frac{\left[\mathrm{ln}(d/r)\right]^2}{2\mathrm{\Delta }^2}\right\},$$
(7)
where $`\mathrm{\Delta }=0.46\delta \mu `$ is the fractional distance error, $`\delta \mu `$ is the distance modulus error, and
$$P(cz|r;𝐩)=\frac{1}{\sqrt{2\pi }\sigma _v(r)}\mathrm{exp}\left\{\frac{\left(cz[r+u(r)]\right)^2}{2\sigma _v(r)^2}\right\},$$
(8)
where $`\sigma _v(r)`$ is the velocity noise and $`u(r)`$ the radial component of the predicted velocity field. Note that we allow $`\sigma _v`$ to vary with position; in practice we can also allow it (in the case of the IRAS model) to vary with local number density. Of course, $`u(r),`$ $`\sigma _v(r),`$ and $`n(r)`$ depend on the parameter vector $`𝐩.`$
Eq. (6), the joint probability distribution of distance and redshift, is not the optimal quantity on which to base likelihood-maximization because it is quite sensitive to the density as well as the velocity model. It is more suitable to use the conditional probability,
$$P(\mathrm{ln}d|cz;𝐩)=\frac{P(\mathrm{ln}d,cz)}{_0^{\mathrm{}}P(\mathrm{ln}d,cz)d(\mathrm{ln}d)}=\frac{_0^{\mathrm{}}r^2n(r)P(\mathrm{ln}d|r)P(cz|r)𝑑r}{_0^{\mathrm{}}r^2n(r)P(cz|r)𝑑r},$$
(9)
which is less sensitive to the density model because of the presence of $`n(r)`$ in both numerator and denominator. The conditional probability that the $`i=1,\mathrm{},N`$ sample galaxies have observed distances $`d_i`$ given that their redshifts are $`cz_i`$ is then
$$P(\mathrm{data};𝐩)=\underset{i=1}{\overset{N}{}}P(\mathrm{ln}d_i|cz_i;𝐩).$$
(10)
The essential step in VELMOD is maximizing the above sample probability with respect to the model parameter vector $`𝐩.`$ In practice, this is done by minimizing the statistic
$$=2\mathrm{ln}P(\mathrm{data};𝐩).$$
(11)
WSDK showed using simulated data sets that minimizing $`,`$ as defined by Eqs. (9) through (11), recovers unbiased values of velocity model parameters.
### 5.2 The SBF Sample
The best current data set to use for constraining the local peculiar velocity field is the SBF sample of Tonry and collaborators (Tonry et al. 1997, 2001). The full sample comprises $`300`$ early-type (mainly E and S0) galaxies out to $`4000\mathrm{km}\mathrm{s}^1,`$ although most are within $`3000\mathrm{km}\mathrm{s}^1.`$ For our fitting, we use a subset of the sample consisting of 281 galaxies that are not in the LG, have $`(VI)`$ colors $`>0.9,`$ and are consistent with our velocity models at the $`3\sigma `$ level. The SBF distances are accurate in most cases to $`\stackrel{<}{}10\%.`$ This accuracy is considerably better than what is available for Tully-Fisher samples; moreover, the SBF method yields a reliable distance error estimate for each galaxy, whereas for Tully-Fisher one generally has only a global scatter which is itself somewhat uncertain. Having a good distance error estimate is crucial if the velocity noise information is to be properly incorporated into the VELMOD procedure.
TBAD00 and Blakeslee et al. (1999) have already used the SBF sample to fit local velocity models. The former study fit the phenomenological model mentioned above, and the latter an IRAS model, to the SBF distances. Neither study, however, used the VELMOD method per se. Tonry et al. used a related approach, one that accounted for small-scale velocity dispersion but not for the role of the volume element (the $`r^2`$ term in §5.1) or of density variations (the $`n(r)`$ term in §5.1). Blakeslee et al. employed the same method that Davis, Nusser, & Willick (1996) used to fit a Tully-Fisher sample at larger distances; as this approach does not fully account for the effects of velocity noise, it is less suitable for the nearby flow field analysis that is of greatest importance here. In what follows, we borrow from both the TBAD00 and the Blakelee et al. (1999) studies, in that we employ similar models of the velocity field, but we use the VELMOD method in order to treat effects neglected in those papers.
Unlike TBAD00 and Blakeslee et al. (1999), we do not assume an absolute calibration of the SBF relation. That is, we use the SBF data only as indicators of distances in $`\mathrm{km}\mathrm{s}^1,`$ not in Mpc. Specifically, Tonry et al. and Blakeslee et al. took the SBF relation to be
$$\overline{M_I}=1.744.5(VI),$$
(12)
where $`\overline{M_I}`$ is the mean fluctuation absolute magnitude of a galaxy of a given $`(VI)`$ color. They then assign an absolute distance $`d_{abs}=10^{\left[0.2(\overline{m_I}\overline{M_I})5\right]}`$ Mpc to a galaxy with measured apparent fluctuation magnitude $`\overline{m_I}.`$ Consequently, when they fit their velocity models, one of the unknown free parameters is necessarily the Hubble constant<sup>6</sup><sup>6</sup>6Indeed, TBAD00 reported $`H_0=77\pm 4\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ and Blakelee et al. (1999) reported $`H_0=74\pm 4\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ (random errors) in addition to their constraints on the velocity field.. In contrast, we write the SBF relation
$$\overline{M_I}=A4.5(VI),$$
(13)
where $`A`$ is a free parameter in our maximum likelihood analysis, and then compute SBF distances in $`\mathrm{km}\mathrm{s}^1`$ according to $`d=10^{\left[0.2(\overline{m_I}\overline{M_I})\right]};`$ it is this value of $`d`$ that enters into calculation of the probability, Eq. (9). Because we never convert SBF distances to Mpc, our analysis of the velocity field using the SBF data does not yield a value for, and is completely unrelated to, the value of the Hubble constant. That is, no assumptions concerning the distance scale go into our velocity field analysis; it is only later, when we apply the VELMOD method to the Cepheid galaxies, that absolute distances enter our analysis. Therefore, it is only the Cepheid distances themselves that directly affect the value of the Hubble constant. (We emphasize, however, that we adopt the Tonry et al. SBF slope ($`4.5`$) as well as their reported distance errors in our analysis.)
It should be noted that our use of a relative rather than an absolute zero point, while conceptually important, is largely a technical distinction. The values of $`H_0`$ obtained by Tonry et al. and Blakeslee et al. are entirely dependent on their specific choice of SBF zero point in Eq. (12), and this value is not very well determined (see Tonry et al. 1997 for an in-depth discussion). A different choice of zero point would have led Tonry et al. and Blakelee et al. to find a different $`H_0,`$ but not a different velocity field. Our approach simply makes this explicit by removing absolute distances altogether from the problem of determining the velocity field.
### 5.3 The IRAS Model
There are any number of “IRAS models” one can apply because, aside from the obvious question of the value of $`\beta ,`$ there are the issues of smoothing scale, filtering, nonlinear corrections to the velocity-density relation, etc. The range of possibilities was discussed by WSDK and WS, to which readers are referred for a detailed discussion of the relative merits of each. Tests in those papers with Tully-Fisher data, and with the SBF data (Willick, Narayanan, Strauss, & Blakeslee, in preparation; hereafter WNSB), have demonstrated clearly that the smallest possible Gaussian smoothing scale for the IRAS data, $`300\mathrm{km}\mathrm{s}^1,`$ yields the most accurate velocity field (see Berlind, Narayanan, & Weinberg 2000 for a theoretical justification of this fact). A second outcome of such tests is that the best fits to both the Tully-Fisher and SBF data are obtained when the linear theory velocity-density relation,
$$𝐯(𝐫)=\frac{\beta }{4\pi }d^3𝐫^{}\frac{\delta _g(𝐫^{})(𝐫^{}𝐫)}{\left|𝐫^{}𝐫\right|^3},$$
(14)
is used, where $`\delta _g`$ is the density contrast of IRAS galaxies smoothed on a 300 km s<sup>-1</sup> scale. Proposed modifications to Eq. (14) to account for nonlinear dynamics, as well as hypothesized nonlinear biasing relations, have not yet been shown to improve the fit, although they typically increase the best value of $`\beta `$ by $`10\%.`$ The reasons for this are unclear at present, but will be discussed in depth by WNSB. For now, we use the linear theory, $`300\mathrm{km}\mathrm{s}^1`$ smoothed IRAS velocity field based on Eq. (14).<sup>7</sup><sup>7</sup>7In the language of WSDK and WS, we note that our adopted model also employs Wiener filtering and Method IV.
Figure 3 shows the main results of applying VELMOD to the SBF data set using the linear IRAS velocity field. The $``$ versus $`\beta `$ plot shows a strong minimum (likelihood maximum) near $`\beta =0.4;`$ the solid curve is a cubic fit to the likelihood points, and its minimum determines the maximum likelihood value of $`\beta `$ and its $`1\sigma `$ uncertainty: $`\beta =0.38\pm 0.06.`$ The $`\beta =0.3`$ likelihood lies within $`4`$ units of $``$ from the minimum and is thus acceptable at the $`2\sigma `$ level, while the $`\beta =0.2`$ and $`\beta =0.5`$ models are about $`3\sigma `$ away from the maximum likelihood value. Note that $`\beta >0.5`$ is ruled out with very high confidence. Indeed, we do not plot results for $`\beta >0.6`$ since the likelihood is so poor. These results suggest that we apply the $`\beta =0.2,0.3,0.4,`$ and $`0.5`$ IRAS models to the Cepheid galaxies when we determine $`H_0`$ below, giving the most weight to the $`\beta =0.3`$ and $`\beta =0.4`$ results.
As is standard for IRAS VELMOD (WSKD), we also allow for a random motion of the LG $`𝐰_{\mathrm{LG}}`$ with respect to the IRAS prediction, which is treated as a free parameter. As found by WSKD and WS, this velocity is generally small $`\stackrel{<}{}150\mathrm{km}\mathrm{s}^1`$ for $`\beta `$ near its optimum value. This result is confirmed here in the lower right panel of Figure 3. The upper right panel shows the variation of the final free parameter, the velocity noise $`\sigma _v,`$ with $`\beta .`$ It, too, generally minimizes near the maximum likelihood $`\beta ,`$ and this occurs here as well. However, note that $`\sigma _v`$ has a considerably larger value, $`185\mathrm{km}\mathrm{s}^1,`$ for the best fit model than was found in the Tully-Fisher VELMOD fits of WS and WSKD, where $`\sigma _v130`$$`150\mathrm{km}\mathrm{s}^1`$ was a more typical value. It is probable that this reflects a real difference between the small-scale velocity noise for late-type (Tully-Fisher) and early-type (SBF) galaxies; we will address this issue in WNSB. When we apply the IRAS (and Tonry) models to the Cepheid galaxies in §6, we will use values of $`\sigma _v`$ typical of the Tully-Fisher spirals, rather than the large $`\sigma _v`$ found here.
As discussed by WSDK and WS, a good method for visually assessing the quality of the IRAS velocity model is to compute smoothed velocity residuals for the sample (see, e.g., Eq. (9) of WS for details on how such residuals are calculated). Sky plots of such residuals are presented in Figure 4 for the $`\beta =0.4`$ model. The Gaussian smoothing scale employed rises from $`250\mathrm{km}\mathrm{s}^1`$ at $`cz\stackrel{<}{}500\mathrm{km}\mathrm{s}^1`$ to $`700\mathrm{km}\mathrm{s}^1`$ at $`3500\mathrm{km}\mathrm{s}^1`$; thus, the smoothed residuals will exhibit coherence over $`30^{}`$ scales. Coherence on larger scales than this signifies a failure of the model.
The residuals in Figure 4 are essentially incoherent on large scales, as indicated by the fact that both starred (outflowing) and open (inflowing) symbols are well mixed throughout the plots. Furthermore, the amplitude of these velocity residuals is almost everywhere $`\stackrel{<}{}100\mathrm{km}\mathrm{s}^1,`$ and is in many places virtually zero (indicated by the smallest points on the figure). These aspects of Figure 4 indicate that the IRAS model provides a good fit to the SBF data. The open squares in Figure 4 indicate the positions of the Cepheid galaxies to which the IRAS and Tonry models will be applied in §5. The squares are seen to lie in regions well sampled by the SBF data, so that the velocity models are well-constrained where the Cepheid galaxies are found.
There is one region, however, where a conspicuous failure of the IRAS model seems to occur: in the lower right quadrant of the $`cz1000\mathrm{km}\mathrm{s}^1`$ map. A group of six SBF galaxies appear to have a coherent flow at $`150\mathrm{km}\mathrm{s}^1`$ relative to the IRAS model there. However, this pattern does not continue into the next redshift interval, where there are more Cepheid galaxies and where the leverage on the $`H_0`$ measurement is greater. Consequently, we do not attempt to correct for this apparent failure of the IRAS model (the Tonry model shows the same discrepancy). However, this anomaly remains unexplained and deserves further attention.
### 5.4 The Tonry Model
Our phenomenological model has the same functional form as the TBAD00 model. However, we have implemented the model in a system of units in which distance is measured in km s<sup>-1</sup>, and used the VELMOD formalism, which differs in a number of ways (§5.1, Batra & Willick 2000 ) from that of TBAD00. To maximize independence from the IRAS results, we assume $`n(r)=`$constant in Eqs. (6) and (9) when we implement the Tonry model.
An important distinction between the Tonry model and the IRAS models is that the former is carried out entirely in the CMB reference frame, i.e., the redshifts used are $`cz_{\mathrm{CMB}},`$ whereas for the IRAS models $`cz_{\mathrm{LG}}`$ is used. The frame of reference is a highly nontrivial issue for local velocity field fits, as redshifts can differ by up to $`600\mathrm{km}\mathrm{s}^1`$ (Table 4, Courteau & van den Bergh 1999). It is thus important to demonstrate that comparable results can be obtained regardless of frame.<sup>8</sup><sup>8</sup>8One can also implement the Tonry model in the LG frame. We have done this, and the value of $`H_0`$ we obtain is unchanged.
The Tonry model assumes the local peculiar velocity field to be dominated by two spherically symmetric attractors, one centered on the Virgo Cluster and one on the Great Attractor. Each attractor is taken to have a mean interior overdensity given by
$$\delta (r)=\frac{\delta _0e^{r/R_{cut}}}{1\gamma /3}\left(\frac{r}{R_c}\right)^3\left[\left(1+(r/R_c)^3\right)^{1\gamma /3}1\right]$$
(15)
where $`r`$ is the distance from the center of the attractor. Note that the attractors have both a “cutoff” radius $`R_{cut}`$ and a “core” radius $`R_c.`$ Each attractor produces an infall velocity given by the Yahil’s (1985) expression,
$$v_{in}(r)=\frac{1}{3}\mathrm{\Omega }_\mathrm{m}^{0.6}r\delta (r)\left[1+\delta (r)\right]^{1/4}$$
(16)
The value of $`\mathrm{\Omega }_\mathrm{m}`$ is not well constrained by the fit, being highly covariant with $`\delta _0`$ (TBAD00), and it thus suffices to fix it at an arbitrary value, which we take to be $`0.2.`$
To roughly account for the effects of mass inhomogeneities other than the attractors, the model also includes velocity dipole and quadrupole fields; the latter is exponentially truncated and centered on the LG. These fields are added vectorially to the infall velocities produced by the attractors.
We fit the Tonry model to the SBF data set, allowing some, but not all, of the model parameters to vary. As discussed in detail by TBAD00, there is significant covariance among the parameters, so to simplify the fit we take the attractor core radii $`R_c`$ and power-law exponents $`\gamma `$ to have the TBAD00 values. The model also contains the velocity dispersions of the tracer galaxies as a function of position (i.e., the quantity $`\sigma _v(r)`$ in Eq. (8)). TBAD00 took the background dispersion to be $`187\mathrm{km}\mathrm{s}^1,`$ and then added this value in quadrature with an additional dispersion, $`\sigma _v^{core},`$ within a distance $`R_c`$ of three clusters: Virgo, the Great Attractor, and Fornax. For these dispersions, too, we adopted the TBAD00 values. (Note, however, that Fornax does not contribute to the overall velocity field.)
Table 6 presents our best fit values for the parameters we allowed to vary, as well as those held fixed at the TBAD00 values. Those parameters which denote positions in space are given in Supergalactic coordinates, in $`\mathrm{km}\mathrm{s}^1`$ units. For this table, the dipole components, in $`\mathrm{km}\mathrm{s}^1`$, are $`(70,230,40)`$; the quadrupole matrix is
$$\left(\begin{array}{ccc}1.79& 2.17& 6.62\\ 2.17& 10.5& 3.99\\ 6.62& 3.99& 8.71\end{array}\right).$$
#### 5.4.1 Comparison of Velocity Models
Figure 5 provides a comparison of the IRAS and Tonry models along the lines of sight toward four SBF galaxies that lie in regions that are also important for the Cepheid analysis of §6. We plot both our own fit of the Tonry model to the SBF data (solid line) and that of TBAD00 (dotted line). Peculiar velocity in the LG frame is plotted as a function of Hubble-flow distance. The IRAS velocity field generally exhibits more “features,” i.e., it is not as smooth as the Tonry models. This results from the fact that the IRAS peculiar velocity field is produced by all mass fluctuations, while the phenomenological model assumes only the existence of two attractors. Another important difference is in the different strength of the Virgo (NGC 4476) and Great (NGC 4616) Attractors. The IRAS velocity field exhibits relatively weak gradients ($`|u^{}(r)|\stackrel{<}{}0.5`$ near Virgo, $`|u^{}(r)|\stackrel{<}{}0.25`$ near the Great Attractor), while the Tonry model has large gradients ($`|u^{}(r)|\stackrel{>}{}1`$) in these regions. The greater influence of the attractors arises because, in the Tonry model, they must account for all features of the velocity field, some of which are in reality due to other mass fluctuations. (On the other hand, the mild gradients in the IRAS $`u(r)`$ near Virgo may be due in part to the undercounting of cluster galaxies by IRAS.) As we shall see §6, the Tonry model does not fit the Cepheid data as well as the IRAS model; as a result, we shall, in the end, adopt the value of $`H_0`$ derived from the IRAS model.
## 6 Application of the Velocity Models to the Cepheid Galaxy Sample
In this section we apply the velocity models of §5 to the Cepheid galaxy sample discussed in §4. Our application is limited to the twenty-seven galaxies listed in Tables 4 and 5 that are not LG members. As noted in §4, 26 of these 27 galaxies have HST Cepheid data that have been analyzed by the $`H_0`$KP team; the one galaxy for which ground-based data are used, NGC 300, was analyzed by $`H_0`$KP team members (Freedman et al. 1992) and is thus expected to be on the same system. The Cepheid sample used here thus constitutes a uniform data set.
We now apply the VELMOD method to the Cepheid data set, as we did in §5 for the SBF data; here we determine $`H_0.`$ There is now one key difference: the Cepheid galaxy distances $`d`$ are in Mpc, and correspondingly, Eq. (7) is rewritten
$$P(\mathrm{ln}d|r)=\frac{1}{\sqrt{2\pi }\mathrm{\Delta }}\mathrm{exp}\left\{\frac{\left[\mathrm{ln}(H_0d/r)\right]^2}{2\mathrm{\Delta }^2}\right\}.$$
(17)
In the exercise of §5 a vector of free parameters, $`𝐩,`$ on which the velocity model depended, was varied to maximize likelihood. Now, $`𝐩`$ is held fixed at the values determined in §5—i.e., the same velocity field is used—and the only parameter that is varied to maximize likelihood is $`H_0.`$ Aside from these differences, the VELMOD procedure applied to the Cepheid galaxies is identical to that applied to the SBF data set.
### 6.1 Results from the IRAS models
To apply the IRAS model we must adopt a functional form and value of $`\sigma _v(r),`$ the small-scale velocity noise as a function of position. Following WS, we write
$$\sigma _v(r)=\sigma _{v,0}+f_v\delta _g(r),$$
(18)
where $`\delta _g`$ is the IRAS galaxy overdensity at position $`r`$ along the line of sight. Thus, $`\sigma _{v,0}`$ is the velocity noise in a mean-density environment, and $`f_v`$ represents the rate of increase of velocity dispersion with density, an effect expected on theoretical grounds and verified in N-body simulations (Kepner, Summers, & Strauss 1997; Strauss, Ostriker, & Cen 1998). For the Cepheid sample we adopt $`\sigma _{v,0}=135\mathrm{km}\mathrm{s}^1`$ and $`f_v=30\mathrm{km}\mathrm{s}^1,`$ similar to the values WSDK and WS found for Tully-Fisher VELMOD. As noted earlier, this value of $`\sigma _v`$ is considerably smaller than the $`185\mathrm{km}\mathrm{s}^1`$ found for the SBF sample, but the Cepheid galaxies are late-type spirals and are more likely to resemble the Tully-Fisher galaxies in their dynamical properties.
There is one subset of galaxies within the Cepheid sample that are not well-described by this model of $`\sigma _v(r),`$ namely, Virgo cluster members. Virgo’s relatively high ($`650\mathrm{km}\mathrm{s}^1`$) velocity dispersion cannot be matched by the linear increase with density at the IRAS smoothing scale of $`300\mathrm{km}\mathrm{s}^1.`$ To account for this we “collapse” Virgo, following WS and WSDK—i.e., we set the redshift of each Virgo galaxy to its mean value of $`cz_{\mathrm{LG},\mathrm{Virg}}=1035\mathrm{km}\mathrm{s}^1`$ (Huchra 1985). To account for uncertainty in this value, we set the $`\sigma _v=30\mathrm{km}\mathrm{s}^1`$ for each collapsed Virgo galaxy. The Cepheid galaxies deemed likely Virgo members were NGC4536, NGC4321, NGC4496A, NGC4535, and NGC4548. Their mean Cepheid distance is 14.7 Mpc. The distance of each one is within $`1.5\sigma `$ of this mean value, each is within 10 of the Virgo core at $`\mathrm{}=283.8^{},`$ $`b=74.5^{},`$ and each has a redshift within $`700\mathrm{km}\mathrm{s}^1`$ of the Virgo mean. Thus, these are all strong candidates for Virgo membership, and collapsing them significantly improves the fit likelihood. However, as we now show, whether or not we collapse Virgo makes little difference to the derived value of $`H_0.`$
Table 7 presents the main results of applying the IRAS velocity models to the Cepheid sample. Column 1 lists the value of $`\beta `$ (the values of $`𝐰_{\mathrm{LG}}`$ used are always those derived from the SBF fit at that value of $`\beta ,`$ and $`\sigma _v(r)`$ is always given as discussed in the previous paragraph); Column 2 lists the maximum likelihood value of $`H_0`$ and its $`1\sigma `$ uncertainty; columns 3 and 4 give the values of the likelihood statistic $``$ and a $`\chi ^2`$ for the fit, which is discussed further below. Columns 5, 6, and 7 repeat the information of 3, 4, and 5 for the case that Virgo is not collapsed, i.e., when the true LG redshifts of the Virgo galaxies are used in the likelihood analysis.
The most striking feature of Table 7 is the robustness of $`H_0`$ with respect to changes in $`\beta .`$ Indeed, the results of the Table can be summarized by the statement $`H_0=85\pm 2.5\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ at 65% ($`1\sigma `$) confidence, irrespective of the value of $`\beta ,`$ provided it is in the range $`0.2`$$`0.5`$ that is allowed by the IRAS velocity model applied to the SBF sample.<sup>9</sup><sup>9</sup>9The application of VELMOD to Tully-Fisher samples by WSDK and WS preferred higher values of $`\beta ,`$ in the range $`0.4`$$`0.6`$ at 95% confidence. We consider the SBF result to be more reliable because of the greater precision of SBF distances. In any case, we note that $`\beta =0.4`$ is allowed by both the SBF and Tully-Fisher data, and thus constitutes the preferred value at present. When Virgo is not collapsed, the likelihood statistics and $`\chi ^2`$ values indicate a much worse fit. This occurs because the Virgo galaxy redshifts scatter so widely about the mean cluster value, and our simple model of velocity noise increase with density does not account for this (possibly because the huge central densities are smoothed out, and because IRAS undercounts cluster cores). It is important to note, however, that not collapsing Virgo affects the derived value of $`H_0`$ very little, increasing it by an amount about equal to the $`1\sigma `$ error estimate. Thus our treatment of Virgo is not crucial to the main conclusion of this paper.
We discuss below the calculation of the confidence intervals on $`H_0.`$ First, we describe the calculation of the $`\chi ^2`$ statistics listed in Table 7.
#### 6.1.1 $`\chi ^2`$ statistic for the velocity fits
The likelihood statistic $``$ does not by itself provide a measure of goodness of fit. WSDK and WS pointed out that it was not possible to construct a rigorous $`\chi ^2`$ statistic for their Tully-Fisher VELMOD fits because the distance and velocity errors were determined as part of likelihood maximization. This is not the case here, however, because (i) we have reliable distance errors for the Cepheid galaxies, and (ii) we have fixed the Cepheid velocity noise a priori (Eq. 18). Thus it makes sense to define and calculate a $`\chi ^2`$ statistic to test the goodness of fit of our velocity models to the Cepheid galaxy data.
We first define the expected distance in $`\mathrm{km}\mathrm{s}^1,`$ or Hubble flow distance, given the redshift and the velocity model,
$$E(r|cz)=\frac{_0^{\mathrm{}}r^3n(r)P(cz|r)𝑑r}{_0^{\mathrm{}}r^2n(r)P(cz|r)𝑑r},$$
(19)
along with the corresponding mean square distance
$$E(r^2|cz)=\frac{_0^{\mathrm{}}r^3n(r)P(cz|r)𝑑r}{_0^{\mathrm{}}r^2n(r)P(cz|r)𝑑r}$$
(20)
and distance error,
$$\delta r=\sqrt{E(r^2|cz)\left[E(r|cz)\right]^2}.$$
(21)
The Cepheid distance in $`\mathrm{km}\mathrm{s}^1`$ is $`H_0d,`$ where $`d`$ is the Cepheid distance in Mpc and $`H_0`$ is the best-fit Hubble constant for the model in question, and the corresponding error is $`H_0d\mathrm{\Delta }`$ (see §5.1). The $`\chi ^2`$ statistic measures the difference between $`H_0d`$ and $`E(r|cz)`$ in units of the overall error:
$$\chi ^2=\underset{i=1}{\overset{27}{}}\frac{\left(E(r|cz)H_0d\right)^2}{(\delta r)^2+(H_0d\mathrm{\Delta })^2}.$$
(22)
It is important to note that the Hubble flow contribution to the error, $`\delta r,`$ is not determined solely by the velocity noise $`\sigma _v.`$ Rather, it is the integrated effect of $`\sigma _v`$ along the line of sight, taking into account the shape of the large-scale velocity field $`u(r).`$ In particular, in regions where $`u^{}(r)<0`$ the effect of velocity noise is enhanced, creating comparatively large $`\delta r`$ (see Appendix A of WS for further discussion).
The fourth and seventh columns of Table 7 list the $`\chi ^2`$ values for the Virgo-collapsed and uncollapsed fits. There are 26 degrees of freedom for the fit—27 galaxies minus one free parameter—so that the expected $`\chi ^2`$ value is 26, with rms dispersion $`\sqrt{52}=7.2.`$ The Virgo-collapsed fits with $`\beta 0.3`$ thus have $`\chi ^2`$ values within $`2\sigma `$ of their expected values. These fits are therefore statistically acceptable, albeit with a larger $`\chi ^2`$ than desirable.
Two points are worth bearing in mind with regard to this statement, however. First, the relatively high $`\chi ^2`$ value is largely due to the influence of one galaxy, NGC 1326A, which deviates from the fit by $`3\sigma .`$ When this object is excluded the computed $`\chi ^2`$ is well within $`1\sigma `$ of the expectation.<sup>10</sup><sup>10</sup>10Moreover, NGC1326A has very little effect on the value of $`H_0;`$ for $`\beta =0.4`$ we find $`H_0=84.2\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ when NGC1326A is excluded. Second, the absolute value of $`\chi ^2`$ is strongly dependent on our adopted $`\sigma _v(r).`$ We chose $`\sigma _v(r)`$ similar to, though somewhat larger than, the value found by WS for a Tully-Fisher sample. If we take $`\sigma _v(r)=185\mathrm{km}\mathrm{s}^1,`$ the value indicated by the SBF fit (§5.3), we obtain $`\chi ^2=29.9`$ for the $`\beta =0.4`$ fit, fully compatible with expected value of $`26\pm 7.2.`$ (The best-fit Hubble constant rises to $`86.8\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ for this choice of $`\sigma _v(r),`$ a change of less than $`1\sigma .`$) We believe that the smaller $`\sigma _v(r)`$ is appropriate for the late-type Cepheids, but this example shows that it is difficult to assign unambiguous significance to our $`\chi ^2`$ statistic. However, the statement that the SBF-constrained IRAS models provide satisfactory fits to the Cepheid sample for $`\beta =0.3`$$`0.5`$ is a reasonable one in view of the above discussion.
#### 6.1.2 Hubble diagrams and confidence intervals
To assess the quality of the fits and the derived values of $`H_0,`$ we plot Hubble diagrams in the left hand panels of Figures 6 and 7. The abcissa is the quantity $`E(r|cz),`$ the expected value of the Hubble flow distance given the LG frame redshift and the IRAS velocity model, as defined by Eq. (19). The vertical error bars are the quantity $`\delta r`$ defined by Eq. (21), and the horizontal error bars are the Cepheid distance uncertainty $`d\mathrm{\Delta }.`$ (We have approximated all errors as being symmetric.) In each Hubble diagram we plot the best-fit value of $`H_0`$ as a solid line through the points, with dashed lines representing the values $`H_0=70\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ and $`H_0=100\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ for comparison.
The diagrams validate the values of $`H_0`$ given in Table 7. The solid lines represent a far better fit to the majority of the data points than do the dashed lines drawn for comparison. The cluster of points very near the best fit Hubble line at a distance of 15 Mpc are the five collapsed Virgo galaxies. Their error bars are much smaller because we assume a velocity uncertainty of only $`30\mathrm{km}\mathrm{s}^1`$ for these objects, associated with the uncertainty in the Virgo redshift. (The error bars are, however, larger than $`30\mathrm{km}\mathrm{s}^1`$ because of the interaction between $`\sigma _v`$ and $`u(r)`$ mentioned above.) It is evident from the diagrams that the fits are not perfect, with larger deviations being seen at the distance of Virgo. This may represent significant departures from the IRAS model in dense regions, a topic we further discuss in §7.1. As noted above, however, our uncertainty about the precise value of $`\sigma _v(r)`$ means we cannot state with assurance that these points do not fit the model.
The Hubble constant errors we have quoted are calculated by calculating the likelihood statistic $``$ for a range of values of $`H_0`$ near the best fit value. (We hold all velocity field parameters, as well as $`\sigma _v(r),`$ constant as we vary $`H_0.`$) As discussed in detail by WSDK, $``$ has the property that, when one fit parameter is varied, increases of $`\mathrm{\Delta }=\pm 1`$ with respect to the minimum yield the $`1\sigma `$ errors on the parameter, while increases of $`\mathrm{\Delta }=\pm 4`$ yield the $`2\sigma `$ errors. This is strictly true if the likelihood is Gaussian, or, equivalently, if the variation of $``$ is parabolic near its minimum. This is a good approximation in the present case. The right hand panels of Figures 6 and 7 plot $``$ versus $`H_0`$ for the $`\beta =0.3`$ and $`\beta =0.4`$ models respectively. The horizontal dashed lines show the $`1`$ and $`2\sigma `$ confidence intervals on $`H_0`$ as determined by the procedure described above. These curves are the origin of the 95% confidence interval on $`H_0`$ of $`5\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ quoted in the Abstract.
### 6.2 Results from the Tonry Model
We apply the phenomenological Tonry model to the Cepheid galaxies using the best-fit parameters arising from the SBF velocity fit. We again collapse Virgo in the Cepheid fit. We now fix the baseline velocity noise to $`\sigma _v(r)=150\mathrm{km}\mathrm{s}^1,`$ allowing it to rise to the SBF values in the cores of the model’s attractors. The baseline noise is a bit higher than was the case for the IRAS fit (Eq. 18). However, the density dependence of $`\sigma (r)`$ for IRAS was such that non-cluster galaxies have very nearly the same velocity noise for the IRAS and Tonry model fits. The value of $`H_0`$ exhibits slight sensitivity to the value of $`\sigma _v,`$ but the variation is small relative to the statistical errors.
Figure (8) shows the Hubble diagram resulting from fitting the Tonry model to the Cepheid data (left panel), along the likelihood versus $`H_0`$ curve (right panel). The maximum likelihood result find $`H_0=91.8\pm 1\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ at 65% ($`1\sigma `$) confidence. This is considerably larger, relative to the errors, than the result from the IRAS model. However, note that the $`\chi ^2`$ for the Tonry model fit, indicated on the Figure, is substantially larger than the value obtained for the IRAS fit.<sup>11</sup><sup>11</sup>11The minimum value of the likelihood statistic $``$ is also markedly larger than for the IRAS fit; however, the absolute likelihood statistics are not necessarily comparable for the two models owing to differences in implementation, such as the adoption of a constant density for the Tonry model. Indeed, it deviates by more than $`3\sigma `$ from the expected value of 26. We conclude that the Tonry model is not an acceptable fit to the Cepheid distances; a corollary is that the discrepancy between the derived values of $`H_0`$ is not meaningful.
Note that the collapsed Virgo galaxies in the left panel of Figure 8 have vertical error bars that are much larger than those for the IRAS fit. This reflects the very strong Virgo infall inherent in the best-fit Tonry model—a triple-valued zone is produced near Virgo, which causes a small $`\sigma _v`$ to translate in to a large $`\delta r.`$ The large $`\chi ^2`$ for the fit probably indicates that the Tonry model attributes too much strength to the Virgo attractor as it attempts to compensate for the missing mass concentrations and voids that the IRAS model contains.
Table 8 lists the data plotted in Figures 6, 7 and 8.
## 7 Discussion and Summary
We have argued in this paper that $`H_0=85\pm 5\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ at 95% confidence, considering random error only. This result, if correct, leads to an expansion timescale $`H_0^1=10.9`$$`12.2`$ Gyr, and thus an expansion age $`t_0=f(\mathrm{\Omega }_\mathrm{m},\mathrm{\Omega }_\mathrm{\Lambda })H_0^1`$ that is shorter still (see the discussion in §1), unless $`\mathrm{\Omega }_\mathrm{m}=1\mathrm{\Omega }_\mathrm{\Lambda }\stackrel{<}{}0.25.`$ For example, for $`H_0=85\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ and an $`\mathrm{\Omega }_\mathrm{m}=0.3,`$ $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ cosmology, $`t_0=11.1`$ Gyr. This expansion age may be compared with the estimated age of the oldest globular clusters, $`t_{}=12.8\pm 1`$ Gyr ($`1\sigma `$ uncertainty; Krauss 1999). At first blush, then, our estimated Hubble constant leads to a universe younger than its oldest stars. Given this logical contradiction, our result obviously requires further scrutiny. We discuss a number of salient issues in this final section.
### 7.1 Why do we disagree with the $`H_0`$KP?
The $`H_0`$KP team reported $`H_0=71\pm 6\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ (Mould et al. 2000). However, of their reported 9% ($`1\sigma `$) error, approximately 6.5% is systematic error due mainly to uncertainty in the distance to the LMC. This systematic error affects our value in precisely the same way as theirs, and thus should not be considered in comparing our $`H_0`$ estimates. The $`1\sigma `$ random error in the $`H_0`$KP Hubble constant is $`4.4\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1,`$ corresponding to a $`2\sigma `$ error of $`9\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1.`$ Thus, the $`H_0`$KP estimate overlaps with ours only at the very edges of our respective $`2\sigma `$ error bars (i.e., at $`80\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$). Since we have used the same Cepheid data set to arrive at our estimates, and therefore share much of their random error, this constitutes a significant disagreement.
There are two principal causes of this disagreement. The first is the difference in Cepheid calibration as discussed in §3. Our OGLE-based PL calibration produces distances that are smaller by $`5\%`$ on average, primarily because the OGLE calibration yields larger reddening and thus larger extinction estimates. When applied to the $`H_0`$KP distances and their procedure for estimating the Hubble constant, the OGLE calibration should bring their value up to $`75\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1,`$ closer to the value we derive, though still inconsistent, given that we use the same data.
The second source of disagreement is more fundamental: the different strategies we have adopted for determining $`H_0.`$ The $`H_0`$KP approach (referred to as “Method I” in §2) has been to use the Cepheid galaxies as calibrators for secondary distance indicators (DIs), especially Type Ia Supernovae (SN Ia), the Tully-Fisher and Fundamental Plane relations, and the SBF relation. The secondary DIs are then applied to galaxies at much larger distances than the Cepheid galaxies themselves, typically in the $`3000`$$`10,000\mathrm{km}\mathrm{s}^1`$ range, where peculiar velocities can be largely neglected. In contrast, we have derived $`H_0`$ from the Cepheid galaxies themselves, correcting the effects of non-Hubble motions with velocity field models (“Method II” of §2).
The two strategies are subject to different pitfalls. Method I can go awry with the propagation of Cepheid errors into the calibration of the secondary DIs. Such errors might occur for a variety of reasons, the foremost being the difficult nature of the measurements. SN Ia are often historical, i.e., occurred many years or decades ago, and the data on their brightnesses may not be consistent with modern methods. And yet, such historical SN Ia must be used in the calibration procedure, SN Ia being rare events and Cepheid galaxies being few in number. Four out of six SN Ia calibrated by Gibson et al. (2000) occurred prior to 1990, and two of these occurred prior to 1975. This small sample of calibrators introduces the biggest uncertainty in the SN-based $`H_0`$ (Suntzeff et al. 1999). The calibration of the SBF and Fundamental Plane (FP) methods suffer from another problem: Cepheids are found in late-type spiral galaxies, whereas SBF applies best, and FP applies only, to early type galaxies. Consequently, the absolute calibration of the FP relation using Cepheid distances (Kelson et al. 2000) must be obtained indirectly, by assuming that the Cepheid calibrators and the FP ellipticals are members of a group lying at a common distance. The SBF relation has been calibrated using a small number of spirals with prominent bulges (Ferrarese et al. 2000a), but possible stellar population differences between spiral bulges and ellipticals make the validity of this calibration uncertain (Tonry et al. 1997; TBAD00).
Neither of the above problems applies to the $`H_0`$KP calibration of the Tully-Fisher relation by Cepheid galaxies (Sakai et al. 2000). However, Sakai et al. obtained their value of $`H_0`$ from a single $`I`$ band Tully-Fisher data set, that of Giovanelli and collaborators (e.g., Giovanelli et al. 1997). Although this data set is of high quality, there are a number of other large Tully-Fisher data sets of recent vintage that were not considered by Sakai et al., such as those collected in the the Mark III Catalog (Willick et al. 1997a). Tully-Fisher measurements are prone to systematic differences in velocity width and photometric measurement conventions, and application of the Cepheid-calibrated Tully-Fisher relation to a wider range of Tully-Fisher data sets is needed; an important first step in this direction has been taken by Tully & Pierce (1999).
A coincidence of multiple secondary DI miscalibrations is very unlikely, at best. Nevertheless, our Method II analysis does not involve secondary DI’s and cannot suffer from the problem of propagated Cepheid calibration errors.<sup>12</sup><sup>12</sup>12It does, of course, potentially suffer from calibration error in the Cepheid PL relation, but this is essentially the problem of the LMC distance. However, our approach is vulnerable to inaccurately modeled peculiar velocities because we measure $`H_0`$ locally (see the discussion in §2). Our peculiar velocity models are “state-of-the-art,” especially the IRAS models, and we have optimized them with respect to the SBF data set, which is the best current sample for constraining the local velocity field. It is evident from our Hubble diagrams, however, that our velocity models are not perfect. A clear indication of this is the ridge of $`5`$ galaxies that lie well below our best-fit Hubble line in Figures 6 and 7. These are objects that lie within, and in the background of, the Virgo-Ursa Major region, and that are falling in toward Virgo or Ursa Major at higher velocity than predicted by the model. None of these objects deviates from our model at more than the $`2\sigma `$ level—indeed, the only $`3\sigma `$ deviant point is NGC 1326A, which is above the $`H_0=100\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ line at a distance of 16.4 Mpc—but it is still disquieting to see the large scatter at distances near and beyond Virgo. This scatter does not invalidate our approach, but reminds us that peculiar velocities are not fully accounted for in our model, and that some caution is needed in interpretation.
### 7.2 Considerations for future work
The discussion above shows that the debate on the Hubble constant will continue. Its ultimate resolution will require that several key needs are met:
1. More nearby galaxies with accurate Cepheid distances.— Our local measurement of $`H_0`$ could be greatly improved with a larger, and more uniformly distributed, sample of Cepheid galaxies. It is to be hoped that further observations by the HST, and later by NGST, will yield such samples. Indeed, such measurements would enable a better grasp of the systematic error in Method II values of $`H_0`$, along the lines of an angular variance analysis suggested by Turner, Cen & Ostriker (1992).
2. Improved models for the local velocity field.— The IRAS model presented here is a reasonable but imperfect fit to both the SBF and the Cepheid distances. Work is currently under way by WNSB to test enhancements of the model. In particular, we intend to further investigate the effects of nonlinear dynamics and nonlinear bias. Another dynamical variable that needs to be better understood is the small-scale velocity noise and its dependence on galaxy density. If these efforts succeed, the uncertainty in $`H_0`$ due to peculiar velocities will be reduced.
3. A re-calibration of secondary DIs.— As noted in §2, both the “distant” and “local” strategies for measuring $`H_0`$ should be pursued, and eventually they should agree. We have pointed to a few issues where the $`H_0`$KP calibrations of secondary DIs could contain subtle errors. For further investigation, one could use the Cepheid distances to calibrate the Tully-Fisher relations for samples not considered by Sakai et al. 2000, in particular, those tabulated in the Mark III Catalog (Willick et al. 1997a), especially after possible calibration errors in that and other catalogs have been corrected via comparison with the uniform, all-sky Shellflow survey recently presented by Courteau et al. (2000).
4. Further exploration of the “Hubble bubble”.— The approach of this paper, and, to a lesser extent, that of the $`H_0`$KP, could overestimate the Hubble constant if the local universe ($`d\stackrel{<}{}30h^1\mathrm{Mpc}`$) is expanding more rapidly than the global average, as has been suggested by Zehavi et al. (1998) on the basis of SN Ia data within $`10,000\mathrm{km}\mathrm{s}^1.`$ Such a situation is certainly possible on theoretical grounds; a fractional mass fluctuation $`\delta _M`$ within a sphere of radius $`R`$ produces a deviation of the Hubble constant within that sphere of $`\delta H_0/H_0\mathrm{\Omega }_\mathrm{m}^{0.6}\delta _M/3`$ (see Turner, Cen & Ostriker 1992 and Tomita 2000 for more detailed theoretical analyses). Typical mass fluctuations on scales $`R\stackrel{>}{}10h^1\mathrm{Mpc}`$ are $`\stackrel{<}{}1`$ in most cosmological scenarios, so that one would expect that the local value of $`H_0`$ on a $`10h^1\mathrm{Mpc}`$ scale to deviate by $`10`$–20% from the global value if $`\mathrm{\Omega }_\mathrm{m}0.3,`$ with smaller deviations on larger scales. However, for our local value of $`H_0`$ to exceed the global value, it would be necessary for our local neighborhood to be underdense relative to the mean. And yet, the 1.2 Jy IRAS density field suggests that the opposite is true—our local region within 20 Mpc is overdense relative to the volume within 100 Mpc well that is well-sampled by IRAS. Moreover, recent Tully-Fisher data (Dale & Giovanelli 2000) do not support the claim of Zehavi et al. (1998) for a local Hubble bubble. It is prudent to view this issue as an open one for now, and to continue to test the relationship between the local and distant Hubble flow. Such tests do not require DIs that are absolutely calibrated, and thus are independent of the question of $`H_0`$ itself.
5. More accurate absolute calibration of the PL relation.— As discussed in §2, the largest systematic error in the analysis of this paper is due to uncertainty in the zero point of the Cepheid PL relation, which is itself due to uncertainty in the distance to the LMC. We have adopted the “canonical” value $`\mu _{\mathrm{LMC}}=18.50`$ in this paper, the same as that adopted by the $`H_0`$KP. Our $`H_0`$ estimate could be range between $`78`$ and $`98\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ depending on the distance to the LMC chosen.
It is, finally, worth taking a moment to consider the question, What if $`H_0`$ really is as large as, say, $`90\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1\mathrm{?}`$ Would that constitute a “crisis for Big Bang cosmology,” a claim that has been heard in some quarters? The answer, for the moment, is clearly “No,” for two reasons. First, the globular clusters could still be as young as $`t_{}=10`$ Gyr. Such a young age is unlikely but possible at the few percent level (Krauss 1999). If this were the case, then $`t_0>t_{}`$ for $`H_0=90\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ and $`\mathrm{\Omega }_\mathrm{m}=0.3,`$ $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7.`$ One would then require only that the globular clusters formed very shortly ($`\stackrel{<}{}10^8`$ yr) after the Big Bang, which is not impossible. Second, even if $`t_{}=13`$ Gyr is correct, one can obtain $`t_0t_{}`$ for $`H_0=90\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ if $`\mathrm{\Omega }_\mathrm{m}=1\mathrm{\Omega }_\mathrm{\Lambda }0.13.`$ A universe of such low density has not been ruled out. In short, a Hubble constant $`90\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ does not pose an insurmountable problem for Big Bang cosmology so long as the ages of the oldest stars and the values of the parameters $`\mathrm{\Omega }_\mathrm{m}`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ remain poorly determined.
### 7.3 Summary
We have presented a new determination of the Hubble constant using Cepheid PL data published by the $`H_0`$KP. Rather than use the nearby ($`d\stackrel{<}{}20`$ Mpc) Cepheid galaxies as calibrators for secondary distance indicators, which are then applied to more distant ($`d\stackrel{>}{}50`$ Mpc) galaxies for which peculiar motions are fractionally small (the $`H_0`$KP strategy), we use Cepheid galaxies directly to measure $`H_0.`$
We first redetermined the Cepheid galaxy distances using a calibration of the PL relation derived from a large sample of LMC Cepheids presented by the OGLE group. Our absolute PL calibration assumed $`\mu _{\mathrm{LMC}}=18.5.`$ (We reemphasize that the $`H_0`$KP group will shortly present their own revision of Cepheid distances in light of the OGLE LMC Cepheid data \[Madore & Freedman 2000, in preparation; Freedman et al. 2000, in preparation\].) We then presented two models of the local peculiar velocity field. The first was obtained from the IRAS galaxy density field using the linear relation between large-scale velocity and density fields and the assumption that IRAS galaxies trace the mass density field up to a linear biasing factor $`b.`$ The IRAS model applies in the Local Group reference frame. The second was the phenomenological model of TBAD00, which applies in the CMB reference frame. The Tonry model assumes the local velocity field is dominated by infall to the Virgo and Great Attractors, along with a dipole and quadrupole term. We optimized each model by fitting it, using the maximum likelihood VELMOD algorithm of WSDK and WS, to a calibration-free 281-galaxy subset of the SBF sample of Tonry et al. (2001), currently the most accurate set of relative-distances for galaxies in the nearby ($`cz\stackrel{<}{}3000\mathrm{km}\mathrm{s}^1`$) universe. In the case of the IRAS model, this optimization consisted mainly of determining the value of $`\beta =\mathrm{\Omega }_\mathrm{m}^{0.6}/b,`$ which was found to be $`0.38\pm 0.06`$ ($`1\sigma `$ error), with $`0.2\beta 0.5`$ allowed at the $`3\sigma `$ level. For the Tonry model the optimization involved constraining several parameters that determine the influence of the Virgo and Great Attractors. The velocity model fits used only relative distances for the SBF galaxies and thus in no way prejudiced our subsequent determination of $`H_0.`$
We then applied the IRAS and Tonry velocity models to 27 Cepheid galaxies, again using the VELMOD algorithm, with the one remaining free parameter now being the Hubble constant. This yielded $`H_0=85\pm 5\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1,`$ essentially independent of the value of $`\beta ,`$ when the IRAS velocity field was used. When the Tonry model was used we obtained $`H_0=92\pm 5\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1.`$ The quoted random errors are at the 95% confidence level. The IRAS model produced a better fit likelihood than the Tonry model, and a Hubble diagram with markedly less scatter. We thus favor the result from IRAS, and adopt the IRAS value of $`H_0`$ for our final conclusion. This value is significantly larger than the $`H_0`$KP result, $`H_0=71\pm 6\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ (Mould et al. 2000).
Until Method I and Method II analyses give consistent results, we find it untenable to state that $`H_0`$ is known to within $`10\%`$. We discussed at length in §7.1 several possible reasons for the difference, as well as (§7.2) a number of lines of further investigation needed to clarify the issue. Should the larger value we quote here turn out to be correct, it would be difficult to reconcile the expansion age of the universe, $`t_0=11.1`$ Gyr for $`\mathrm{\Omega }_\mathrm{m}=0.3,`$ $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7,`$ with the estimated ages of oldest globular clusters, $`12.8\pm 1`$ Gyr (Krauss 1999). However, the remaining uncertainty in these stellar ages, as well as in the cosmological density parameters, is sufficiently large that a Hubble constant as large or even somewhat larger than what we have argued for here does not yet pose a logical inconsistency for Big Bang cosmology.
Sadly, Jeff died in a car accident only two days after we (PB and JAW) first huddled over the APJ referee’s report. I’d like to think that some part of Jeff’s wisdom shows in my own halting revisions to this paper. I would like to thank Stephane Courteau & Michael Strauss for helping with the completion of this work; I am also grateful for Tod Lauer’s referee’s report. I would have liked to thank JAW for his time—so valuable to me now. Jeff’s original acknowledgements follow below. This paper would not have been possible without the generous assistance of a number of individuals who shared data and expertise. We are particularly grateful to Wendy Freedman for timely input regarding the evolving calibration of the PL relation and helpful comments on an initial draft of this paper. $`H_0`$KP team members Laura Ferrarese and Brad Gibson also provided useful input on calibration issues. In addition, we acknowledge the entire $`H_0`$KP team for their epochal achievement in putting together the HST Cepheid database. Pierre Lanoix is thanked for his efforts in assembling and maintaining the Extragalactic Cepheid Database. Jeff Newman kindly provided his Cepheid data for NGC 4258 in advance of publication. Special thanks are due to Vijay Narayanan and Michael Strauss for producing IRAS-predicted peculiar velocities and densities used in this paper, and for helpful comments on the paper. Finally, we owe a large debt to John Tonry and his collaborators for their fine SBF data set, which enabled us to constrain the local velocity field. We are particularly grateful to John Blakeslee who shared the SBF data, and his wisdom on how to use it, with us prior to its publication. JAW was supported by a Cottrell Scholarship of Research Corporation and a Terman Fellowship from Stanford University, and, during the initial phase of this work, NSF grant AST96-17188.
## Appendix A Calculation of random distance errors
Although the systematic zero point errors in the Cepheid PL relation dominate the distance error budget, the weights we assign to the galaxies in the Hubble constant fit must be determined by random error only. To calculate this error we must account for how PL scatter couples with the reddening determination. The result of this coupling is that the distance error is markedly larger than one might naively expect.
We denote the I and V band galaxy (random) distance modulus errors $`\delta \mu _I`$ and $`\delta \mu _V.`$ These errors are due to PL scatter, and we assume them to be distributed as Gaussian random variables with mean zero and rms dispersions $`\sigma _I/\sqrt{N_{\mathrm{Ceph}}}`$ and $`\sigma _V/\sqrt{N_{\mathrm{Ceph}}},`$ where $`\sigma _I`$ and $`\sigma _V`$ are the I and V band rms PL scatter. The modulus errors induce an error $`\delta (VI)=\delta \mu _V\delta \mu _I`$ in the mean $`(VI)`$ color, which leads in turn to an error in the assumed reddening given by
$$\delta E(BV)=\frac{\delta (VI)}{R_VR_I},$$
(A1)
where the $`R_XA_X/E(BV)`$ are given in Table 2. The corresponding error in the distance modulus we assign to the galaxy is
$$\delta \mu ^{red}=\frac{1}{2}(R_V+R_I)\delta E(BV)=\frac{1}{2}\frac{R_V+R_I}{R_VR_I}\delta (VI),$$
(A2)
where we have assumed that the V and I band data are equally weighted in the distance determination, as they are in the present application. The negative sign in Eq. (A2) arises because PL errors that make the galaxy appear redder result in its being overcorrected for extinction, i.e., assigned too small a distance modulus.
Eq. (A2) represents only that part of the distance modulus error due to reddening determination error. To it we must add the direct error due to PL scatter, $`\delta \mu ^{direct}=(\delta \mu _V+\delta \mu _I)/2.`$ Adding the two, and writing the result in terms of the independent variables $`\delta \mu _V`$ and $`\delta \mu _I,`$ we obtain the overall galaxy distance modulus error,
$$\delta \mu =\frac{1}{2}\left[(1\mathrm{\Delta }_f)\delta \mu _V+(1+\mathrm{\Delta }_f)\delta \mu _I\right],$$
(A3)
where we have defined
$$\mathrm{\Delta }_f\frac{R_V+R_I}{R_VR_I}=4.0625.$$
(A4)
The numerical value of $`\mathrm{\Delta }_f`$ follows from the values of $`R_V`$ and $`R_I`$ in Table 2. Eq. (A3) tells us that $`\delta \mu `$ is normally distributed with mean zero and rms dispersion
$$\sigma _\mu =\sqrt{\left(\frac{1\mathrm{\Delta }_f}{2}\right)^2\sigma _V^2/N_{\mathrm{Ceph}}+\left(\frac{1+\mathrm{\Delta }_f}{2}\right)^2\sigma _I^2/N_{\mathrm{Ceph}}}.$$
(A5)
If we furthermore take $`\sigma _V\sigma _I=\sigma _{\mathrm{Ceph}},`$ the last equation reduces to
$$\sigma _\mu =\frac{\sigma _{\mathrm{Ceph}}}{2\sqrt{N_{\mathrm{Ceph}}}}\sqrt{(1\mathrm{\Delta }_f)^2+(1+\mathrm{\Delta }_f)^2}=2.96\frac{\sigma _{\mathrm{Ceph}}}{\sqrt{N_{\mathrm{Ceph}}}}.$$
(A6)
The corresponding expression for the rms error in the mean reddening for the galaxy is
$$\sigma _{E(BV)}=\frac{\sqrt{2}}{R_VR_I}\frac{\sigma _{\mathrm{Ceph}}}{\sqrt{N_{\mathrm{Ceph}}}}=1.1\frac{\sigma _{\mathrm{Ceph}}}{\sqrt{N_{\mathrm{Ceph}}}}.$$
(A7)
We use Eqs. (A6) and (A7), with $`\sigma _{\mathrm{Ceph}}=0.15`$ mag, to calculate the distance modulus and reddening errors in columns 2 and 4 of Table 5. (The distance errors in column 3 are computed directly from the modulus errors in the standard fashion.)
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# Iordanskii and Lifshitz-Pitaevskii Forces in the Two-Fluid Model
## 1 A HIERARCHY OF LENGTH SCALES
A comprehensive theory of vortex dynamics in quantum fluids must address phenomena occurring at three different length scales. At the smallest, most microscopic scale, a fully quantum mechanical treatment of a vortex, including its internal structure, interaction with elementary excitations such as phonons and rotons, and interaction with disorder, is required
At the next—what we shall refer to as intermediate—length scale, one would like to regard a vortex as a classical object, subject to a classical equation of motion. Of course, there is no guarantee that we can do this, and, indeed, certain pathologies result from our insistence in doing so, but experience has shown that this classical picture is extremely successful.
There are different ways to formulate the classical approach. If we let $`𝐑(t)`$ denote the position in the $`xy`$ plane of the center of an isolated straight vortex line, say, as a function of time, we could hypothesize a phenomenological equation of motion of the form
$$M\frac{d^2𝐑}{dt^2}=\eta \frac{d𝐑}{dt}\gamma \frac{d𝐑}{dt}\times 𝐞_z+𝐟_\mathrm{p}(𝐑)+𝐟_\mathrm{d}(𝐑).$$
(1)
Here we have taken the circulation vector $`𝐊`$ of the vortex, a vector parallel to the vortex with a magnitude equal to the circulation, to be along the $`z`$ direction. The first two terms on the right-hand-side of (1) are to include all forces linear in the vortex velocity. The coefficients $`M`$, $`\eta `$, and $`\gamma `$, which describe the vortex effective mass per unit length, viscous damping force per unit length, and nondissipative transverse force per unit length, respectively, are to be determined from a microscopic theory, as are the “pinning” and “driving” forces per unit length $`𝐟_\mathrm{p}`$ and $`𝐟_\mathrm{d}`$. The latter may depend on the normal and superfluid densities and velocities, and therefore on the vortex position $`𝐑`$, but do not depend on the vortex velocity. In a superfluid, $`𝐟_\mathrm{p}`$ might describe the force exerted by an externally imposed wire, and $`𝐟_\mathrm{d}`$ would include the superfluid-velocity-dependent part of the Magnus force (see below) and possibly other vortex-velocity-independent contributions.
At the most macroscopic scale one needs to understand how the forces arising at the intermediate scale act to determine the bulk, experimentally observable properties of quantum fluids. The length scale that usually defines this regime is the characteristic inter-vortex distance, and one is interested in coarse-grained properties of the quantum fluid at scales larger than that distance. In superfluids, this is the regime where the concept of mutual friction applies. Mutual friction refers to a momentum transfer, with both longitudinal and transverse components, between the normal and superfluid parts of a quantum fluid. Although such an interaction is absent in the two-fluid model itself, the presence of vortices mediate a bulk momentum exchange. Similarly, in superconductors one needs to understand how the Lorentz force, pinning forces, and other intermediate-scale forces conspire to determine, say, the Hall effect in the mixed state, which depends on the collective motion of a macroscopic number of vortices.
Our recent work has focused mostly on the intermediate length-scale regime, namely, the determination of the various forces that act on vortices. After a brief general review of that work we shall discuss our new results on the problem of the Iordanskii and Lifshitz-Pitaevskii forces, which have been subjects of considerable controversy.
## 2 THE TAN THEORY
The Thouless-Ao-Niu (TAN) theory and its generalization to ultraclean type-II superconductors start with the appropriate exact many-body Hamiltonian and include a pinning potential centered at $`𝐑`$, which is also the position of the vortex. The vortex is then dragged with a velocity $`𝐕`$, and a transverse force $`𝐟_{}=𝐞_z\times 𝐕𝑑𝐥𝐣_\mathrm{c}`$ is found. Here $`𝐣_\mathrm{c}`$ is the canonical momentum density, and the line integral is taken around a large circle enclosing the vortex. Evidently, the transverse force is independent of the detailed microscopic structure of the vortex and its interaction with the normal fluid.
In a neutral Bose or Fermi superfluid, $`𝐣_\mathrm{c}=\rho _\mathrm{s}𝐯_\mathrm{s}+\rho _\mathrm{n}𝐯_\mathrm{n},`$ and therefore
$$𝑑𝐥𝐣_\mathrm{c}=\rho _\mathrm{s}K_\mathrm{s}+\rho _\mathrm{n}K_\mathrm{n}.$$
(2)
In the normal fluid component, viscosity causes vortex-like circulation to diffuse away to the outer boundary, so $`K_\mathrm{n}`$ vanishes in equilibrium, and
$$𝐟_{}=\rho _\mathrm{s}𝐊_\mathrm{s}\times 𝐕.$$
(3)
In a charged superfluid or superconductor, however, the Hamiltonian contains a current-current interaction term that modifies the velocity operator. The canonical momentum density can be expressed in terms of the physical, gauge-invariant momentum density $`𝐣=𝐣_\mathrm{c}+\frac{e}{c}n𝐀`$, resulting in
$$𝑑𝐥𝐣_\mathrm{c}=𝑑𝐥\left(𝐣\frac{e}{c}n𝐀\right)=\rho \mathrm{\Phi }_0,\mathrm{\Phi }_0hc/2e.$$
(4)
The second equality in (4) follows from the Meissner effect, which causes $`𝐣`$ to vanish at large distances, and from flux quantization. Therefore,
$$𝐟_{}=\rho 𝐊_\mathrm{s}\times 𝐕,$$
(5)
where $`\rho =\rho _\mathrm{s}+\rho _\mathrm{n}`$ is the total mass density of the fluid.
## 3 A HIERARCHY OF CONTROVERSIES
We turn now to a brief discussion of some of the controversies in the theory of intermediate-scale vortex dynamics, focusing on transverse forces, and organized according to the complexity of the quantum fluid in question.
It seems appropriate to start by recalling a result of classical hydrodynamics, the Magnus force: When a vortex with circulation $`𝐊`$ is dragged with velocity $`𝐕`$ through an ideal fluid of mass density $`\rho `$, a transverse force per unit length $`𝐟_{}=\rho 𝐊\times 𝐕`$ acts on the object doing the dragging. The force is similar to the lift force on an airplane wing. It is clear from Galilean invariance that if the fluid far from the vortex is not at rest, but has a velocity $`𝐯`$, then the force is instead
$$𝐟_{}=\rho 𝐊\times (𝐕𝐯).$$
(6)
By analogy, it is reasonable in a neutral Bose superfluid at zero temperature to expect that $`𝐟_{}=\rho _\mathrm{s}𝐊_\mathrm{s}\times (𝐕𝐯_\mathrm{s}),`$ where $`K_\mathrm{s}`$ is the quantized circulation. Of course, writing $`\rho _\mathrm{s}`$ here instead of $`\rho `$ is arbitrary, because these are equal at zero temperature.
The controversy begins in the finite-temperature neutral Bose superfluid, because Galilean invariance allows for a transverse force of the form
$$𝐟_{}=a𝐊_\mathrm{s}\times (𝐕𝐯_\mathrm{s})+b𝐊_\mathrm{s}\times (𝐕𝐯_\mathrm{n}),$$
(7)
where $`a`$ and $`b`$ are parameters. The first term in (7) has a simple classical interpretation: It describes a Magnus-type force originating from the superfluid component of the fluid (remembering that it is the superfluid that is circulating). Thus classical reasoning would suggest that $`a=\rho _\mathrm{s}`$, and Wexler has recently given an elegant proof of this. To our knowledge there have been no criticisms of Wexler’s theory.
From this classical point-of-view, again keeping in mind that it is the superfluid that is circulating here, the second term in (7) would be of a nonhydrodynamic origin. If present, it would describe a transverse interaction between the vortex and the excitations—phonons and rotons—of the fluid.
Iordanskii predicted just such an interaction with phonons, and Lifshitz and Pitaevskii (following earlier work by Hall and Vinen) predicted one with rotons. Note, however, that in these works $`𝐯_\mathrm{s}`$ and $`𝐯_\mathrm{n}`$ refer to flow velocities near the vortex line, whereas our quantities are asymptotic values. But the TAN result (3) implies $`a+b=\rho _\mathrm{s}`$. When combined with Wexler’s theory, this implies $`b=0`$, i.e., that there are no transverse Iordanskii and Lifshitz-Pitaevskii forces. This apparent discrepancy has motivated us to understand better the interaction between a quantized vortex and the normal fluid in a neutral Bose system, and to calculate the coefficient $`b`$ in (7) directly. We shall return to this direct calculation below in Section 4
The next controversy concerns the transverse force in a neutral Fermi superfluid, also described in the TAN theory. According to Kopnin and Kravtsov, low-energy quasiparticles in the vortex core also contribute to the transverse force linear in $`𝐕`$, a contribution not found by TAN.
Vortices in superconductors inherit all of the above controversies and have additional complexity of their own. We will not have space to discuss them further.
## 4 VORTEX DYNAMICS IN THE TWO-FLUID MODEL
Two ingredients are needed for a microscopic theory of the coefficient $`b`$ in Eqn. (7). First, one has to solve a vortex-excitation scattering problem. It is probably correct to use the Gross-Pitaevskii equation to do this, even though mean field theory is not expected to hold inside the vortex, because the scattering potential is long-ranged. In fact, the potential is sufficiently long-ranged, and the forward scattering sufficiently singular, that there are several mathematically incorrect scattering calculations present in the literature.
In the scale of lengths greater than the size of the vortex core and less than the mean free path for scattering of excitations by one another, if there is such a range, the dynamics of excitations moving in the slowly varying background of the rotating superfluid is well defined, and in this region the traditional arguments give a transverse force on the excitations equal to $`\rho _\mathrm{n}𝐊_\mathrm{s}\times (𝐕𝐯_\mathrm{n})`$, where $`𝐯_\mathrm{n}`$ is the average normal fluid velocity at a distance of the order of the mean free path from the vortex core, as well as a dissipative longitudinal force.
In the region well beyond a mean free path from the vortex core, flow velocities are varying slowly and two-fluid hydrodynamics, which incorporates the basic conservation laws, provides an accurate description. In particular, the two-fluid version of the Navier–Stokes equation gives the force–momentum balance. In this region we have studied how forces generated within a mean free path of the vortex core can be transmitted to large distances. Within a linearized approximation to the two-fluid model the force due to the superfluid flow relative to the vortex is just the ordinary superfluid Magnus force $`\rho _\mathrm{s}𝐊_\mathrm{s}\times (𝐕𝐯_\mathrm{s})`$, where $`𝐯_\mathrm{s}`$ is the asymptotic value of the superfluid velocity. This is given for superfluids, as it is for ordinary fluids, by the cross terms, between circulation and the superfluid flow, in the momentum flow tensor and the Bernoulli pressure. There seems to be no possible modification of this even for a fermionic superfluid, such as is suggested by the work of Kopnin and Kravtsov, if there is no bulk interaction with a stationary background.
For the normal fluid contributions our considerations are quite similar to those of Hall and Vinen. In the linear regime the force is transmitted by the viscous force (and an accompanying pressure term), which is generated by a flow velocity in the direction of the force that increases as the logarithm of the distance from the vortex core. The asymptotic value of the normal fluid velocity is given by a combination of the value $`𝐯_\mathrm{n}`$ close to the core, and this logarithmically growing term. The logarithmically growing term has to be cut off at a distance $`R_c`$ which is determined by the Reynolds number of the flow, by the spacing between vortices, or by the diffusion length at the frequency of vortex oscillation. In the limit considered by TAN, with vanishingly small normal fluid velocity, and a single vortex in an infinite medium, this cut-off goes to infinity, and $`𝐯_\mathrm{n}`$ is negligibly small in comparison, so the asymptotic velocity is determined by the logarithmic term, and the force is parallel to the normal fluid flow.
Realistically, the logarithm which the cut-off introduces is not necessarily large compared with other parameters in the problem. Even when the Iordanskii force is not put in explicitly at the inner boundary, there is some transverse force due to normal fluid flow, but this is proportional to $`[\mathrm{ln}(R_c/\lambda )]^2v_\mathrm{n}`$, where $`\lambda `$ is the mean free path. It is not clear how to match conditions close to the core with conditions in the hydrodynamic region, but if it is assumed that the Iordanskii force causes $`𝐯_\mathrm{n}`$ close to the core to match the viscous force in the hydrodynamic region, and to be parallel to it, the relative importance of the logarithmic term will depend on the ratio of the kinematic viscosity $`\eta /\rho _\mathrm{n}`$ to the quantum of circulation—a large kinematic viscosity will diminish the importance of the logarithmic term in the flow velocity.
## ACKNOWLEDGMENTS
This work was supported in part by NSF grant No. DMR-9813932, by a Research Innovation Award from the Research Corporation, and by a Sarah Moss Fellowship.
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# Divergence of a quantum thermal state on Kerr space-time
## Abstract
We present a simple proof, using the conservation equations, that any quantum stress tensor on Kerr space-time which is isotropic in a frame which rotates rigidly with the angular velocity of the event horizon must be divergent at the velocity of light surface. We comment on our result in the light of the absence of a ‘true Hartle-Hawking’ vacuum for Kerr.
preprint: OUTP-00-23-P:gr-qc/0005108
One of the fundamental results of quantum field theory in Kerr space-time is a theorem of Kay and Wald , that there does not exist a Hadamard state which respects the symmetries of the space-time and is regular everywhere outside and on the event horizon. This means that there is no ‘true Hartle-Hawking’ (HH) vacuum on Kerr space-time. The HH vacuum on Schwarzschild black holes has been extensively studied, since it is the state which lends itself most readily to numerical computations, due to its regularity and high degree of symmetry. For the same reason the construction of a state on Kerr space-time with most (but not all) of the properties of the HH state remains an important open question.
The construction of quantum states on Kerr space-time is a delicate matter due to the presence of the super-radiant modes (see for details of this procedure). There is a consensus in the literature that the (past) Boulware vacuum $`|B^{}`$ is defined by taking a basis of modes which are positive frequency with respect to the Killing vector $`/t`$ at $`^{}`$ and with respect to the Killing vector $`/t+\mathrm{\Omega }_H/\varphi `$ (where $`\mathrm{\Omega }_H`$ is the angular velocity of the horizon) at $`^{}`$. This state corresponds to an absence of particles coming up from $`^{}`$ or in from $`^{}`$, and contains, at $`^+`$, an outward flux of particles due to the Unruh-Starobinskii effect (spontaneous emission in superradiant modes). There are two attempts in the literature to define a state analogous to the HH state , one due to Frolov and Thorne and the other due to Candelas, Chrzanowski and Howard which we shall denote by $`|FT`$ and $`|CCH`$, respectively. We refer the reader to for details of the construction of these states, which are not important here. Both these candidate states are thermal in nature (although the energy with respect to which the modes are thermalized differ), but they have different symmetry and regularity properties. The state $`|FT`$ is invariant under simultaneous $`t`$, $`\varphi `$ reversal, while $`|CCH`$ is not.
Our interest in this letter is in the difference in expectation values of the quantum stress tensor in the Boulware and candidate HH vacua. By considering such a difference any renormalization terms cancel and we may deal effectively deal with the bare operators. We shall denote this tensor as $`\widehat{T}_{\mu \nu }_{HHB}`$ regardless of which candidate HH vacuum we are considering, since our result in this paper is quite general and does not depend on the details of the construction of the thermal states. In , $`\widehat{T}_{\mu \nu }_{HHB}`$ represents a precisely thermal atmosphere of quanta at the Hawking temperature which rotates rigidly with the same angular velocity as the event horizon, $`\mathrm{\Omega }_H`$. Therefore, $`\widehat{T}_{\mu \nu }_{HHB}`$ should be isotropic in a frame which also rotates rigidly with angular velocity $`\mathrm{\Omega }_H`$. This is in agreement with the calculation of , where the difference in expectation values of the stress tensor in the Boulware $`|B^{}`$ and $`|CCH`$ states was calculated at the event horizon for an electromagnetic field. They found that this tensor was isotropic in the Carter tetrad . Due to the rotation of the black hole, this is the same, at the event horizon, as the tensor being isotropic in a tetrad which is rigidly rotating with the same angular velocity as the black hole. The calculation in is valid only close to the event horizon, and does not provide any information about the isotropy (or otherwise) of $`\widehat{T}_{\mu \nu }_{HHB}`$ away from the horizon.
Frolov and Thorne then used an argument based on the details of the quantum field modes to show that this stress tensor will fail to be regular at the velocity of light surface $`𝒮_L`$. This is the surface, where, in order to co-rotate with the event horizon, an observer must travel at the speed of light. We shall now show that it is an elementary consequence of the conservation equations that $`\widehat{T}_{\mu \nu }_{HHB}`$, with the assumed isotropy property, will fail to be regular at $`𝒮_L`$. In addition, we will also find that $`\widehat{T}_{\mu \nu }_{HHB}`$ has the same, divergent, form at the event horizon as found by . Our metric has signature $`(+++)`$ and we use units in which $`G=c=\mathrm{}=k_b=1`$ throughout. Greek letters will denote co-ordinate components, while bracketed Roman letters denote tetrad components.
Firstly, we begin by writing the Kerr metric in the unconventional form
$$ds^2=\alpha ^2dt^2+\rho ^2\mathrm{\Delta }^1dr^2+\rho ^2d\theta ^2+\stackrel{~}{\mathrm{\Omega }}^2(d\varphi \mathrm{\Omega }dt)^2,$$
(1)
where $`\rho ^2=r^2+a^2\mathrm{cos}^2\theta `$, $`\mathrm{\Delta }=r^22Mr+a^2`$, with $`M`$ the mass of the black hole and $`a`$ the angular momentum per unit mass, as viewed from infinity. The other functions appearing in the metric (1) are given by:
$`\alpha ^2`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }\rho ^2}{\left(r^2+a^2\right)^2\mathrm{\Delta }a^2\mathrm{sin}^2\theta }},`$ (2)
$`\mathrm{\Omega }`$ $`=`$ $`{\displaystyle \frac{2Mra}{\left(r^2+a^2\right)^2\mathrm{\Delta }a^2\mathrm{sin}^2\theta }},`$ (3)
$`\stackrel{~}{\mathrm{\Omega }}^2`$ $`=`$ $`{\displaystyle \frac{1}{\rho ^2}}\left[\left(r^2+a^2\right)^2\mathrm{\Delta }a^2\mathrm{sin}^2\theta \right]\mathrm{sin}^2\theta .`$ (4)
Here $`\alpha `$ is the lapse function, vanishing on the event horizon $`r=r_H`$ (when $`\mathrm{\Delta }=0`$), and $`\mathrm{\Omega }`$ is the angular velocity of a locally non-rotating observer (LNRO), which is equal to the angular velocity $`\mathrm{\Omega }_H`$ of the event horizon when $`r=r_H`$. An observer who is rotating with angular velocity $`\mathrm{\Omega }_H`$ has the Lorentz factor $`\gamma `$ relative to a LNRO at the same values of $`r`$ and $`\theta `$, where
$$\gamma =\left(1v^2\right)^{\frac{1}{2}},v=\alpha ^1\left(\mathrm{\Omega }_H\mathrm{\Omega }\right)\stackrel{~}{\mathrm{\Omega }}.$$
(5)
At the event horizon, an LNRO has angular velocity $`\mathrm{\Omega }_H`$ and in this case $`\gamma =1`$, whilst at $`𝒮_L`$, $`v1`$ and $`\gamma \mathrm{}`$, as expected.
An orthonormal frame which rotates rigidly with angular velocity $`\mathrm{\Omega }_H`$ has basis 1-forms given by:
$`e_{(t)i}dx^i`$ $`=`$ $`\alpha ^1\gamma \left[\alpha ^2dt\stackrel{~}{\mathrm{\Omega }}^2(\mathrm{\Omega }_H\mathrm{\Omega })(d\varphi \mathrm{\Omega }dt)\right],`$ (6)
$`e_{(r)i}dx^i`$ $`=`$ $`(\rho ^2\mathrm{\Delta }^1)^{\frac{1}{2}}dr,`$ (7)
$`e_{(\theta )i}dx^i`$ $`=`$ $`\rho d\theta ,`$ (8)
$`e_{(\varphi )i}dx^i`$ $`=`$ $`\stackrel{~}{\mathrm{\Omega }}\gamma (d\varphi \mathrm{\Omega }_Hdt).`$ (9)
We shall consider a stress tensor $`\widehat{T}_{\mu \nu }_{HHB}`$ which is isotropic in this frame, so that with respect to the basis of 1-forms (9) the tetrad components are:
$$\widehat{T}_{(b)}^{(a)}_{HHB}=f(r,\theta )\mathrm{diag}\{3,1,1,1\},$$
(10)
where we have used the Killing vector symmetries of the geometry to assume that $`\widehat{T}_{\mu \nu }_{HHB}`$ does not depend on either $`t`$ or $`\varphi `$ . In order to find the unknown function $`f(r,\theta )`$, we solve the conservation equations for $`\widehat{T}_{\mu \nu }_{HHB}`$. The simplest way to do this, which avoids the calculation of Ricci rotation coefficients for the tetrad (9), is to convert the tetrad components back to Boyer-Lindquist co-ordinate components and then solve the conservation equations in the form :
$$_\nu \left(\widehat{T}_\mu ^\nu _{HHB}\sqrt{g}\right)=\frac{1}{2}\sqrt{g}\left(_\mu g_{\lambda \sigma }\right)\widehat{T}^{\lambda \sigma }_{HHB}.$$
(11)
The $`\mu =t`$ and $`\mu =\varphi `$ equations are trivial, and the $`\mu =r`$ and $`\mu =\theta `$ equations give, respectively,
$`_r\left(f(r,\theta )\rho ^2\mathrm{sin}\theta \right)`$ $`=`$ $`{\displaystyle \frac{1}{2}}f(r,\theta )\rho ^2\mathrm{sin}\theta _r\left(\mathrm{log}|\alpha ^8\gamma ^8\rho ^4\mathrm{sin}^2\theta |\right)`$ (12)
$`_\theta \left(f(r,\theta )\rho ^2\mathrm{sin}\theta \right)`$ $`=`$ $`{\displaystyle \frac{1}{2}}f(r,\theta )\rho ^2\mathrm{sin}\theta _\theta \left(\mathrm{log}|\alpha ^8\gamma ^8\rho ^4\mathrm{sin}^2\theta |\right).`$ (13)
These two equations are compatible and determine $`f(r,\theta )`$ up to an arbitrary constant. The result is
$$f(r,\theta )=k\alpha ^4\gamma ^4$$
(14)
where $`k`$ is an arbitrary constant. This implies that the tetrad components of $`\widehat{T}_{\mu \nu }_{HHB}`$ diverge as $`\mathrm{\Delta }^2`$ as the event horizon is approached, which is the behaviour found in .
In order to consider the regularity (or otherwise) of $`\widehat{T}_{\mu \nu }_{HHB}`$ at the event horizon or $`𝒮_L`$, we must first convert the tetrad components to Boyer-Lindquist components, since the tetrad (9) is not regular either at the event horizon or $`𝒮_L`$. The non-zero components of $`\widehat{T}_{\mu \nu }_{HHB}`$ are:
$`\widehat{T}_{tt}_{HHB}`$ $`=`$ $`\left[3\alpha ^2+\stackrel{~}{\mathrm{\Omega }}^2\left(4\gamma ^2\mathrm{\Omega }_H^23\mathrm{\Omega }^2\right)\right]f(r,\theta ),`$ (15)
$`\widehat{T}_{t\varphi }_{HHB}`$ $`=`$ $`\stackrel{~}{\mathrm{\Omega }}^2\left(4\gamma ^2\mathrm{\Omega }_H3\mathrm{\Omega }\right)f(r,\theta ),`$ (16)
$`\widehat{T}_{\varphi \varphi }_{HHB}`$ $`=`$ $`\stackrel{~}{\mathrm{\Omega }}^2\left(4\gamma ^23\right)f(r,\theta ),`$ (17)
$`\widehat{T}_{rr}_{HHB}`$ $`=`$ $`\rho ^2\mathrm{\Delta }^1f(r,\theta ),`$ (18)
$`\widehat{T}_{\theta \theta }_{HHB}`$ $`=`$ $`\rho ^2f(r,\theta ).`$ (19)
Using these components, the regularity of $`\widehat{T}_{\mu \nu }_{HHB}`$ on the event horizon can be considered by first transforming to Kruskal co-ordinates, as in . Considering the Kruskal components shows that $`\widehat{T}_{\mu \nu }_{HHB}`$ is not regular on either the past or the future event horizon, due to the $`\mathrm{\Delta }^2`$ divergence of $`f`$ as the horizon is approached. This is precisely the behaviour expected of the Boulware vacuum close to the event horizon, and does not preclude the possibility that a candidate HH state may be regular on some section of the event horizon. In it was argued that, of the two candidates for the analogue of the HH vacuum in Kerr, $`|FT`$ is regular only at the pole of the event horizon, and $`|CCH`$ is regular on the future (but not on the past) event horizon. Both these cases are compatible with our results here, provided that the divergences (where they exist) are of lower order than the $`\mathrm{\Delta }^2`$ expected for the Boulware vacuum.
The Boyer-Lindquist co-ordinate system is regular at $`𝒮_L`$, and so (19) reveals that the components of $`\widehat{T}_{\mu \nu }_{HHB}`$ diverge at least as fast as $`\gamma ^4`$ as the velocity of light surface is approached and $`\gamma \mathrm{}`$. Thus, if $`B^{}|\widehat{T}_{\mu \nu }|B^{}_{ren}`$ is regular at the velocity of light surface, then the expectation value of the stress tensor in the HH vacuum diverges there. This is certainly the case for slowly rotating black holes when $`𝒮_L`$ is far from the horizon and $`B^{}|\widehat{T}_{\mu \nu }|B^{}_{ren}`$ is given by the Unruh-Starobinskii effect. In fact, the Boulware vacuum is expected to be regular everywhere away from the event horizon (where it diverges). We conclude that if we have a state $`|H`$ such that $`\widehat{T}_{\mu \nu }_{HHB}`$ is isotropic in the tetrad (9), then the state $`|H`$ will not be regular on $`𝒮_L`$.
Thus we have shown that a stress tensor which is isotropic with respect to a frame which rotates rigidly with the angular velocity of the event horizon must diverge as $`\mathrm{\Delta }^2`$ as the event horizon is approached, and must also fail to be regular at $`𝒮_L`$. The question is therefore, is this isotropy condition likely to be satisfied by a physical stress tensor? The answer is far from obvious as the peak of the thermal spectrum at the Hawking temperature is at a wavelength comparable to the radius of curvature of the space-time near the horizon. In the absence of a full numerical calculation in Kerr, we note that the corresponding property of $`\widehat{T}_{\mu \nu }_{HHB}`$ for Schwarzschild black holes was conjectured in , and subsequently shown to hold to a good approximation in . The authors of used the thermal properties of their state $`|FT`$ to show that the corresponding stress tensor $`\widehat{T}_{\mu \nu }_{HHB}`$ is isotropic in the rigidly rotating tetrad, although in we have argued that, despite its attractive symmetry properties (in particular, simultaneous $`t`$, $`\varphi `$ reversal invariance) this state is fundamentally flawed. The alternative state $`|CCH`$, which is not invariant under simultaneous $`t`$, $`\varphi `$ reversal, is more workable, and satisfies the required isotropy property, at least near the event horizon . Note that, in contrast to the situation in Schwarzschild, the Boulware vacuum in Kerr is not invariant under simultaneous $`t`$, $`\varphi `$ reversal. Therefore, we would expect any analogue HH state such that $`\widehat{T}_{\mu \nu }_{HHB}`$ has our conjectured isotropy must also not be $`t`$, $`\varphi `$ reversal invariant.
It should be stressed that it is crucial to our analysis that the stress tensor $`\widehat{T}_{\mu \nu }_{HHB}`$ is isotropic in the rigidly rotating frame (9) rather than any other tetrad. For example, a stress tensor which is isotropic with respect to the Carter tetrad everywhere outside the event horizon will also display the divergence near the event horizon we have exhibited here, but will be regular elsewhere. This would be anticipated from the fact that the Carter tetrad rotates with an angular velocity which, although not the same as the angular velocity of the LNROs, decreases as we move away from the event horizon. Therefore, observers who are rotating with the same angular velocity as the Carter tetrad will always have a finite Lorentz factor relative to the LNROs (compare (5)). The point of is that only observers who have angular velocity $`\mathrm{\Omega }_H`$ (whatever their position outside the event horizon) will see an isotropic thermal distribution of particles, which is the reason for our conjectured isotropy condition here. The argument in uses the thermal properties of their state $`|FT`$, and does not apply to the other candidate HH state, $`|CCH`$, due to the different thermalization of the modes . We hope to return to this question subsequently, by examining the isotropy (or otherwise) of these candidate HH vacua using numerical calculations.
In the light of the Kay-Wald theorem , which states that there is no true HH state on Kerr spacetime, our result shows that, if it is possible to define a state on Kerr which has most of the properties of the HH state, then either the stress tensor $`\widehat{T}_{\mu \nu }_{HHB}`$ corresponding to this state will fail to be isotropic in the rigidly rotating tetrad, or the state will cease to be regular on $`𝒮_L`$. Furthermore, we have shown that the divergence at $`𝒮_L`$ is an elementary consequence of the conservation equations.
The work of E.W. is supported by a fellowship from Oriel College, Oxford, and she would like to thank the Department of Physics, University of Newcastle, Newcastle-upon-Tyne, for hospitality during the completion of this work.
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# Crossing probabilities on same-spin clusters in the two-dimensional Ising model
## 1 Introduction
The probability $`\pi _h(r)`$ of crossing on open sites inside a rectangle of aspect ratio $`r`$ has been measured at $`p_c`$ for several models of percolation in and . These simulations support hypotheses of universality and conformal invariance of this function $`\pi _h(r)`$ and of several others. Cardy’s contemporaneous work offered a prediction for $`\pi _h`$ using conformal field theory. His analytic expression agrees with the simulations within statistical errors and provides further support that these crossing probabilities are order parameters with the usual critical properties.
This paper presents similar evidence, both numerical and analytic, for crossing probabilities on same-spin clusters of the two-dimensional Ising model at criticality. These crossing probabilities are not traditional order parameters for the Ising model. It is not a priori obvious that they are not identically zero or one in any dimension $`d`$. In dimension two however the simulations carried out in indicates clearly that they are non-trivial functions and that they are likely to satisfy the same hypotheses (or even more restrictive ones) of universality and conformal invariance. Moreover physical quantities for which conformal field theory gives quantitative predictions that can be readily verified by simulation are always welcome.
Let a triangular lattice be oriented in such a way that sites on horizontal lines are at a distance of one mesh unit. A rectangle of height $`V`$ and width $`H`$ is superimposed on the lattice. The width $`H`$ is measured in mesh units but $`V`$ is the number of horizontal lines in the rectangle. A configuration of the Ising model has a horizontal crossing if there exists a path made of edges between nearest neighbor sites going from the leftmost inner column to the rightmost one and visiting only plus spins. (Because of the relative position of the rectangle and the lattice, the vertical columns are made of sites in zigzag.) Let $`\pi _h(V,H)`$ be the probability at the critical temperature of such a crossing. One defines similarly a vertical crossing and the associated probability $`\pi _v(V,H)`$. It is a well-known fact from percolation theory that, on the triangular lattice, a horizontal crossing on plus spins (or open sites) exist if and only if there is no vertical crossing on minus spins (closed sites). (One can convince oneself easily of this simple fact using a drawing.) Since plus and minus spins are equiprobable, this observation implies $`\pi _h(V,H)+\pi _v(V,H)=1`$. In this paper we will ultimately be interested in the limit $`\pi _h(r)=lim_{V,H\mathrm{},r=\sqrt{3}V/2H}\pi _h(V,H)`$ and similarly for $`\pi _v`$. The infinite lattice limit of the previous relation is known as the duality relation: $`\pi _h(r)+\pi _v(r)=1`$. If $`\pi _h`$ and $`\pi _v`$ are rotational invariants, the latter relation can be written as $`\pi _h(r)+\pi _h(1/r)=1`$. This duality relation will hold for other (regular) lattices if the functions $`\pi _h`$ and $`\pi _v`$ are universal but its discrete equivalent ($`\pi _h(V,H)+\pi _v(V,H)=1`$) hold strictly for neither the square nor the hexagonal lattices.
The present paper discusses both a prediction for $`\pi _h`$ extending Cardy’s approach and precise measurements of the function $`\pi _h`$ for $`40`$ values of its parameter $`r`$. The agreement will be seen to be excellent, that is, perfect within statistical errors. The first section covers the theoretical prediction, the second the details of the simulation and the comparison of the data with the prediction.
## 2 A theoretical prediction based <br>on conformal field theory
### 2.1 Cardy’s prediction for percolation
Cardy’s prediction for $`\pi _h`$ for two-dimensional percolation proceeds in two steps. He first identifies the probability $`\pi _h`$ with the difference of two partition functions with boundary conditions for the $`1`$-state Potts model. He then uses the conformal field theory at $`c=0`$ to obtain an analytic expression for this difference. Here is a (very rapid) presentation of these two steps.
The partition function $`Z(q)`$ of the $`q`$-state Potts model on a finite rectangular domain is the sum over all configurations $`\sigma `$ of $`e^{\beta H(\sigma )}`$ where $`H(\sigma )=J_{x,y}(1\delta _{\sigma (x),\sigma (y)})`$. The sum in $`H(\sigma )`$ runs over immediate neighbor pairs $`x,y`$. This can be rewritten as $`Z(q)=_Rp^{B(R)}(1p)^{BB(R)}q^{N_c(R)}`$ where $`p=1e^{\beta J}`$. The sum is over all subsets $`R`$ of the set of edges of the lattice in the rectangular domain. The integer $`B`$ counts the edges in the lattice, $`B(R)`$ those in the subset $`R`$ and $`N_c(R)`$ the clusters in $`R`$. If $`q=1`$ (the value for percolation), this sum is 1 as desired. Let $`Z_{\alpha \beta }`$ be the partition function of the Potts model for configurations whose spins on the left side of the rectangle are in the state $`\alpha \{1,2,\mathrm{},q\}`$, those on the right in the state $`\beta `$, and the others free. Cardy’s first crucial observation is that $`\pi _h=(Z_{\alpha \alpha }Z_{\alpha \beta })|_{q=1}`$ where $`\alpha \beta `$. (The difference, done for a “generic” $`q`$, contains precisely the configurations that have a cluster intersecting the left and right sides.) The problem of calculating $`\pi _h`$ is therefore transformed into that of calculating partition functions.
The possibility of calculating partition functions on finite domains with given boundary conditions also originates from works by Cardy (see for example ). In the case of $`Z_{\alpha \beta }`$, for example, the four sides of the rectangle are submitted respectively to the boundary conditions $`\alpha `$, free, $`\beta `$ and free. Cardy argues that such a partition function is proportional to the $`4`$-point correlation function $`\varphi (z_1)\varphi (z_2)\varphi (z_3)\varphi (z_4)`$ in the conformal field theory associated to percolation whose central charge is $`c=0`$. The $`z_i`$ are the vertices of the rectangle in the complex plane and $`\varphi `$ is the field that changes the boundary conditions in this theory. For percolation, this field is identified to $`\varphi _{1,2}`$ and it has conformal weight $`h=0`$, a necessary condition for $`\varphi (z_1)\varphi (z_2)\varphi (z_3)\varphi (z_4)`$ to be scale invariant. (The indices on $`\varphi _{1,2}`$ refer to the labels of Kac table. See .) The rules to find correlation functions in a conformal field theory are well-known and with this identification between partition functions with boundary conditions and $`4`$-point correlation functions, the problem of calculating $`\pi _h`$ amounts to solving an ordinary differential equation.
Two obstacles appear in applying these ideas to the Ising model. First the Ising model is the $`q=2`$ Potts model and the difference $`(Z_{\alpha \alpha }Z_{\alpha \beta })|_{q=2}`$ cannot be interpreted as the crossing probability since the factor $`q^{N_c(R)}`$ in the sum gives different weights to the various configurations that have a crossing from left to right. Second although the operator for the Ising model that changes the boundary state from free to a given state $`\alpha `$ is still $`\varphi _{1,2}`$, as in percolation, its conformal weight $`h_{1,2}`$ is now $`\frac{1}{16}`$ and its $`4`$-point correlation is not anymore invariant under a conformal mapping $`zw`$ but picks up the usual jacobian factors: $`\varphi (w_1)\varphi (w_2)\varphi (w_3)\varphi (w_4)=(_i|w^{}(z_i)|^{h_{1,2}})\varphi (z_1)\varphi (z_2)\varphi (z_3)\varphi (z_4)`$. These prefactors (and those additional coming from the presence of vertices along the boundary) seem to contradict the (strict) conformal invariance observed by simulation in .
In and the probability $`\pi _{hv}`$ of having simultaneous horizontal and vertical crossings was also obtained numerically. Watts was able to extend Cardy’s argument to obtain a prediction that fits extremely well the data. His work is of particular interest to us as he is able to write $`\pi _{hv}`$ again as a partition function but with more general boundary conditions. Whether a conformal boundary operator accomplishes the change between these more complex boundary conditions is not clear. Nonetheless the expression of $`\pi _{hv}`$ as a partition function allows us to expect, Watts argues, that this probability is given by some $`4`$-point correlation function. He seeks it in the $`h=0`$ sector.
### 2.2 The differential equation for $`\pi _h`$
Cardy’s argument for percolation cannot be extended to the Ising model. One can still hope to relate crossing probabilities like $`\pi _h`$ to $`4`$-point correlation functions as Watts did for $`\pi _{hv}`$ of percolation. If such a relationship exists, the choice can be narrowed to $`4`$-point functions of the identity family as these are the only ones in the $`c=\frac{1}{2}`$ conformal field theory that are invariant under conformal map $`zw`$ like $`\pi _h(r)`$. An obvious objection will be that the $`4`$-point function of the primary field in the identity family is identically $`1`$ (or a constant). For the calculation at hand it might well be that this primary field must be interpreted as one whose correlation functions satisfy only one of the differential equations corresponding to the two leading singular vectors. This milder requirement does not force the function to be a constant as will be seen immediately.
The Verma module $`V_{(c=\frac{1}{2},h=0)}`$ of the Virasoro algebra has a maximal proper submodule $`M`$ generated by two singular vectors, one at level $`1`$, the other at level $`6`$. The first of these is $`L_1|0`$ and the corresponding differential equation implies that the $`4`$-point function is a constant. We shall drop this requirement. The other singular vector is
$`(L_{1}^{}{}_{}{}^{6}10L_{1}^{}{}_{}{}^{4}L_2+{\displaystyle \frac{43}{3}}L_{1}^{}{}_{}{}^{2}L_{2}^{}{}_{}{}^{2}{\displaystyle \frac{100}{27}}L_{2}^{}{}_{}{}^{3}+{\displaystyle \frac{97}{2}}L_{1}^{}{}_{}{}^{3}L_3`$
$`{\displaystyle \frac{337}{6}}L_1L_2L_3+{\displaystyle \frac{3185}{144}}L_{3}^{}{}_{}{}^{2}{\displaystyle \frac{381}{2}}L_{1}^{}{}_{}{}^{2}L_4+{\displaystyle \frac{1265}{18}}L_2L_4`$
$`+{\displaystyle \frac{19309}{36}}L_1L_5{\displaystyle \frac{9005}{12}}L_6)|c=\frac{1}{2},h=0.`$
If $`f(z)`$ is the $`4`$-point function with $`z=(z_1z_2)(z_3z_4)/(z_1z_3)(z_2z_4)`$, the differential equation is
$$\begin{array}{cc}\hfill \frac{1}{72}\left(12z\right)\left(686z^2\left(1z\right)^2+73z\left(1z\right)+25\right)\frac{d}{dz}f(z)& \\ \hfill +\frac{1}{144}z\left(1z\right)\left(25141z^2\left(1z\right)^22986z\left(1z\right)171\right)\frac{d^2}{dz^2}f(z)& \\ \hfill +\frac{1}{27}z^2\left(1z\right)^2\left(12z\right)\left(2083595z\left(1z\right)\right)\frac{d^3}{dz^3}f(z)& \\ \hfill +\frac{1}{6}z^3\left(1z\right)^3\left(137737z\left(1z\right)\right)\frac{d^4}{dz^4}f(z)& \\ \hfill +10\left(12z\right)z^4\left(1z\right)^4\frac{d^5}{dz^5}f(z)& \\ \hfill +z^5\left(1z\right)^5\frac{d^6}{dz^6}f(z)& =0.\hfill \end{array}$$
(1)
This differential equation has three (regular) singular points at $`0,1`$ and $`\mathrm{}`$. It is invariant under any permutation of these three points. (Invariance under $`z1z`$ is clear: invariance under $`z1/z`$ requires some work.) The exponents at any of these points are $`0`$, $`\frac{1}{6}`$ twice degenerate, $`\frac{1}{2}`$, $`\frac{5}{3}`$ and $`\frac{5}{2}`$. The monodromy matrices around the three singular points are similar due to the symmetry of the equation but they cannot be diagonalized simultaneously.
The cross-ratio $`z=(z_1z_2)(z_3z_4)/(z_1z_3)(z_2z_4)`$ is related to the aspect ratio $`r=\sqrt{3}V/2H`$ of the rectangle. If the four points $`z_1,z_2,z_3`$ and $`z_4`$ are chosen along the real axis at $`\frac{1}{k},1,1,\frac{1}{k}`$, then
$$k=\frac{1\sqrt{z}}{1+\sqrt{z}}.$$
A Schwarz-Christoffel transformation can be used to map the upper plane onto a rectangle with the images of the $`z_i`$ at the vertices. The aspect ratio is then given as
$$r=\frac{K(1k^2)}{2K(k^2)}$$
where $`K`$ is the complete elliptic integral of the first kind. The very short but wide rectangles ($`r=\sqrt{3}V/2H0^+`$) corresponds to $`z0^+`$, the tall and narrow ($`r+\mathrm{}`$) to $`z1^{}`$ and the square to $`r=1,z=\frac{1}{2}`$. The function $`r(z)`$ has the property $`r(z)=r(\frac{1}{z})`$ and the symmetry $`z\frac{1}{z}`$ of the differential equation is thus welcome. If $`\pi _h`$ does not depend on the relative angle between the rectangle and the lattice, that is if $`\pi _h`$ is a rotational invariant, then the duality relation implies $`\pi _h(1)=\frac{1}{2}`$ and it can be put in the form $`(\frac{1}{2}\pi _h(r))=(\frac{1}{2}\pi _h(\frac{1}{r}))`$. Fortunately the function $`r(z)`$ is such that $`r(z)=1/r(1z)`$ and the duality simply states that the function $`f(z)=\frac{1}{2}\pi _h(r(z))`$ is odd with respect to the axis $`z=\frac{1}{2}`$. An odd subspace of the solution space of the differential equation (odd with respect with $`z=\frac{1}{2}`$) exists due to the symmetry $`z1z`$ and it is of dimension $`3`$.
We explored several paths to cast the solutions of eq. (1) into analytically tractable forms. One of them was to write the lhs of eq. (1) as $`_{1i6}(z(1z))^{a_i}\frac{d}{dz}`$ like Watts did. But there is no real solutions for the $`a_i`$’s in the present case. The most natural path however is the screening operator method (, see also ). The pertinent field is $`\varphi _{2,3}`$ (with Kac’s labels) and the 4-point correlation calls for three contour integrals (a charge $`Q_+Q_{}^2`$ must be added at infinity to assure neutrality). Integral representations of six linearly independent solutions can be obtained in this straightforward (but probably tedious) way. The main problem is therefore whether there is sufficient physical information on $`f`$ to fix the linear combination. As argued above the odd parity of $`f`$ reduces the space of solutions to a three-dimensional subspace. The condition $`f(0)=\frac{1}{2}`$ (i.e. $`\pi _h(0)=1\pi _h(1)=0`$) is one further linear constraint. The function $`f(z)`$ is monotone increasing but this condition will restrict $`f`$ to an open set of the $`3`$-d subspace rather than decrease the dimension. The numerical data presented below indicate that $`\pi _h(z)z^{\frac{1}{6}}`$ as $`z0^+`$. This is rather striking in view of the two-fold degeneracy of the exponent $`\frac{1}{6}`$. Two solutions associated to this exponent can be chosen to behave as $`z^{\frac{1}{6}}`$ and $`z^{\frac{1}{6}}\mathrm{log}z`$ for small $`z`$ and the latter, if present in $`f`$, should dominate the former whenever $`z`$ is close to $`0`$. The fact that it is not seen in the simulation could be interpreted physically as a manifestation of the power law behavior of critical correlation functions at short distance. Imposing that the behavior in $`z^{\frac{1}{6}}\mathrm{log}z`$ be absent of $`f`$ would add one linear constraint. With all these constraints we would still be left with a one-dimensional subspace in the space of odd solutions. We have not found any further constraints to fix completely the function $`f`$. This is why we resorted to a numerical fit (see Paragraph 3.2).
As we were trying to solve analytically the differential equation, Marc-André Lewis suggested to us to look for a solution of the form $`ez^d{}_{2}{}^{}F_{1}^{}(a,b,c;z)`$, with $`a,b,c,d`$ and $`e`$ constants, that would be odd around $`z=\frac{1}{2}`$ and reproduce the asymptotic behavior of the data. Cardy’s prediction for percolation is of this form and the suggestion is natural in this sense. Such a function exists but it does not satisfy the differential equation. However it follows so closely the data of (the worst gap is $`1`$%) that we decided to improve these measurements to answer the question: which of the hypergeometric function or of the solution of the differential equation, if any, describes the data.
## 3 Improved measurements of $`\pi _h`$
### 3.1 Finite size effects and power law behavior
The probabilities $`\pi _h,\pi _v`$ and $`\pi _{hv}`$ at $`81`$ values of $`r`$ ($`[0.136,7.351]`$) were measured in for the three regular lattices on rectangles containing around $`40000`$ sites. For these sizes, departure from the duality relation is small but still noticeable for the square and the hexagonal lattices. No attempt was made there to use various sizes in order to reduce finite-size effects. For the new runs to be presented here we chose to concentrate on the triangular lattice and use various sizes to approximate the function $`\pi _h`$ (and the other two, $`\pi _v`$ and $`\pi _{hv}`$) in the limit when the mesh goes to zero.
A power law for the finite-size behavior of critical data is an accepted hypothesis and our measurements rest upon it. It states that, for sufficiently large size,
$$|\pi _h(V,H)\pi _h(r)|aV^\beta $$
with $`\beta `$ a negative constant and $`r=\sqrt{3}V/2H`$. We shall use several linear sizes of the form $`V=2^iV_0`$ and $`H=2^iH_0`$. Writing $`\pi _h(i)`$ for $`\pi _h(2^iV_0,2^iH_0)`$ and supposing that $`\pi _h(i)`$ is decreasing, we can write $`\pi _h(i)\pi _h(i+1)aV_0^\beta 2^{i\beta }(12^\beta )`$ and therefore express $`\pi _h(r)`$ as $`\pi _h(i)aV_0^\beta 2^{i\beta }`$. To determine the constants $`a`$ and $`\beta `$ requires at least $`3`$ rectangle sizes as only the differences $`(\pi _h(i)\pi _h(i+1))`$ can be used.
What are the right rectangle sizes and what is the required precision on each $`\widehat{\pi }_h(i)`$? For the two extreme rectangles that we are planning to measure, $`r_{\text{wide}}=\frac{\sqrt{3}}{2}\frac{V_0}{H_0}=\frac{\sqrt{3}}{2}\frac{4}{24}0.1443`$ and $`r_{\text{tall}}=\frac{\sqrt{3}}{2}\frac{32}{4}6.928`$, and for a rectangle close to a square $`r_{\text{sq}}=\frac{\sqrt{3}}{2}\frac{16}{14}0.9897`$, we obtained $`\pi _h(i)`$, $`i=0,1,2,3,4`$. Each increment corresponds to an increase by a factor of 2 of the linear size. For example $`\pi _h(0)`$ was measured on a rectangle of $`4\times 24`$ sites and $`\pi _h(4)`$ on $`64\times 384`$ sites for the wide rectangle. The results appear in Table 1. The samples were large, at least $`250\times 10^6`$ configurations. The digits after the vertical bar gives the statistical error on the digits just before; for example, the first element in the table ($`0.029269|21`$) means that $`\widehat{\pi }_h(0)`$ is $`0.029269`$ with the 95%-confidence interval being $`[0.029248,0.029290]`$. The differences between $`\widehat{\pi }(3)`$ and $`\widehat{\pi }(4)`$ are however small. In fact the monotonicity of $`\widehat{\pi }_h(i)`$ for $`r_{\text{tall}}=6.928`$ is broken for $`i=4`$ even though the error bars do allow for the power law to hold. Larger samples would definitely be required for the large lattices. Fortunately the precision on the measurements for $`i=0,1,2`$ and the fact that the power law seems to hold for very small sizes (4 sites in one direction!) allow for good estimates of $`\pi _h(r)`$ without these larger lattices. Table 2 shows estimates of $`\pi _h`$ for $`r_{\text{wide}}`$, $`r_{\text{tall}}`$ and $`r_{\text{sq}}`$ using the power law hypothesis and a subset of the measurements of Table 1. The notation $`i`$$`j`$ means that $`\widehat{\pi }_h(i),\widehat{\pi }_h(i+1),\mathrm{},\widehat{\pi }_h(j)`$ were used to obtain $`\widehat{\pi }_h(r)`$. Using only the three smallest lattices, the three largest or the five ones lead to estimates $`\widehat{\pi }_h(r)`$ that differ by less than $`4`$ units on the fourth significant digits. We therefore decided to use only three sizes for each $`r`$ considered and choose the pairs $`(V_0,H_0)`$ in such a way that $`V_0,H_04`$ and that $`V_0H_096`$. All samples were larger or equal to $`10^8`$. Tables 3 gives the results for $`\pi _h`$, $`\pi _v`$ and $`\pi _{hv}`$ at $`40`$ values of $`r[0.1443,6.928]`$.
There is always, on a finite lattice, the problem of determining $`r`$ from the numbers $`V`$ and $`H`$. The two simplest choices are the aspect ratios of the smallest or the largest rectangles that include the sites considered and only those. Even though it is not a natural choice, the aspect ratio of the tallest and narrowest rectangle would be another convention. The method we have used to determine $`\pi _h(r)`$ overcomes this imprecision due to convention. For any convention the aspect ratio for a rectangular subset of the triangular lattice will be $`\frac{\sqrt{3}}{2}(V+\mathrm{\Delta }_V)/(H+\mathrm{\Delta }_H)`$ with $`\mathrm{\Delta }_V`$ and $`\mathrm{\Delta }_H`$ dependent on the convention but independent of $`V`$ and $`H`$. The ratio $`r`$ at which $`\pi _h`$ is measured is therefore $`lim_i\mathrm{}\frac{\sqrt{3}}{2}(2^iV_0+\mathrm{\Delta }_V)/(2^iH_0+\mathrm{\Delta }_H)=\frac{\sqrt{3}}{2}V_0/H_0`$, independent of $`\mathrm{\Delta }_V`$ and $`\mathrm{\Delta }_H`$, that is independent of the convention. This is another advantage of using several lattices for a given $`r`$.
Even though $`\pi _h`$ (and $`\pi _v`$) is invariant under rotation, finite-size effects are not. The differences between $`\widehat{\pi }_h(0)=0.02927`$, $`\widehat{\pi }_h(1)=0.02558`$ and $`\widehat{\pi }_h(2)=0.02381`$ for $`r=0.1443`$ are much bigger than those between $`\widehat{\pi }_v(0)=0.02174`$, $`\widehat{\pi }_v(1)=0.02205`$ and $`\widehat{\pi }_v(2)=0.02214`$ for $`1/r=1/6.928=0.1443`$. The statistical errors on $`\widehat{\pi }_h(0.1443)`$ and $`\widehat{\pi }_v(6.928)`$ are therefore different even if both numbers turn out to be very close ($`0.02215`$ and $`0.02218`$). In the worst cases the $`95`$%-confidence interval amounts to less than $`2`$ units on the third significant digit (e.g. $`\widehat{\pi }_h(0.1443)=0.02215\pm 0.00015`$). At the center of the range of $`r`$ the error on $`\widehat{\pi }_h(r)`$ decreases to $`4`$ units on the fourth digit and it is even smaller for large $`r`$. Over the whole range it is smaller than $`4\times 10^4`$ for both $`\pi _h`$ and $`\pi _v`$.
The improvement upon previous measurements found in can be checked easily. Among the 40 values of $`r`$ used there are nine pairs $`((V_a,H_a),(V_b,H_b))`$ such that
$$r_a=\frac{\sqrt{3}}{2}\frac{V_a}{H_a}=\left(\frac{\sqrt{3}}{2}\frac{V_b}{H_b}\right)^1=\frac{1}{r_b}.$$
The pairs $`(r_a,1/r_a)`$ are those corresponding to the following lines of Table 3: $`(1,40),(5,36),(7,34),(8,33),(10,31),(11,30),(12,29),(18,23),(20,21)`$. The measurements for these pairs should satisfy $`\widehat{\pi }_h(r_a)=\widehat{\pi }_v(r_b)`$, $`\widehat{\pi }_v(r_a)=\widehat{\pi }_h(r_b)`$ and $`\widehat{\pi }_{hv}(r_a)=\widehat{\pi }_{hv}(r_b)`$ within statistical errors. This turns out to be the case. (See Figure 1.) There are in total 27 independent comparisons. Their relative errors is always less than $`2\times 10^3`$. The largest occur when the quantities $`\widehat{\pi }`$ being compared are themselves very small, like $`\widehat{\pi }_h(\frac{\sqrt{3}}{2}\frac{1}{6})=0.02215`$ and $`\widehat{\pi }_v(\frac{\sqrt{3}}{2}8)=0.02218`$. In all cases the absolute value of these differences are less than $`5\times 10^4`$. Not only are these variations small, they are of both signs. This fact is a further indication that the power law hypothesis provides a very good approximation. Suppose indeed that, at the sizes used, correction terms are required: $`|\pi _h(i)\pi _h(r)|aV^\beta (1+\frac{b}{V}+\mathrm{})`$. These new terms would lead to a systematic error that is not seen here.
### 3.2 A prediction for $`\pi _h`$ for the critical Ising model
The first easy comparison between the theory developed in the first Section and the new data lies in the asymptotic behavior of $`\pi _h`$ as $`r`$ approaches $`0`$ and $`+\mathrm{}`$. If $`\pi _h`$ is a solution of the differential equation (1), then $`\mathrm{log}\pi _h(r)\lambda \pi /r`$ as $`r0`$ and $`\mathrm{log}(1\pi _h(r))\lambda \pi r`$ as $`r+\mathrm{}`$ with $`\lambda `$ one of the exponents. Using the ten extreme values of $`r`$ on each side of the measured interval, we obtain for $`\pi _h`$
$`\mathrm{log}\widehat{\pi }_h(r)`$ $`\underset{r0}{}0.16648\pi {\displaystyle \frac{1}{r}}`$
$`\mathrm{log}(1\widehat{\pi }_h(r))`$ $`\underset{r\mathrm{}}{}0.16657\pi r.`$
The slopes obtained using the data for $`\pi _v`$ are $`0.16647\pi `$ and $`0.16654\pi `$. These numbers are very close to the exponent $`\frac{1}{6}`$, the smallest non-vanishing exponent of the differential equation (1). This is remarkable! Clearly $`\lambda `$ is to be interpreted as a critical exponent of the Ising model. But none of the usual exponents of the Ising model contains the prime number $`3`$ in their denominators and scaling laws involve only products and integral linear combinations of these exponents. If this new critical exponent can be deduced from the usual ones, it will not be by traditional scaling laws.
The second test is to obtain a solution of the differential equation (1) that describes the data. Let $`f_i,i=0,1,\mathrm{},5`$, be a basis of solutions for equation (1) defined by their behavior at $`z=\frac{1}{2}`$: $`f_i^{(j)}(z=\frac{1}{2})=\delta _{ij}`$. Since $`f(z)=\frac{1}{2}\pi _h(r(z))`$ is odd with respect to $`z=\frac{1}{2}`$, it lies in the subspace of functions of the form $`af_1+bf_3+cf_5`$. We determine the constants $`a,b`$ and $`c`$ by requiring that $`L=(_i(\frac{1}{2}\widehat{\pi }_h(r_i)f(z(r_i)))^2)^{\frac{1}{2}}`$ be minimum. The sum is over the 40 data. The three solutions $`f_1,f_3`$ and $`f_5`$ were obtained numerically. Both Matlab and Mathematica give similar fits. These softwares have internal parameters controlling the required accuracy of the integration. These parameters can be pushed to a point where stronger requirement does not lead to any significant improvement on the minimum of $`L`$. Figure 2 has been drawn using the values $`a,b`$ and $`c`$ obtained for control parameters beyond this point. The largest among the differences $`|\frac{1}{2}\widehat{\pi }_h(r_i)f(z(r_i))|,i=1,\mathrm{},40`$, is $`3.6\times 10^4`$, smaller than the statistical error, and the standard deviation is $`1.5\times 10^4`$. Similar results are obtained for $`\pi _v`$. The agreement is therefore excellent. As a comparison it is instructive to redo Cardy’s calculation testing his prediction for percolation. Since the publication of , better data were obtained for percolation by sites on a square lattice for $`81`$ rectangles with at least $`10^6`$ sites . The samples contained over $`10^6`$ configurations. For these, the statistical errors are approximately $`10^4`$ at the extremities of the interval ($`r[0.142,7.067]`$) and $`10^3`$ in the middle. (At the extremities these data are therefore more precise than the present ones for the Ising model and at the center they are less precise.) The largest departure from duality is $`4\times 10^4`$, almost exactly what is seen in Figure 1 for the present data. Comparing his prediction $`\pi _h^{\text{perco}}(r(z))=3\mathrm{\Gamma }(\frac{2}{3})z^{\frac{1}{3}}{}_{2}{}^{}F_{1}^{}(\frac{1}{3},\frac{2}{3},\frac{4}{3},z)/\mathrm{\Gamma }(\frac{1}{3})^2`$ with the data leads to $`\mathrm{max}_{1i81}|\pi _h^{\text{perco}}(r(z_i))\widehat{\pi }_h(r_i)|=7.8\times 10^4`$ and to a standard deviation of $`4.2\times 10^5`$. These results are similar to those just reported.
We mentioned earlier the possibility of describing the data with a hypergeometric function. The function
$$g(z)=\frac{6\mathrm{\Gamma }(\frac{1}{3})}{\mathrm{\Gamma }(\frac{1}{6})^2}z^{\frac{1}{6}}{}_{2}{}^{}F_{1}^{}(\frac{1}{6},\frac{5}{6},\frac{7}{6};z)$$
is odd around $`z=\frac{1}{2}`$, when shifted by $`\frac{1}{2}`$, and behaves as $`z^{\frac{1}{6}}`$ and $`(1z)^{\frac{1}{6}}`$ when $`z0^+`$ and $`z1^{}`$. Several data are now more than $`2\times 10^3`$ apart from the corresponding values of $`g`$, a gap barely visible on a figure, but clearly out of any reasonable confidence interval. We may accept the solution of the differential equation as an analytic prediction for $`\pi _h`$ but we must reject the function $`g(z)`$.
Measurements were also made of the probability $`\pi _{hv}`$. As its behavior at $`z=0`$ and $`z=1`$ is also described by the exponent $`\frac{1}{6}`$, one may hope that the even subspace (around $`z=\frac{1}{2}`$) of the differential equation may contain a solution matching the data. This is not the case. The best fit in this subspace lies up to $`6\times 10^3`$ away from the data, an unacceptable gap. (This disagreement has an advantage. It shows that the success of the fit for $`\pi _h`$ is not a consequence of the large freedom that a three-parameter fit gives.) Watts used the third singular vector to describe $`\pi _{hv}`$ for percolation. In the Verma module $`V_{c=0,h=0}`$, this vector is at level five, leading to a differential equation of order $`5`$. However the third one in $`V_{c=\frac{1}{2},h=0}`$ is at level 11 and the associated operator $`P(z,\frac{d}{dz})`$ is a polynomial of order $`11`$ in $`\frac{d}{dz}`$. It can be cast as
$`P(z,{\displaystyle \frac{d}{dz}})=`$ $`{\displaystyle \underset{1i11,i\text{ odd}}{}}((z(1z))^{i1}p_i(z(1z)){\displaystyle \frac{d^i}{dz^i}}`$
$`+{\displaystyle \underset{2i10,i\text{ even}}{}}(12z)((z(1z))^{i1}p_i(z(1z)){\displaystyle \frac{d^i}{dz^i}}`$
with
$`p_1(u)`$ $`=\frac{1}{81}u\left(48420465u+120702u^2+134456u^3+326536u^4\right)`$
$`p_2(u)`$ $`=\frac{1}{81}\left(48465267u+555942u^2+2442422u^3+9038782u^4\right)`$
$`p_3(u)`$ $`=\frac{1}{162}\left(443181002950u5132553u^212317152u^3+245463307u^4\right)`$
$`p_4(u)`$ $`=\frac{1}{486}\left(4128398111249u42237350u^2+749236363u^3\right)`$
$`p_5(u)`$ $`=\frac{1}{1944}\left(3492203+28198986u1347384726u^2+5106251212u^3\right)`$
$`p_6(u)`$ $`=\frac{7}{1944}\left(36833927634444u+153070553u^2\right)`$
$`p_7(u)`$ $`=\frac{11}{1296}\left(5545519066926u+29200567u^2\right)`$
$`p_8(u)`$ $`=\frac{11}{54}\left(11842+75835u\right)`$
$`p_9(u)`$ $`=\frac{11}{18}\left(757+3457u\right)`$
$`p_{10}(u)`$ $`=\frac{110}{3}`$
$`p_{11}(u)`$ $`=1.`$
The differential equation $`P(z,\frac{d}{dz})f(z)=0`$ is again symmetric under any permutation of the three regular singular points $`0,1`$ and $`\mathrm{}`$. The exponents (with their degeneracies) are $`0(2),\frac{1}{6}(2),\frac{1}{2}(1),1(1),\frac{5}{3}(2),\frac{5}{2}(1),\frac{14}{3}(1)`$ and $`6(1)`$. The even subspace of the differential equation is of dimension 6 and the fit will therefore contain 6 parameters. Using the data of the last column of Table 3 and numerical integration of the differential equation, we obtain a best fit (in the same sense as the one used for $`\pi _h`$) that has a largest difference of $`2.7\times 10^4`$ and a standard deviation of $`1.3\times 10^4`$. These are excellent results well within the experimental windows. However we have tried to fit similarly the function $`h(z)=\kappa z^{\frac{1}{6}}(1z)^{\frac{1}{6}}`$ with the constant $`\kappa `$ chosen such that $`h(z=\frac{1}{2})=\widehat{\pi }_{hv}(z=\frac{1}{2})`$. While this function is not a solution, the 6-dimensional subspace of (numerical) even solutions contains a solution $`\stackrel{~}{h}`$ that approaches it extremely well (namely $`\mathrm{max}_{z[0,1]}|\stackrel{~}{h}(z)h(z)|10^6`$). Consequently the large dimension of the even subspace allows for functions that are not solution to be fitted within statistical errors and the fit for $`\pi _{hv}`$ is much less convincing than the one above for $`\pi _h`$ or than Watts’ prediction.
## Acknowledments
It is a pleasure to thank R.P. Langlands, M.-A. Lewis and G. Watts for helpful discussions, A. Bourlioux for introducing us to the control of errors in integrating differential equations and H. Pinson and Ph. Zaugg for a careful reading of the manuscript.
E. L. gratefully acknowledges a fellowship from the NSERC Canada Scholarships Program and Y. S.-A. support from NSERC (Canada) and FCAR (Québec).
## List of captions
Table 1: $`\widehat{\pi }_h`$ for three aspect ratios $`r`$ and five sizes.
Table 2: Estimates $`\widehat{\pi }_h`$ using subsets of available data.
Table 3: The measurements $`\widehat{\pi }_h`$, $`\widehat{\pi }_v`$ and $`\widehat{\pi }_{hv}`$.
Figure 1: Differences $`\widehat{\pi }_h(r_a)\widehat{\pi }_v(r_b)`$ and $`\widehat{\pi }_{hv}(r_a)\widehat{\pi }_{hv}(r_b)`$ for several pairs $`(r_a,r_b=r_a^1)`$.
Figure 2: The prediction $`\mathrm{log}\pi _h/(1\pi _h)`$ as a function of $`\mathrm{log}r`$ together with the 40 measurements $`\widehat{\pi }_h(r_i)`$.
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# Quantum Spin Dynamics (QSD) : VII. Symplectic Structures and Continuum Lattice Formulations of Gauge Field Theories
## 1 Introduction
Quantum General Relativity (QGR) has matured over the past decade to a mathematically well-defined theory of quantum gravity. In contrast to string theory, by definition, GQR is a manifestly background independent, diffeomorphism invariant and non-perturbative theory. The obvious advantage is that one will never have to postulate the existence of a non-perturbative extension of the theory, which in string theory has been called the still unknown M(ystery)-Theory.
The disadvantage of a non-perturbative and background independent formulation is, of course, that one is faced with new and interesting mathematical problems so that one cannot just go ahead and “start calculating scattering amplitudes”: As there is no background around which one could perturb, rather the full metric is fluctuating, one is not doing quantum field theory on a spacetime but only on a differential manifold. Once there is no (Minkowski) metric at our disposal, one loses familiar notions such as causality structure, locality, Poincaré group and so forth, in other words, the theory is not a theory to which the Wightman axioms apply. Therefore, one must build an entirely new mathematical apparatus to treat the resulting quantum field theory which is drastically different from the Fock space picture to which particle physicists are used to.
As a consequence, the mathematical formulation of the theory was the main focus of research in the field over the past decade. The main achievements to date are the following (more or less in chronological order) :
* Kinematical Framework
The starting point was the introduction of new field variables for the gravitational field which are better suited to a background independent formulation of the quantum theory than the ones employed until that time. In its original version these variables were complex valued, however, currently their real valued version, considered first in for classical Euclidean gravity and later in for classical Lorentzian gravity, is preferred because to date it seems that it is only with these variables that one can rigorously define the dynamics of Euclidean or Lorentzian quantum gravity .
These variables are coordinates for the infinite dimensional phase space of an $`SU\left(2\right)`$ gauge theory subject to further constraints besides the Gauss law, that is, a connection and a canonically conjugate electric field. As such, it is very natural to introduce smeared functions of these variables, specifically Wilson loop and electric flux functions. (Notice that one does not need a metric to define these functions, that is, they are background independent). This had been done for ordinary gauge fields already before in and was then reconsidered for gravity (see e.g. ).
The next step was the choice of a representation of the canonical commutation relations between the electric and magnetic degrees of freedom. This involves the choice of a suitable space of distributional connections and a faithful measure thereon which, as one can show , is $`\sigma `$-additive. The proof that the resulting Hilbert space indeed solves the adjointness relations induced by the reality structure of the classical theory as well as the canonical commutation relations induced by the symplectic structure of the classical theory can be found in . Independently, a second representation, called the loop representation, of the canonical commutation relations had been advocated (see e.g. and especially and references therein) but both representations were shown to be unitarily equivalent in (see also for a different method of proof).
This is then the first major achievement : The theory is based on a rigorously defined kinematical framework.
* Geometrical Operators
The second major achievement concerns the spectra of positive semi-definite, self-adjoint geometrical operators measuring lengths , areas and volumes of curves, surfaces and regions in spacetime. These spectra are pure point (discete) and imply a discrete Planck scale structure. It should be pointed out that the discreteness is, in contrast to other approaches to quantum gravity, not put in by hand but it is a prediction !
* Regularization- and Renormalization Techniques
The third major achievement is that there is a new regularization and renormalization technique for diffeomorphism covariant, density-one-valued operators at our disposal which was successfully tested in model theories . This technique can be applied, in particular, to the standard model coupled to gravity and to the Poincaré generators at spatial infinity . In particular, it works for Lorentzian gravity while all earlier proposals could at best work in the Euclidean context only (see, e.g. and references therein). The algebra of important operators of the resulting quantum field theories was shown to be consistent . Most surprisingly, these operators are UV and IR finite ! Notice that this result, at least as far as these operators are concerned, is stronger than the believed but unproved finiteness of scattering amplitudes order by order in perturbation theory of the five critical string theories, figuratively speaking, we claim that our perturbation series converges. The absence of the divergences that usually plague interacting quantum fields propagating on a Minkowski background can be understood intuitively from the diffeomorphism invariance of the theory : “short and long distances are gauge equivalent”. We will elaborate more on this point in future publications. The classical limit of the above mentioned operators will be studied in our companion paper .
* Spin Foam Models
After the construction of the densely defined Hamiltonian constraint operator of , a formal, Euclidean functional integral was constructed in and gave rise to the so-called spin foam models (a spin foam is a history of a graph with faces as the history of edges) . Spin foam models are in close connection with causal spin-network evolutions , state sum models and topological quantum field theory, in particular BF theory . To date most results are at a formal level and for the Euclidean version of the theory only but the programme is exciting since it may restore manifest four-dimensional diffeomorphism invariance which in the Hamiltonian formulation is somewhat hidden.
* Finally, the fifth major achievement is the existence of a rigorous and satisfactory framework for the quantum statistical description of black holes which reproduces the Bekenstein-Hawking Entropy-Area relation and applies, in particular, to physical Schwarzschild black holes while stringy black holes so far are under control only for extremal charged black holes.
Summarizing, the work of the past decade has now culminated in a promising starting point for a quantum theory of the gravitational field plus matter and the stage is set to address physical questions. In particular, one would like to make contact with the language that particle physicists are more familiar with, that is, perturbation theory. In other words, one should be able to define something like gravitons and photons propagating on a fluctuating quantum spacetime. By this we mean the following :
Suppose we want to study the semi-classical limit of our quantum gravity theory, that is, a limit in which the gravitational field behaves almost classical. This does not mean that we want to treat gravity as a background field , rather we take all the quantum fluctuations into account but try to find a state with respect to which those fluctuations (around the Minkowski metric) are minimal. With respect to such a “background state” one can study relative excitations of the gravitational field (gravitons) or of matter fields (such as photons).
In order to do this we must first develop an appropriate semi-classical framework which we will do in . But even before doing this we must examine the following issue which seems not to have been sufficiently appreciated throughout the literature so far :
Namely, the quantum theory is based on certain configuration and conjugate momentum variables respectively, specifically holonomy – and electric flux variables. We stress that we use here non-standard flux variables not previously considered in the literature. These non-local functions on the classical phase space are, in particular, used to regularize more complicated composite operators such as the geometrical operators mentioned above. It is already quite remarkable that one can remove the regulator without encountering any UV divergencies ! However, in order to be convinced that this regulator-free operator really has the correct classical limit one has to check, for instance, that it has the correct expectation values with respect to semi-classical states. The question arises how such semi-classical states should look like. Now, since the final regulator-free operator is actually only densely defined (since it is usually unbounded) one has to employ states which are semi-classical and simultaneously belong to a dense subspace of the Hilbert space. The states which belong to the domain of definition of the operator are labelled by graphs $`\gamma `$. Given such a graph $`\gamma `$ one can define unambiguously holonomies along its edges as the basic configuration operators labelled by $`\gamma `$, however, the associated (conjugate) momentum operators are largely ambiguous in the sense that any choice of surfaces which are mutually disjoint and are intersected by precisely one edge of $`\gamma `$ gives rise to completely identical Poisson brackets between the canonical variables.
This then leads to the following problem : Suppose we define a semi-classical state by requiring that the expectation value of the basic holonomy and electric flux operators associated with $`\gamma `$ take certain values and satisfy a minimal uncertainty property. Given a point in the classical phase space those values should be the values of the holonomy and flux functions evaluated at that point. However, this makes sense only when we specify the surfaces with respect to which we calculate the flux.
We are thus led to invent a new kind of generalized projective family labelled not only by graphs but also by so-called “dual” faces. In particular, we wish to do this already at the classical level by introducing a new kind of generalized projective family of symplectic manifolds. The idea behind all of this is that these symplectic manifolds enable us to discuss in a clean way the quantization procedure and its inverse, the process of taking the classical limit :
* Classical Regularization
Suppose we are given a function on the classical phase space $`(M,\mathrm{\Omega })`$, usually a function $`F\left(m\right)`$ of the connection and the electric field, $`m=(A,E)`$. Here $`M`$ denotes the set of connections and electric fields respectively (a differentiable manifold modelled on a Banach space, see below) and $`\mathrm{\Omega }`$ is a strong symplectic structure on $`M`$. As we cannot define $`\widehat{A},\widehat{E}`$ on our Hilbert space directly as operators, we must first find a substitute $`F_\gamma \left(m\right)`$ for $`F`$ which can be written entirely in terms of holonomy and flux variables associated with $`\gamma `$. These variables coordinatize a symplectic manifold $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$. We will say that $`F_\gamma \left(m\right)`$ is a substitute for $`F\left(m\right)`$ provided that 1) $`F_\gamma `$ converges to $`F`$ pointwise on $`M`$ as $`\gamma \mathrm{}`$ (the graph becomes infinitely fine, we will specify this limit below) and 2) that the Hamiltonian vector field of $`F_\gamma `$ with respect to $`\mathrm{\Omega }_\gamma `$ converges pointwise on $`M`$ to that of $`F`$ with respect to $`\mathrm{\Omega }`$.
* Regularized Operators
The classical phase spaces $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ turn out to be (in)finite direct products of (copies of) cotangent spaces over the gauge group $`G`$ equipped with a non-standard symplectic structure and allow for a bona fide quantization by usual geometrical quantization techniques. By substituting classical variables for operators defined on a subspace $`_\gamma `$ of the Hilbert space and Poisson brackets with respect to $`\mathrm{\Omega }_\gamma `$ by commutators we obtain an operator $`\widehat{F}_\gamma `$ unambiguously defined on $`_\gamma `$ up to factor ordering ambiguities. Thus, the phase spaces $`M_\gamma `$ are much better suited for the quantization of interesting functions $`F`$ on $`M`$ as they are automatically finite and we have always control that the quantization has the correct classical limit on $`M_\gamma `$. In other words, quantization and regularization can be neatly separated as individual processes.
* Unregularized Operator
It turns out that for a large class of functions $`F`$ including the ones of physical interest the family of operators $`\widehat{F}_\gamma `$ so obtained provides an operator $`\widehat{F}`$ consistently defined on a dense subspace of the whole Hilbert space in the sense that its restriction to $`_\gamma `$ coincides with $`\widehat{F}_\gamma `$. This will be our candidate for a well-defined continuum operator.
* Classical Limit
In order to study the classical limit of $`\widehat{F}`$ we introduce a generalized projective family of semi-classical states $`\psi _{\gamma ,m}^t_\gamma `$ labelled by the graph $`\gamma `$, a point in $`mM`$ and a classicality parameter $`t\mathrm{}`$. We say that $`\widehat{F}_\gamma `$ is a quantization of $`F_\gamma \left(m\right)`$ provided that $`lim_{t0}<\psi _{\gamma ,m}^t,\widehat{F}_\gamma \psi _{\gamma ,m}^t>=F_\gamma \left(m\right)`$ for each $`mM`$ and that $`\widehat{F}`$ is a quantization of $`F`$ provided that $`lim_{t0}[lim_\gamma \mathrm{}<\psi _{\gamma ,m}^t,\widehat{F}_\gamma \psi _{\gamma ,m}^t>]=F\left(m\right)`$ for each $`mM`$.
Theses four steps provide then a closed path of how to go from a classical phase space function to an operator and back. As we see, this procedure requires as a classical cornerstone the analysis of the phase spaces $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ which is the subject of the present paper. In particular, one must show that these symplectic manifolds contain a generalized projective sequence that can be identified with $`(M,\mathrm{\Omega })`$.
The outline of the paper is as follows :
In section two we recall a working collection of material from the kinematical framework of the theory.
In section three we derive from the symplectic manifold $`(M,\mathrm{\Omega })`$ for gauge theories with compact gauge groups (in any dimension and on any (globally hyperbolic) manifold) a generalized projective family of (in)finite dimensional symplectic manifolds $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ labelled by graphs $`\gamma `$ embedded in that manifold. We show that the generalized projective limit symplectic manifold of a certain generalized projective sequence agrees with the standard symplectic manifolds $`(M,\mathrm{\Omega })`$ for gauge theories (weighted Sobolev spaces). The purpose of doing this is that the generalized family of symplectic manifolds is much better suited to quantization than the standard gauge theory phase space as outlined above.
In section four we propose a substitute $`G_\gamma `$ for an important function $`G`$ on the phase space of any gauge theory, namely the Gauss constraint and show that $`G_\gamma `$ converges to $`G`$ pointwise on $`M`$ in the generalized projective limit.
In section five we derive the quantization of $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ and $`G_\gamma `$. We show that $`\widehat{G}_\gamma `$ is a consistently defined system of cylindrical projections of an operator $`\widehat{G}`$ whose constraint algebra closes without anomalies. Finally we sketch the last step of the above programme applied to $`\widehat{G}`$ concerning the classical limit. The proof that this step can be completed will be found in .
In section six we complement earlier results obtained by Ashtekar, Corichi and Zapata :
These authors considered certain classical functions $`F`$ on $`M`$ and quantized them using $`\mathrm{\Omega }`$ as a starting point. They obtained operators $`\widehat{F}`$ this way which do not commute on certain states $`f_\gamma _\gamma `$ although the classical functions $`F`$ have vanishing Poisson brackets (with respect to $`\mathrm{\Omega }`$) among each other. This seeming quantum “anomaly” was explained by pointing out that the connection and electric field of the theory are smeared with distributional rather than smooth test functions. If one uses a smearing with smooth functions then the “anomaly” vanishes, allowing the interpretation that the Poisson brackets of the unsmeared fields is non-vanishing, proportional to a distribution with support contained in a measurable subset of Lebesgue measure zero which is therefore detectable only when smearing with distributional smearing functions. This interpretation therefore removes the apparent contradiction. However, then one notices that this extended Poisson bracket does not close in an obvious way (the Jacobi identity is not obeyed in an obvious way). This was shown not to be an obstacle to quantization by recalling that it is not necessary to base the quantization on Poisson brackets but that one can instead base it on the Lie algebra of vector fields on $`M`$ which always obey the Jacobi identity and is always closed.
We show that the non-commutativity of these operators has a natural explanation from the point of view of the symplectic manifolds $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ :
1) We observe that we can find, for each of the above choices of $`\gamma `$, functions $`F_\gamma F`$ which can be considered as functions on $`M_\gamma `$ as well. Furthermore, the functions $`F_\gamma `$ do have non-vanishing Poisson brackets among each other, both with respect to $`\mathrm{\Omega }`$ and with respect to $`\mathrm{\Omega }_\gamma `$ (actually, their brackets with respect to $`\mathrm{\Omega }_\gamma `$ follow from those with resepct to $`\mathrm{\Omega }`$).
2) The quantization of these new functions is such that $`\widehat{F}_\gamma `$ and $`\widehat{F}`$ agree on $`_\gamma `$.
3) The commutator algebra of the $`\widehat{F}`$ on $`_\gamma `$ is precisely the one to be expected from the Poisson bracket structure of the $`F_\gamma `$.
In conclusion, the unexpected non-commutativity observed in can be related to a quantization ambiguity. If we insist on a Poisson bracket – comutator correspondence principle, however, then one cannot accept the $`F`$ as classical limit of $`\widehat{F}`$ but must instead consider the $`F_\gamma `$.
Finally, in an appendix we write the symplectic structure $`\mathrm{\Omega }_\gamma `$ for $`G=U\left(1\right),SU\left(2\right)`$ in the language of differential forms which could be useful for future research.
## 2 Preliminaries
In this section we will recall the main ingredients of the mathematical formulation of diffeomorphism invariant quantum field theories of connections with local degrees of freedom in any dimension and for any compact gauge group. See and references therein for more details.
Let $`G`$ be a compact gauge group, $`\mathrm{\Sigma }`$ a $`D`$dimensional manifold which admits a principal $`G`$bundle with connection over $`\mathrm{\Sigma }`$. Let us denote the pull-back to $`\mathrm{\Sigma }`$ of the connection by local sections by $`A_a^i`$ where $`a,b,c,..=1,..,D`$ denote tensorial indices and $`i,j,k,..=1,..,dim\left(G\right)`$ denote indices for the Lie algebra of $`G`$. We will denote the set of all smooth connections by $`𝒜`$ and endow it with a globally defined metric topology of the Sobolev kind
$$d_\rho [A,A^{}]:=\sqrt{\frac{1}{N}_\mathrm{\Sigma }d^Dx\sqrt{det\left(\rho \right)\left(x\right)}\text{tr}\left(\left[A_aA_a^{}\right]\left(x\right)\left[A_bA_b^{}\right]\left(x\right)\right)\rho ^{ab}\left(x\right)}$$
(2.1)
where $`\text{tr}\left(\tau _i\tau _j\right)=N\delta _{ij}`$ is our choice of normalization for the generators of a Lie algebra $`Lie\left(G\right)`$ of rank $`N`$ and our conventions are such that $`[\tau _i,\tau _j]=2f_{ij}^k\tau _k`$ define the structure constants of $`Lie\left(G\right)`$. Here $`\rho _{ab}`$ is a fiducial metric on $`\mathrm{\Sigma }`$ of everywhere Euclidean signature. In what follows we assume that either $`D2`$ ( for $`D=2`$, (2.1) depends only on the conformal structure of $`\rho `$ and cannot guarantee convergence for arbitrary fall-off conditions on the connections) or that $`D=2`$ and the fields $`A`$ are Lebesgue integrable.
Let $`\mathrm{\Gamma }_0^\omega `$ be the set of all piecewise analytic, oriented graphs $`\gamma `$ embedded into $`\mathrm{\Sigma }`$ and denote by $`E\left(\gamma \right)`$ and $`V\left(\gamma \right)`$ respectively its sets of oriented edges $`e`$ and vertices $`v`$ respectively. One can extend the framework to certain, tame piecewise smooth graphs but the description becomes more complicated and we refrain from doing this here. More important is the extension to infinite piecewise analytical graphs $`\mathrm{\Gamma }_\sigma ^\omega `$ about which much will be said in the first reference of . For the purpose of this paper it will be sufficient to stick to $`\mathrm{\Gamma }_0^\omega `$ which is sufficient, e.g. if $`\mathrm{\Sigma }`$ is compact. All the properties that are derived here for $`\mathrm{\Gamma }_0^\omega `$ readily extend to $`\mathrm{\Gamma }_\sigma ^\omega `$ as one can easily check.
We denote by $`h_e\left(A\right)`$ the holonomy of $`A`$ along $`e`$ and say that a function $`f`$ on $`𝒜`$ is cylindrical with respect to $`\gamma `$ if there exists a function $`f_\gamma `$ on $`G^{\left|E\left(\gamma \right)\right|}`$ such that $`f=p_\gamma ^{}f_\gamma =fp_\gamma `$ where $`p_\gamma \left(A\right)=\left\{h_e\left(A\right)\right\}_{eE\left(\gamma \right)}`$. The set of such functions is denoted by $`\mathrm{\Phi }_\gamma `$. Holonomies are invariant under reparameterizations of the edge and in this article we take edges always to be analytic diffeomorphisms between $`[0,1]`$ and a one-dimensional submanifold of $`\mathrm{\Sigma }`$. Gauge transformations are functions $`g:\mathrm{\Sigma }G;xg\left(x\right)`$ and they act on holonomies as $`h_eg\left(e\left(0\right)\right)h_eg\left(e\left(1\right)\right)^1`$.
A particularly useful set of cylindrical functions are the so-called spin-netwok functions . A spin-network function is labelled by a graph $`\gamma `$, a set of irreducible representations $`\stackrel{}{\pi }=\left\{\pi _e\right\}_{eE\left(\gamma \right)}`$ (choose from each equivalence class of equivalent representations once and for all a fixed representant), one for each edge of $`\gamma `$, and a set $`\stackrel{}{c}=\left\{c_v\right\}_{vV\left(\gamma \right)}`$ of contraction matrices, one for each vertex of $`\gamma `$, which contract the indices of the tensor product $`_{eE\left(\gamma \right)}\pi _e\left(h_e\right)`$ in such a way that the resulting function is gauge invariant. We denote spin-network functions as $`T_I`$ where $`I=\{\gamma ,\stackrel{}{\pi },\stackrel{}{c}\}`$ is a compound label. One can show that these functions are linearly independent.
The set of finite linear combinations of spin-network functions forms an Abelian algebra $``$ of functions on $`𝒜`$. By completing it with respect to the sup-norm topology it becomes an Abelian C algebra (here the compactness of $`G`$ is crucial). The spectrum $`\overline{𝒜}`$ of this algebra, that is, the set of all algebraic homomorphisms $`\text{ }\mathrm{C}`$ is called the quantum configuration space. This space is equipped with the Gel’fand topology, that is, the space of continuous functions $`C^0\left(\overline{𝒜}\right)`$ on $`\overline{𝒜}`$ is given by the Gel’fand transforms of elements of $``$. Recall that the Gel’fand transform is given by $`\stackrel{~}{f}\left(\overline{A}\right):=\overline{A}\left(f\right)\overline{A}\overline{𝒜}`$. It is easy to see that $`\overline{𝒜}`$ with this topology is a compact Hausdorff space. Obviously, the elements of $`𝒜`$ are contained in $`\overline{𝒜}`$ and one can show that $`𝒜`$ is even dense . Generic elements of $`\overline{𝒜}`$ are, however, distributional.
The idea is now to construct a Hilbert space consisting of square integrable functions on $`\overline{𝒜}`$ with respect to some measure $`\mu `$. Recall that one can define a measure on a locally compact Hausdorff space by prescribing a positive linear functional $`\chi _\mu `$ on the space of continuous functions thereon. The particular measure we choose is given by $`\chi _{\mu _0}\left(\stackrel{~}{T}_I\right)=1`$ if $`I=\{\left\{p\right\},\stackrel{}{0},\stackrel{}{1}\}`$ and $`\chi _{\mu _0}\left(\stackrel{~}{T}_I\right)=0`$ otherwise. Here $`p`$ is any point in $`\mathrm{\Sigma }`$, $`0`$ denotes the trivial representation and $`1`$ the trivial contraction matrix. In other words, (Gel’fand transforms of) spin-network functions play the same role for $`\mu _0`$ as Wick-polynomials do for Gaussian measures and like those they form an orthonormal basis in the Hilbert space $`:=L_2(\overline{𝒜},d\mu _0)`$ obtained by completing their finite linear span $`\mathrm{\Phi }`$.
An equivalent definition of $`\overline{𝒜},\mu _0`$ is as follows :
$`\overline{𝒜}`$ is in one to one correspondence, via the surjective map $`H`$ defined below, with the set $`\overline{𝒜}^{}:=\text{Hom}(𝒳,G)`$ of homomorphisms from the groupoid $`𝒳`$ of composable, holonomically independent, analytical paths into the gauge group. The correspondence is explicitly given by $`\overline{𝒜}\overline{A}H_{\overline{A}}\text{Hom}(𝒳,G)`$ where $`𝒳eH_{\overline{A}}\left(e\right):=\overline{A}\left(h_e\right)=\stackrel{~}{h}_e\left(\overline{A}\right)G`$ and $`\stackrel{~}{h}_e`$ is the Gel’fand transform of the function $`𝒜Ah_e\left(A\right)G`$. Consider now the restriction of $`𝒳`$ to $`𝒳_\gamma `$, the groupoid of composable edges of the graph $`\gamma `$. One can then show that the projective limit of the corresponding cylindrical sets $`\overline{𝒜}_\gamma ^{}:=\text{Hom}(𝒳_\gamma ,G)`$ coincides with $`\overline{𝒜}^{}`$. Moreover, we have $`\left\{\left\{H\left(e\right)\right\}_{eE\left(\gamma \right)};H\overline{𝒜}_\gamma ^{}\right\}=\left\{\left\{H_{\overline{A}}\left(e\right)\right\}_{eE\left(\gamma \right)};\overline{A}\overline{𝒜}\right\}=G^{\left|E\left(\gamma \right)\right|}`$. Let now $`f`$ be a function cylindrical over $`\gamma `$ then
$$\chi _{\mu _0}\left(\stackrel{~}{f}\right)=_{\overline{𝒜}}𝑑\mu _0\left(\overline{A}\right)\stackrel{~}{f}\left(\overline{A}\right)=_{G^{\left|E\left(\gamma \right)\right|}}_{eE\left(\gamma \right)}d\mu _H\left(h_e\right)f_\gamma \left(\left\{h_e\right\}_{eE\left(\gamma \right)}\right)$$
where $`\mu _H`$ is the Haar measure on $`G`$. As usual, $`𝒜`$ turns out to be contained in a measurable subset of $`\overline{𝒜}`$ which has measure zero with respect to $`\mu _0`$.
This concludes the definition of the kinematical Hilbert space $``$, of the quantum configuration space $`\overline{𝒜}`$ and of the classical configuration space. What about the classical and quantum phase space ? This question has actually so far not been analysed satisfactorily in the literature, partial results are scattered over a number of papers. We therefore begin the next section with this issue.
## 3 Symplectic Manifolds Labelled by Graphs
### 3.1 Standard Continuum Symplectic Structures
Let us first recall the usual infinite dimensional symplectic geometry that underlies gauge theories.
Let $`F_i^a`$ be a Lie algebra valued vector density test field of weight one and let $`f_a^i`$ be a Lie algebra valued covector test field. Let, as before $`A_a^i`$ be a the pull-back of a connection to $`\mathrm{\Sigma }`$ and consider a vector bundle of electric fields, that is, of Lie algebra valued vector densities of weight one whose bundle projection to $`\mathrm{\Sigma }`$ we denote by $`E_i^a`$. We consider the smeared quantities
$$F\left(A\right):=_\mathrm{\Sigma }d^DxF_i^aA_a^i\text{ and }E\left(f\right):=_\mathrm{\Sigma }d^DxE_i^af_a^i$$
(3.1)
While both are diffeomorphism covarinat it is only the latter which is gauge covariant, one reason to consider the singular smearing functions discussed below. The choice of the space of pairs of test fields $`(F,f)𝒮`$ depends on the boundary conditions on the space of connections and electric fields which in turn depends on the topology of $`\mathrm{\Sigma }`$ and will not be specified in what follows.
Consider the set $`M`$ of all pairs of smooth functions $`(A,E)`$ on $`\mathrm{\Sigma }`$ such that (3.1) is well defined for any $`(F,f)𝒮`$. We wish to endow it with a manifold structure and a symplectic structure, that is, we wish to turn it into an infinite dimensional symplectic manifold.
We define a topology on $`M`$ through the metric :
$`d_{\rho ,\sigma }[(A,E),(A^{},E^{})]`$ (3.2)
$`:=`$ $`\sqrt{{\displaystyle \frac{1}{N}}{\displaystyle _\mathrm{\Sigma }}d^Dx\left[\sqrt{det\left(\rho \right)}\rho ^{ab}\text{tr}\left(\left[A_aA_a^{}\right]\left[A_bA_b^{}\right]\right)+{\displaystyle \frac{\sigma _{ab}\text{tr}\left(\left[E^aE^a\right]\left[E^bE^b\right]\right)}{\sqrt{det\left(\sigma \right)}}}\right]}`$
where $`\rho _{ab},\sigma _{ab}`$ are again fiducial metrics on $`\mathrm{\Sigma }`$ of everywhere Euclidean signature. Their fall-off behaviour has to be suited to the boundary conditions of the fields $`A,E`$ at spatial infinity. Notice that the metric (3.2) is gauge invariant (and thus globally defined) and diffeomorphism covariant and that $`d_{\rho ,\sigma }[(A,E),(A^{},E^{})]=d_\rho [A,A^{}]+d_\sigma [E,E^{}]`$ (recall (2.1)).
Now, while the space of electric fields in Yang-Mills theory is a vector space, the space of connections is only an affine space. However, as we have also applications in general relativity with asymptotically Minkowskian boundary conditions in mind, also the space of electric fields will in general not be a vector space. Thus, in order to induce a norm from (3.2) we proceed as follows : Consider an atlas of $`M`$ consisting only of $`M`$ itself and choose a fiducial background connection and electric field $`A^{\left(0\right)},E^{\left(0\right)}`$ (for instance $`A^{\left(0\right)}=0`$). We define the global chart
$$\phi :M;(A,E)(AA^{\left(0\right)},EE^{\left(0\right)})$$
(3.3)
of $`M`$ onto the vector space of pairs $`(AA^{\left(0\right)},EE^{\left(0\right)})`$. Obviously, $`\phi `$ is a bijection. We topologize $``$ in the norm
$$(AA^{\left(0\right)},EE^{\left(0\right)})_{\rho \sigma }:=\sqrt{d_{\rho \sigma }[(A,E),(A^{\left(0\right)},E^{\left(0\right)})]}$$
(3.4)
The norm (4) is of course no longer gauge and diffeomorphism covariant since the fields $`A^{\left(0\right)},E^{\left(0\right)}`$ do not transform, they are backgrond fields. We need it, however, only in order to encode the fall-off behaviour of the fields which are independent of gauge – and diffeomorphism covariance.
Notice that the metric induced by this norm coincides with (3.2). In the terminology of weighted Sobolev spaces the completion of $``$ in the norm (3.4) is called the Sobolev space $`H_{0,\rho }^2\times H_{0,\sigma ^1}^2`$, see e.g. . We will call the completed space $``$ again and its image under $`\phi ^1`$, $`M`$ again (the dependence of $`\phi `$ on $`(A^{\left(0\right)},E^{\left(0\right)})`$ will be suppressed). Thus, $``$ is a normed, complete vector space, that is, a Banach space, in fact it is even a Hilbert space. Moreover, we have modelled $`M`$ on the Banach space $``$, that is, $`M`$ acquires the structure of a (so far only topological) Banach manifold. However, since $`M`$ can be covered by a single chart and the identity map on $``$ is certainly $`C^{\mathrm{}}`$, $`M`$ is actually a smooth manifold. The advantage of modelling $`M`$ on a Banach manifold is that one can take over almost all the pleasant properties from the finite dimensional case to the infinite dimensional one (in particular, the inverse function theorem).
Next we study differential geometry on $`M`$ with the standard techniques of calculus on infinite dimensional manifolds (see e.g. ). We will not repeat all the technicalities of the definitions involved, the interested reader is referred to the literature quoted.
* A function $`f:M\text{ }\mathrm{C}`$ on $`M`$ is said to be differentiable at $`m`$ if $`g:=f\phi ^1:\text{ }\mathrm{C}`$ is differentiable at $`u=\phi \left(m\right)`$, that is, there exist bounded linear operators $`Dg_u,Rg_u:\text{ }\mathrm{C}`$ such that
$$g\left(u+v\right)g\left(u\right)=\left(Dg_u\right)v+\left(Rg_u\right)v\text{ where }\underset{v0}{lim}\frac{\left|\left(Rg_u\right)v\right|}{v}=0.$$
(3.5)
$`Df_m:=Dg_u`$ is called the functional derivative of $`f`$ at $`m`$ (notice that we identify, as usual, the tangent space of $`M`$ at $`m`$ with $``$). The definition extends in an obvious way to the case where $`\mathrm{C}`$ is replaced by another Banach manifold. The equivalence class of functions differentiable at $`m`$ is called the germ $`G\left(m\right)`$ at $`m`$. Here two functions are said to be equivalent provided they coincide in a neighbourhood containing $`m`$.
* In general, a tangent vector $`v_m`$ at $`mM`$ is an equivalence class of triples $`(U,\phi ,v_m)`$ where $`(U,\phi )`$ is a chart of the atlas of $`M`$ containing $`m`$ and $`v_m`$. Two triples are said to be equivalent provided that $`v_m^{}=D\left(\phi ^{}\phi ^1\right)_{\phi \left(m\right)}v_m`$. In our case we have only one chart and equivalence becomes trivial. Tangent vectors at $`m`$ can be considererd as derivatives on the germ $`G\left(m\right)`$ by defining
$$v_m\left(f\right):=\left(Df_m\right)v_m=\left(D\left(f\phi ^1\right)_{\phi \left(m\right)}\right)v_m$$
(3.6)
Notice that the definition depends only on the equivalence class and not on the representant. The set of vectors tangent at $`m`$ defines the tangent space $`T_m\left(M\right)`$ of $`M`$ at $`m`$.
* The cotangent space $`T_m^{}\left(M\right)`$ is the topological dual of $`T_m\left(M\right)`$, that is, the set of continuous linear functionals on $`T_m\left(M\right)`$. It is obviously isomorphic with $`^{}`$, the topological dual of $``$. Since our model space $``$ is reflexive (it is a Hilbert space) we can naturally identify tangent and cotangent space (by the Riesz lemma) which also makes the definition of contravariant tensors less ambiguous. We will, however, not need them for what follows. Similarly, one defines the space of $`p`$covariant tensors at $`mM`$ as the space of continuous $`p`$linear forms on the $`p`$fold tensor product of $`T_m\left(M\right)`$.
* So far the fact that $``$ is a Banach manifold was not very crucial. But while the tangent bundle $`T\left(M\right)=_{mM}T_m\left(M\right)`$ carries a natural manifold structure modelled on $`\times `$ for a general Fréchet space (or even locally convex space) $``$ the cotangent bundle $`T^{}\left(M\right)=_{mM}T_m^{}\left(M\right)`$ carries a manifold structure only when $``$ is a Banach space as one needs the inverse function theorem to show that each chart is not only a differentiable bijection but that also its inverse is differentiable. In our case again there is no problem. We define differentiable vector fields and $`p`$covariant tensor fields as cross sections of the corresponding fibre bundles.
* A differential form of degree $`p`$ on $`M`$ or $`p`$form is a cross section of the fibre bundle of completely skew continuous $`p`$linear forms. Exterior product, pull-back, exterior differential, interior product with vector fields and Lie derivatives are defined as in the finite dimensional case.
###### Definition 3.1
Let $`M`$ be a differentiable manifold modelled on a Banach space $``$. A weak respectively strong symplectic structure $`\mathrm{\Omega }`$ on $`M`$ is a closed 2-form such that for all $`mM`$ the map
$$\mathrm{\Omega }_m:T_m\left(M\right)T_m^{}\left(M\right);v_m\mathrm{\Omega }(v_m,.)$$
(3.7)
is an injection respectively a bijection.
Strong symplectic structures are more useful because weak symplectic structures do not allow us to define Hamiltonian vector fields through the definition $`DL+i_{\chi _L}\mathrm{\Omega }=0`$ for differentiable $`L`$ on $`M`$ and Poisson brackets through $`\{f,g\}:=\mathrm{\Omega }(\chi _f,\chi _g)`$. Thus we define finally a strong symplectic structure for our case by
$$\mathrm{\Omega }((f,F),(f^{},F^{})):=_\mathrm{\Sigma }d^Dx\left[F_i^af_a^iF_i^af_a^i\right]\left(x\right)$$
(3.8)
for any $`(f,F),(f^{},F^{})`$. To see that $`\mathrm{\Omega }`$ is a strong symplectic structure we observe first that the integral kernel of $`\mathrm{\Omega }`$ is constant so that $`\mathrm{\Omega }`$ is clearly exact, so, in particular, closed. Next, let $`\theta ^{}`$. To show that $`\mathrm{\Omega }`$ is a bijection it suffices to show that it is a surjection (injectivity follows trivially from linearity). We must find $`(f,F)`$ so that $`\theta (.)=\mathrm{\Omega }((f,F),.)`$. Now by the Riesz lemma there exists $`(f_\theta ,F_\theta )`$ such that $`\theta (.)=<(f_\theta ,F_\theta ),.>`$ where $`<.,.>`$ is the inner product induced by (3.4). Comparing (3.4) and (3.8) we see that we have achieved our goal provided that the functions
$$F_i^a:=\rho ^{ab}\sqrt{det\left(\rho \right)}f_{b\theta }^i,f_a^i:=\frac{\sigma _{ab}}{\sqrt{det\left(\rho \right)}}F_{i\theta }^b$$
(3.9)
are elements of $``$. Inserting the definitions we see that this will be the case provided that the functions $`\rho ^{cd}\sigma _{ca}\sigma _{db}/\sqrt{det\left(\rho \right)}`$ and $`det\left(\rho \right)\sigma _{cd}\rho ^{ca}\rho ^{db}/\sqrt{det\left(\sigma \right)}`$ respectively fall off at least as $`\sigma _{ab}/\sqrt{det\left(\sigma \right)}`$ and $`\rho ^{ab}\sqrt{det\left(\rho \right)}`$ respectively. In physical applications these metrics are usually chosen to be of the form $`1+O\left(1/r\right)`$ where $`r`$ is an asymptotical radius function so that these conditions are certainly satisfied. Therefore, $`(f,F)`$ and our small lemma is established.
Let us compute the Hamiltonian vector field of a function $`L`$ on our $`M`$. By definition for all $`(f,F)`$ we have at $`m=(A,E)`$
$$DL_m(f,F)=_\mathrm{\Sigma }d^Dx[\left(DL_m\right)_i^af_a^i+\left(DL_m\right)_a^iF_i^a]=_\mathrm{\Sigma }d^Dx[\left(\chi _{Lm}\right)_i^af_a^i\left(\chi _{Lm}\right)_a^iF_i^a$$
(3.10)
thus $`\left(\chi _L\right)_i^a=\left(DL\right)_i^a`$ and $`\left(\chi _L\right)_a^i=\left(DL\right)_a^i`$. Obviously, this defines a bounded operator on $``$ if and only if $`L`$ is differentiable. Finally, the Poisson bracket is given by
$$\{L,L^{}\}_m=\mathrm{\Omega }(\chi _L,\chi _L^{})=_\mathrm{\Sigma }d^Dx\left[\left(DL_m\right)_a^i\left(DL_m^{}\right)_i^a\left(DL_m\right)_i^a\left(DL_m^{}\right)_a^i\right]$$
(3.11)
It is easy to see that $`\mathrm{\Omega }`$ has the symplectic potential $`\mathrm{\Theta }`$, a one-form on $`M`$, defined by
$$\mathrm{\Theta }_m\left((f,F)\right)=_\mathrm{\Sigma }d^DxE_i^af_a^i$$
(3.12)
since
$$D\mathrm{\Theta }_m((f,F),(f^{},F^{})):=\left(D\left(\mathrm{\Theta }_m\right)(f,F)\right)(f^{},F^{})\left(D\left(\mathrm{\Theta }_m\right)(f^{},F^{})\right)(f,F)$$
and $`DE_i^a\left(x\right)_m(f,F)=F_i^a\left(x\right)`$ as follows from the definition.
Coming back to the choice of $`𝒮`$, it will in general be a subspace of $``$ so that (3.1) still converges. We can now compute the Poisson brackets between the functions $`F\left(A\right),E\left(f\right)`$ on $`M`$ and find
$$\{E\left(f\right),E\left(f^{}\right)\}=\{F\left(A\right),F^{}\left(A\right)\}=0,\{E\left(f\right),A\left(F\right)\}=F\left(f\right)$$
(3.13)
Remark :
In physicists notation one often writes $`\left(DL_m\right)_a^i\left(x\right):=\frac{\delta L}{\delta A_a^i\left(x\right)}`$ etc. and one writes the symplectic structure as $`\mathrm{\Omega }=d^DxDE_i^a\left(x\right)DA_a^i\left(x\right)`$.
### 3.2 New Symplectic Structure
As outlined in the introduction we wish to define symplectic manifolds $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ for every $`\gamma \mathrm{\Gamma }`$ which “in some sense come from $`(M,\mathrm{\Omega })`$”. In order to do this we need besides graphs new extended objects associated with them. This is the topic of the following subsubsection. The class of graphs $`\gamma `$ that we have in mind consists of the set $`\mathrm{\Gamma }_\sigma ^\omega `$ of piecewise analytic, $`\sigma `$-finite graphs. These are graphs with an at most countable number of edges and such that for every compact subset $`U`$ of the locally compact manifold $`\mathrm{\Sigma }`$ the restriction $`U\gamma `$ is a piecewise analytic, finite graph in $`\mathrm{\Gamma }_0^\omega `$. The precise definition of these graphs and their properties as well as the extension of the quantum kinematical framework to them needs the framework of the infinite tensor product of Hilbert spaces and will therefore be postponed to the first reference of . For the rest of this paper one may without doing any mistake think of $`\gamma `$ as an element of $`\mathrm{\Gamma }_0^\omega `$.
#### 3.2.1 Dual Polyhedral Decompositions
###### Definition 3.2
A polyhedral decomposition $`P`$ of $`\mathrm{\Sigma }`$ is a subdivision of $`\mathrm{\Sigma }`$ into closed compact regions $`\mathrm{\Delta }`$, that is $`\mathrm{\Sigma }=_{\mathrm{\Delta }P}\mathrm{\Delta }`$, each of which is diffeomorphic to a polyhedron in flat space and intersects any other polyhedron only in the set of points of their common boundary (of codimension at least one).
Note that we allow for decompositions with a countably infinite number of polyhedra in case that $`\mathrm{\Sigma }`$ is not compact. We also do not insist that the decomposition be convex as this would require the choice of a background metric (rather of a diffeomorphism invariance class because convexity is defined by geodesity of curves which is a diffeomorphism invariant notion).
###### Definition 3.3
Let $`P`$ be a polyhedronal decomposition of $`\mathrm{\Sigma }`$, pick any polyhedron $`\mathrm{\Delta }P`$ and consider its boundary $`\mathrm{\Delta }`$.
i) A face $`S`$ of $`\mathrm{\Delta }`$ (a maximal, connected, analytic subset of its boundary of codimension one) is called a “standard face” provided that
1) $`S`$ is a submanifold of $`\mathrm{\Sigma }`$ of codimension one, diffeomorphic to a standard cube $`[1,1]^{D1}`$ in $`\text{ }\mathrm{R}^{D1}`$.
2) $`S`$ has no boundary, $`S=\mathrm{}`$, i.e. it is open.
3) $`S`$ is contained in the domain of a chart of $`\mathrm{\Sigma }`$.
4) $`S`$ is maximal, that is, there does not exist any $`S^{}\mathrm{\Delta }`$ properly containing $`S`$ which satisfies 1), 2) and 3).
ii) As $`S`$ is an open submanifold of $`\mathrm{\Sigma }`$ of codimension one and for some $`\mathrm{\Delta }P,S\mathrm{\Delta }`$ is contained in the domain of a chart $`(U,\phi )`$ of an atlas of $`\mathrm{\Sigma }`$ there exists an open subset $`VU`$ containing $`S`$, divided by $`S`$ into two halves and a diffeomorphism $`\phi ^{}`$ that maps $`V,S`$ respectively to $`\text{ }\mathrm{R}^D`$ and the hypersurface $`x^D=0`$ respectively. An orientation of $`S`$ is given by a choice of which of the half spaces or “sides” given by the set of points satisfying $`x^D>0`$ or $`x^D<0`$ we call “up” or “down”.
iii) A polyhedronal decomposition of $`\mathrm{\Sigma }`$ is said to be oriented provided the collection of all its standard faces have been oriented.
This definition just formalizes the intuitive idea of a face of a polyhedron with a regular shape. Notice that a face is always shared by precisely two polyhedra of the decomposition. In D=1,2,3 respectively simplicial polyhedra are given by closed lines, triangles and tetrahedra respectively and their faces are open points, lines and triangles respectively. The notion of an orientation is clearly independent of the chart employed. Also, faces are always orientable even if $`\mathrm{\Sigma }`$ is not orientable and even if $`\mathrm{\Sigma }`$ is orientable the orientation of a face can be opposite to the orientation induced from $`\mathrm{\Sigma }`$ on $`S`$ (as a submanifold).
###### Definition 3.4
Given a graph $`\gamma \mathrm{\Gamma }`$ we say that an oriented, polyhedral decomposition $`P`$ of $`\mathrm{\Sigma }`$ is dual to $`\gamma `$ provided that :
1) Given an edge $`e`$ of $`\gamma `$ there exists precisely one standard face $`S_e`$ from the collection of all faces of all the polyhedra of the decomposition which is intersected by $`e`$.
2) The edge $`e`$ intersects $`S_e`$ transversally, that is, a) $`\{p_e\}:=eS_e`$ consists of a single point which is an interior point of $`e`$ (it is necessarily an interior point of $`S_e`$ since $`S_e`$ is open) and b) there exists an open neighbourhood $`V`$ of $`p_e`$ and a diffeomorphism which maps $`V`$ to $`\text{ }\mathrm{R}^D`$, $`p_e`$ to the origin of $`\text{ }\mathrm{R}^D`$, $`S_eV`$ to the plane $`x^D=0`$ and $`eV`$ to the line $`x^1=..=x^{D1}=0`$.
3) The orientations of $`S_e`$ and $`e`$ agree, that is, if under the diffeomorphism outlined in 2) the edge points into the half space $`x^D>0`$ or $`x^D<0`$ respectively then that half space corresponds to the “up” side of $`S_e`$.
4) The decomposition is irreducible, that is, one cannot reduce the number of polyhedra of the decomposition (by removing faces) without destroying at least one of the properties 1)-3).
Dual decompositions certainly exist in any dimension : A first example is given by a cubic lattice in any dimension (every vertex is 2D-valent) where we assume that $`\mathrm{\Sigma }`$ has no boundary (periodic boundary conditions). The dual lattice is unique up to diffeomorphisms and corresponds to D-cubes around every vertex. A second example (again with periodic boundary conditions) is given by a simplicial lattice (every vertex is (D+1)-valent) and corresponds to a simplicial decomposition with D-simplices around every vertex. Again the dual lattice is uniquely determined up to diffeomorphisms. Whether or not this uniqueness is a general feature of dual decompositions is not clear but we have the following.
###### Theorem 3.1
Given a graph $`\gamma \mathrm{\Gamma }`$, a dual, oriented, polyhedral decomposition $`P_\gamma `$ of it exists and it is unique up to diffeomorphisms and up to the number of possibilities to obey the Euler (or Dehn-Sommerfeld in $`D3`$) relation for the various polyhedra under the reduction process.
Proof of Theorem 3.1 :
Existence :
Given a graph $`\gamma `$ consider for each vertex $`v`$ a closed neighbourhood $`U_v`$ of $`v`$ such that the various $`U_v`$ are mutually disjoint to begin with. Moreover, we can choose them to be closed balls such that their boundaries have the topology of the sphere. Now distort the $`U_v`$ along the edges $`e`$ incident at $`v`$ without changing their topology until for any edge $`e`$ with end points $`v,v^{}`$ the $`U_v,U_v^{}`$ intersect in precisely one point $`p_e`$ which is obviously an interior point of $`e`$, otherwise they are mutually non-intersecting. Now blow up the $`U_v`$ even more, keeping the $`p_e`$ fixed until each $`U_v`$ looks like a solid ball with a spherical boundary except for $`n\left(v\right)`$ cusps $`S_e`$ of the topology of cubes $`[1,1]^{D1}`$ with the $`p_e`$ as an interior point. $`S_e`$ is the only set in which $`U_v,U_v^{},\{v,v^{}\}=e`$ intersect. By shifting the $`p_e`$ and varying the size of the cusps we can achieve that the $`S_e`$ are mutually non-intersecting, contained in the domain of a chart and intersect $`e`$ transversally. We can therefore equip them with an orientation that agrees with that of $`e`$. Altogether, the $`U_v`$ have now mutated to polyhedrons with $`n\left(v\right)`$ open, smooth standard faces $`S_e,ve`$ and the additional closed connected component of its boundary consisting of $`S_v:=U_v_{ve}S_e`$ . Finally, consider the polyhedronal decomposition of $`\mathrm{\Sigma }`$ consisting of the $`U_v`$ and the remainder $`\mathrm{\Sigma }_vU_v`$. This decomposition satisfies all the properties of a dual decomposition except, possibly, for the irreducibility requirement. To satisfy it, we remove the faces $`S_v`$ by letting the $`S_e`$ share their boundaries, as long as compatible with the Euler relation between the number of connected components of all possible subcomplexes of a polyhedron. (For instance, in D=3 the relation is given by $`f=ke+\chi `$ where $`f,k,e`$ denotes the number of faces, links and corners of a polyhedron and $`\chi `$ is essentially the Euler characteristic of the manifold that the polyhedron triangulates). That this is always possible follows by the lemma of choice (or Zorn’s lemma).
Uniqueness :
The constructive proof just given just has just fixed the topology of the final dual polyhedronal decomposition reached and is therefore unique at most up to diffemorphisms. Morover, the number of possible irreducible decompositions reachable obviously equals the number of possible solutions to the just mentioned Euler relation.
$`\mathrm{}`$
The non-uniqueness of $`P_\gamma `$ does not affect us because we will use only the $`S_e`$ which are determined up to diffeomorphisms.
Notice that on the other hand, if we are given an oriented cellular decomposition of $`\mathrm{\Sigma }`$ into polyhedra with open faces, there is, up to diffeomorphism equivalence, only one graph such that the decomposition is dual to it. This follows easily from the fact that there is no choice but choosing an interior point of each polyhedron and connecting common faces between polyhedra with transversal edges running between the corresponding interior points with the obvious orientation.
From now on we pick for each graph $`\gamma `$ an oriented, polyhedral, dual decomposition $`P_\gamma `$. By theorem 3.1 we can do this in such a way that, if $`\gamma \gamma ^{}`$ are diffeomorphic, then $`P_\gamma ,P_\gamma ^{}`$ are diffeomorphic. Notice, however, that it is not possible to require that for any diffeomorphism $`\phi `$ we have $`P_{\phi \left(\gamma \right)}=\phi \left(P_\gamma \right)`$ since there are many diffeomorphisms, say $`\phi _1,\phi _2`$ such that $`\phi _1\left(\gamma \right)=\phi _2\left(\gamma \right)`$ but $`\phi _1\left(P_\gamma \right)\phi _2\left(P_\gamma \right)`$.
Also, a word of caution is appropriate with respect to refinements : The set $`\mathrm{\Gamma }`$ is partially ordered by inclusion and for each pair $`\gamma ,\gamma ^{}`$ there exists a bigger (refined) graph $`\stackrel{~}{\gamma }`$ containing both of them, for instance the graph $`\gamma \gamma ^{}`$. However, in general it will not be true that there exists a refined graph such that $`P_\gamma ,P_\gamma ^{}`$ are both contained in $`P_{\stackrel{~}{\gamma }}`$. This can happen only if the graphs under consideration have a high degree of symmetry as e.g. cubic lattices as one can show.
#### 3.2.2 The Graph Phase Space From the Continuum Phase Space
We are now ready to derive a symplectic manifold $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ for each $`\gamma \mathrm{\Gamma }`$ from the standard symplectic manifold $`(M,\mathrm{\Omega })`$ considered in the beginning of this section.
###### Definition 3.5
i) Let $`S_0`$ be the interior of the subset $`[1,1]^{D1}\text{ }\mathrm{R}^{D1}`$ in the $`x^D=0`$ plane with normal orientation into the $`x^D>0`$ direction. Let $`p_0=0`$ be the origin in $`\text{ }\mathrm{R}^D`$ and $`e_0`$ be the interval $`[1/2,1/2]`$ of the $`x^D`$-axis. Let $`x_0S_0`$ and define $`\rho _0(x_0)`$ to be the straight line within $`S_0`$ connecting $`p_0`$ and $`x_0`$.
ii) Given a graph $`\gamma `$ and a dual polyhedronal decomposition $`P_\gamma `$, we call a collection of paths $`\mathrm{\Pi }_\gamma :=\{\rho _e(x)S_e;xS_e\}_{eE(\gamma )}`$ adapted to $`\gamma ,P_\gamma `$ provided there exists for each $`eE(\gamma )`$ an analytic diffeomorphism $`\phi _e:\text{ }\mathrm{R}^D\mathrm{\Sigma }`$ such that
$$(e,S_e,p_e,x,\rho _e\left(x\right))=(\phi _e\left(e_0\right),\phi _e\left(S_0\right),\phi _e\left(p_0\right),\phi _e\left(x_0\right),\phi _e\left(\rho _0\left(x_0\right)\right)$$
(3.14)
We will denote the set of triples $`(\gamma ,P_\gamma ,\mathrm{\Pi }_\gamma )`$ by $`_{}`$ or $``$ where $`=0,\sigma `$ depends on the class of the graphs that we allow. The elements $`l`$ will be called structured graphs. We also use the notation $`l=(\gamma (l),P_\gamma (l),\mathrm{\Pi }_\gamma (l))`$.
iii) Given a structured graph $`l`$, let w.l.g. $`p_e=e(1/2)`$.
Then we define the following function on $`(M,\mathrm{\Omega })`$
$$P_i^e(A,E):=\frac{1}{N}\text{tr}\left(\tau _ih_e(0,1/2)\left[_{S_e}h_{\rho _e\left(x\right)}E\left(x\right)h_{\rho _e\left(x\right)}^1\right]h_e(0,1/2)^1\right)$$
(3.15)
where $`h_e(s,t)`$ denotes the holonomy of $`A`$ along $`e`$ between the parameter values $`s<t`$, $``$ denotes the Hodge dual, that is, $`E`$ is a $`(D1)`$form on $`\mathrm{\Sigma }`$and $`E^a:=E_i^a\tau _i`$.
Notice that in contrast to similar variables used earlier in the literature the function $`P_i^e`$ is gauge covariant. Namely, under gauge transformations it transforms as $`P^eg\left(e\left(0\right)\right)P^eg\left(e\left(0\right)\right)^1`$, the price to pay being that $`P^e`$ depends on both $`A`$ and $`E`$ and not only on $`E`$. As we will see shortly, this is actually an advantage. Of course, the notation (3.15) is abusing as $`P^e`$ not only depends on $`e`$ but actually on $`S_e,\rho _e\left(x\right),xS_e`$. In the sequel, unless confusion can arise we will continue abusing notation and write $`\gamma `$ instead of $`l`$.
The problem with the functions $`h_e\left(A\right)`$ and $`P_i^e(A,E)`$ on $`M`$ is that they are not differentiable on $`M`$, that is, $`Dh_e,DP_i^e`$ are nowhere bounded operators on $``$ as one can easily see. The reason for this is, of course, that these are functions on $`M`$ which are not properly smeared with functions from $`𝒮`$, rather they are smeared with distributional test functions with support on $`e`$ or $`S_e`$ respectively. Nevertheless one would like to base the quantization of the theory on these functions as basic variables because of their gauge and diffeomorphism covariance. Indeed, under diffeomorphisms the structured graph $`l`$ is simply replaced by
$$\phi ^1\left(l\right)=(\phi ^1\left(e\right),\phi ^1\left(S_e\right),\left\{\phi ^1\left(\rho _e\left(x\right)\right);xS_e\right\})_{eE\left(\gamma \right)}$$
(3.16)
which is a structured graph again and in this sense $`h_eh_{\phi ^1\left(e\right)},P_i^eP_i^{\phi ^1\left(e\right)}`$. Furthermore, their quantizations are properly represented on the Hilbert space described in section 2 as we will see. The fact that the smearing dimensions of $`h_e`$ and $`P_i^e`$ add to $`D`$ raises some hope that one can still derive a bona fide Poisson algebra among these variables. We therefore define
###### Definition 3.6
The set of pairs $`(h_e(A),P_i^e(A,E))_{eE(\gamma )}`$ as $`(A,E)`$ varies over $`M`$ will be denoted by $`\overline{M}_{\gamma |M}`$. We also define $`\overline{M}_\gamma =(G\times Lie(G))^{|E(\gamma )|}`$.
It is easy to see that $`\overline{M}_{\gamma |M}`$ is generically a proper subset of $`\overline{M}_\gamma `$. Indeed, since the edges $`e`$ are mutually disjoint among each other except for the vertices we can find a smooth connection with support in disjoint open neighbourhoods $`U_e`$, one for each $`e`$, such that $`eU_e`$ is an open segment of $`e`$. The holonomy along those segments can be given independent values since we can vary the behaviour of $`A`$ in each $`U_e`$ independently and arbitrarily without destroying smoothness. Similar considerations hold for the momenta $`P_i^e`$. The range of these values is, however, constrained by the boundary conditions imposed by the fact that $`(A,E)`$ are points in a classical phase space subject to fall-off conditions. The bar in the notation $`\overline{M}_\gamma `$ accounts for the fact that the points of $`\overline{M}_\gamma `$ do not satisfy such regularity assumptions similar as in the case of $`\overline{𝒜}_\gamma `$.
We equip a subset $`M_\gamma `$ of $`\overline{M}_\gamma `$ with the following natural topology : Let $`(u_i,\varphi _i)_{iI}`$ be an atlas of $`G`$ where $`I`$ is a finite index set (always possible since $`G`$ is compact). For instance the $`u_i`$ could be preimages of open sets in $`\text{ }\mathrm{R}^{dim\left(G\right)}`$ under the exponential map which is locally a diffeomorphism between $`G`$ and $`Lie\left(G\right)`$. Since $`Lie\left(G\right)`$ is topologically trivial we can construct an atlas of $`G\times Lie\left(G\right)`$ by $`\left(U_i=u_i\times Lie\left(G\right),\mathrm{\Phi }_i=\varphi _i\times \text{id}\right)`$ where id denotes the identity map of $`\text{ }\mathrm{R}^{dim\left(G\right)}`$. Then $`M_\gamma `$ can be given the differentiable structure defined by the atlas
$$(\times _{eE\left(\gamma \right)}U_{i_e},\times _{eE\left(\gamma \right)}\mathrm{\Phi }_{i_e})_{i_eI}$$
(3.17)
which displays $`M_\gamma `$ as a Banach manifold modelled on $`_\gamma =\text{ }\mathrm{R}^{2dim\left(G\right)\left|E\left(\gamma \right)\right|}`$. $`_\gamma `$ is equipped with the norm
$$\{x_e,y_e\}_{eE\left(\gamma \right)}_{\rho ,\sigma }:=\sqrt{\underset{eE\left(\gamma \right)}{}\left[\rho _e^{ij}x_e^ix_e^j+\sigma _{ij}^ey_e^iy_e^j\right]}$$
(3.18)
where $`\rho _e^{ij},\sigma _{ij}^e`$ are metrics of Euclidean signature for each $`e`$. Obviously, then $`M_\gamma `$ is a proper subset of $`\overline{M}_\gamma `$.
Remark :
A connection between (3.18) and (3.2) can be given on certain graphs as follows : The idea is that (3.18) is a discretization of (3.2) so that they eqaul each other in the limit of an infinitely fine graph. Consider for simplicity $`\mathrm{\Sigma }=\text{ }\mathrm{R}^3`$ (the generalization to arbitrary $`\mathrm{\Sigma }`$ is straightforward with the tools of the next subsection) and consider a lattice $`\gamma `$ of regular cubic topology. Then edges can be labelled as $`te_I(v,t)`$ where $`v=e_I(v,0)`$ is a vertex and $`I=1,..,D`$. Let edges be images of the interval $`t[0,ϵ]`$ and define $`Y_I^a\left(v\right)=\dot{e}_I^a(v,0),n_a^I\left(v\right)=\frac{1}{\left(D1\right)!}ϵ_{ab_1..b_{D1}}ϵ^{IJ_1..J_{D1}}Y_{J_1}^{b_1}\left(v\right)..Y_{J_{D1}}^{b_{D1}}\left(v\right)`$. We now choose $`\rho _{ij}^e=\delta _{ij}f\left(e\right),\sigma _e^{ij}=\delta ^{ij}g\left(e\right)`$ for some functions $`f,g:E\left(\gamma \right)\text{ }\mathrm{R}`$ and choose the $`e_I\left(v\right)`$ in such a way that
$$\underset{I}{}f\left(e_I\left(v\right)\right)\left)Y^a_I\right(v\left)Y^b_I\right(v)=ϵ^{D2}(\sqrt{det\left(\rho \right)}\rho ^{ab}\left)\right(v)\text{ and }_Ig(e_I\left(v\right)\left)n_a^I\right(v\left)n_b^I\right(v)=ϵ^{\left(D2\right)}(\sigma _{ab}/\sqrt{det\left(\sigma \right)}\left)\right(v)$$
the metrics of (3.2) result. It is then easy to see that with $`x_i^e:=\frac{2}{N}\text{tr}\left(\tau _ih_e\left(A\right)h_e\left(A^{}\right)^1\right),y_e^i:=P_i^e(A,E)P_i^e(A^{},E^{})`$ the resulting metric (3.18) converges to the metric (3.2).
In order to proceed and to give $`M_\gamma `$ a symplectic structure derived from that of $`M`$ one must regularize the elementary functions so that one can use the symplectic structure $`\mathrm{\Omega }`$, then study the limit in which the regulator is removed and hope that the result is a well-defined symplectic structure $`\mathrm{\Omega }_\gamma `$. We choose the following regularization :
Given an edge $`e`$ we define a tube $`T^e`$ around $`e`$ to be a foliation of $`D1`$dimensional surfaces which are topologically discs. Let $`f_ϵ:DT_ϵ^e`$ be a one-parameter family of smooth functions which tends to the $`\delta `$-distribution on the unit disc $`D`$. Recall that the holonomy of a smooth connection can be written as the path ordered exponential
$$h_e\left(A\right)=1+\underset{n=1}{\overset{\mathrm{}}{}}_0^1dt_n_0^{t_n}dt_{n1}.._0^{t_2}dt_1A\left(t_1\right)..A\left(t_n\right)$$
(3.19)
where $`A\left(t\right)=\dot{e}^a\left(t\right)A_a^i\left(e\left(t\right)\right)\tau _i/2`$. We define the holonomy smeared over the tube $`T_ϵ^e`$ by
$$h_e^ϵ\left(A\right)=1+\underset{n=1}{\overset{\mathrm{}}{}}_0^1dt_n_0^{t_n}dt_{n1}.._0^{t_2}dt_1_Dd^{D1}y_nf_ϵ\left(y_n\right).._Dd^{D1}y_1f_ϵ\left(y_1\right)A_{y_1}\left(t_1\right)..A_{y_n}\left(t_n\right)$$
(3.20)
where $`A_y\left(t\right)=\dot{e}_y^a\left(t\right)A_a^i\left(e_y\left(t\right)\right)\tau _i/2`$ and $`D\times [0,1]T^e;(y,t)e_y\left(t\right)`$ defines the tube $`T^e`$ with the convention that $`e_0\left(t\right)=e\left(t\right)`$.
Likewise, let $`R^e`$ be a region foliated by surfaces diffeomorphic to $`S^e`$. Let $`X:VS^e;u:=(u_1,..,u_{D1}))X(u)`$ be a parameterization of $`S^e`$ where $`V`$ is an open submanifold of $`\text{ }\mathrm{R}^D`$ of dimension $`D1`$. Furhermore, let $`[1,1]\times VR^e;(s,u)X_s\left(u\right)`$ define $`R^e`$ with the convention that $`X_0\left(u\right)=X\left(u\right)`$. Let $`g_ϵ`$ be a one parameter family of smooth functions which tends to the $`\delta `$-distribution on the interval $`[1,1]`$ as $`ϵ0`$. Also, define $`\rho _e^s\left(x\right)`$ to be paths in $`X_s\left(V\right)=S_e^s`$ between $`p_e\left(s\right)=X_s\left(0\right)=eX_s\left(V\right)`$ and $`xX_s\left(V\right)`$ and let $`e_s`$ be a reparameterization of $`e`$ such that $`e_s\left(1/2\right)=p_e^s`$. Here we assume w.l.g. that all the surfaces $`S_e^s,s[1,1]`$ satisfy the conditions that the surface $`S_e`$ has to satisfy. We can now define
$$P_{i,s}^e(A,E):=\frac{1}{N}\text{tr}\left(\tau _ih_{e_s}(0,1/2)\left[_{S_e^s}h_{\rho _{e_s}\left(x\right)}E\left(x\right)h_{\rho _{e_s}\left(x\right)}^1\right]h_{e_s}(0,1/2)^1\right)$$
(3.21)
and then
$$P_i^{eϵ}:=_1^1𝑑sg_ϵ\left(s\right)P_{i,s}^e$$
(3.22)
Notice that the holonomies involved in (3.22) remain unsmeared as compared to (3.20). We could improve this by an additional smearing, preserving gauge covariance, but it would just blow up the subsequent calculations and would not change the end result. The careful reader is invited to fill in the missing details.
Apart from these details, the functions (3.20) and (3.22) are now written as functions of the variables $`F\left(A\right),E\left(f\right)`$ where $`F,f`$ are of the form $`F_i^a\left(x\right)=\chi _{T^e}\left(x\right)\left(f_ϵ\left(y\right)\dot{e}_y^a\left(t\right)\delta _i^j\right)_{x=e_y\left(t\right)}`$ and $`f_a^i\left(x\right)=\chi _{R^e}\left(x\right)(g_ϵ\left(s\right)ϵ_{ab_1..b_{D1}}X_{s,u_1}^{b1}\left(u\right)..X_{s,u_{D1}}^{b2}\left(u\right)\delta _i^j)_{x=X_s\left(u\right)}`$ for some $`j`$ and are thus certainly elements of $`𝒮`$. It follows that the smeared functions are functionally differentiable. Moreover, by construction, the smeared objects converge pointwise on $`M`$ to the unsmeared objects, that is
$$\underset{ϵ0}{lim}|\left(h_e^ϵ\left(A\right)\right)_{AB}\left(h_e\left(A\right)\right)_{AB})|=lim_{ϵ0}|P^{eϵ}_i(A,E)P^e_i(A,E)|=0$$
(3.23)
for all $`(A,E)M,i=1,..,dim\left(G\right),A,B`$ where $`A,B,..`$ denote group indices.
###### Theorem 3.2
The smeared variables allow us to define the following bracket $`\{.,.\}_\gamma `$ on $`M_\gamma `$ :
$`\{h_e,h_e^{}\}_\gamma `$ $`:=`$ $`\underset{ϵ_1,ϵ_20}{lim}\{h_e^{ϵ_1},h_e^{}^{ϵ_2}\}=0`$ (3.24)
$`\{P_i^e,h_e^{}\}_\gamma `$ $`:=`$ $`\underset{ϵ_10}{lim}\underset{ϵ_20}{lim}\{P_i^{eϵ_1},h_e^{}^{ϵ_2}\}=\delta _e^{}^e{\displaystyle \frac{\tau _i}{2}}h_e`$ (3.25)
$`\{P_i^e,P_j^e^{}\}_\gamma `$ $`:=`$ $`\underset{ϵ_1,ϵ_20}{lim}\{P_i^{eϵ_1},P_j^{e^{}ϵ_2}\}=\delta ^{ee^{}}f_{ij}^kP_k^e`$ (3.26)
where $`\{.,.\}`$ is the bracket on $`M`$ and convergence is meant here and in what follows to be pointwise on $`M`$.
Notice that we do not yet call $`\{.,.\}_\gamma `$ a Poisson bracket since we must check that it qualifies as a (strong) symplectic structure. This we will do in a separate step.
Proof of Theorem 3.2 :
\[1.\]
Recalling (3.13) the first identity (3.24) follows easily from the fact that $`\{h_e^{ϵ_1},h_e^{}^{ϵ_2}\}=0`$ at every finite $`ϵ_1,ϵ_2`$. It is for this reason that we do not have to smear the $`h_{\rho _e^s\left(x\right)},h_{e_s}`$ involved in (3.21) in addition to $`E`$ in order to define the brackets, there would be no extra contribution due to this fact.
\[2.\]
The second identity (3.25) is significantly more involved. We first prove the following lemma.
###### Lemma 3.1
For any $`f_a^i𝒮`$ and any path $`e`$ we have
$$\{E\left(f\right),h_e\}:=\underset{ϵ0}{lim}\{E\left(f\right),h_e^ϵ\}=_0^1𝑑t\dot{e}^a\left(t\right)f_a^i\left(e\left(t\right)\right)h_e(0,t)\frac{\tau _i}{2}h_e(t,1)$$
(3.27)
Proof of Lemma 3.1 :
We have by definition
$`\{E\left(f\right),h_e^ϵ\left(A\right)\}`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle _0^1}dt_n{\displaystyle _0^{t_n}}dt_{n1}..{\displaystyle _0^{t_2}}dt_1{\displaystyle _D}d^{D1}y_nf_ϵ\left(y_n\right)..{\displaystyle _D}d^{D1}y_1f_ϵ\left(y_1\right)\times `$ (3.28)
$`\times {\displaystyle \underset{k=1}{\overset{n}{}}}A_{y_1}\left(t_1\right)..\{E\left(f\right),A_{y_k}\left(t_k\right)\}..A_{y_n}\left(t_n\right)`$
$`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle _0^1}dt_n{\displaystyle _0^{t_n}}dt_{n1}..{\displaystyle _0^{t_2}}dt_1{\displaystyle _D}d^{D1}y_nf_ϵ\left(y_n\right)..{\displaystyle _D}d^{D1}y_1f_ϵ\left(y_1\right)\times `$
$`\times {\displaystyle \underset{k=1}{\overset{n}{}}}f_a^i\left(e_{y_k}\left(t_k\right)\right)\dot{e}_{y_k}^a\left(t_k\right)A_{y_1}\left(t_1\right)..{\displaystyle \frac{\tau _i}{2}}..A_{y_n}\left(t_n\right)`$
Relabel $`T_1=t_1,..,T_{k1}=t_{k1},t=t_k,T_k=t_{k+1},..,T_{n1}=t_n`$ and $`z_1=y_1,..,z_{k1}=y_{k1},y=y_k,z_k=y_{k+1},..,z_{n1}=y_n`$ then (3.28) becomes
$`\{E\left(f\right),h_e^ϵ\left(A\right)\}`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle _0^1}dT_{n1}{\displaystyle _0^{T_{n1}}}dT_{n2}..{\displaystyle _0^{T_{k+1}}}dT_k{\displaystyle _0^{T_k}}dt{\displaystyle _0^t}dT_{k1}..{\displaystyle _0^{T_2}}dT_1\times `$ (3.29)
$`\times {\displaystyle _D}d^{D1}yf_ϵ\left(y\right)f_a^i\left(e_y\left(t\right)\right)\dot{e}_y^a\left(t\right){\displaystyle _D}d^{D1}z_{n1}f_ϵ\left(z_{n1}\right)..{\displaystyle _D}d^{D1}z_1f_ϵ\left(z_1\right)\times `$
$`\times {\displaystyle \underset{k=1}{\overset{n}{}}}A_{z_1}\left(T_1\right)..{\displaystyle \frac{\tau _i}{2}}..A_{z_{n1}}\left(T_{n1}\right)`$
$`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle _0^1}dt{\displaystyle _D}d^{D1}yf_ϵ\left(y\right)f_a^i\left(e_y\left(t\right)\right)\dot{e}_y^a\left(t\right)\times `$
$`\times {\displaystyle _t^1}dT_{n1}{\displaystyle _t^{T_{n1}}}dT_{n2}..{\displaystyle _t^{T_{k+1}}}dT_k{\displaystyle _0^t}dT_{k1}{\displaystyle _0^{T_{k1}}}dT_{k2}..{\displaystyle _0^{T_2}}dT_1\times `$
$`\times {\displaystyle _D}d^{D1}z_{n1}f_ϵ\left(z_{n1}\right)..{\displaystyle _D}d^{D1}z_1f_ϵ\left(z_1\right){\displaystyle \underset{k=1}{\overset{n}{}}}A_{z_1}\left(T_1\right)..{\displaystyle \frac{\tau _i}{2}}..A_{z_{n1}}\left(T_{n1}\right)`$
where in the last step we have used the fact that $`1T_{n1}..T_ktT_{k1}..T_10`$ in order to make the range of integration of $`t`$ independent of the $`T_k`$. We can now easily take the limit $`ϵ0`$ with the result
$`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle _0^1}𝑑tf_a^i\left(e\left(t\right)\right)\dot{e}^a\left(t\right){\displaystyle _t^1}𝑑T_{n1}{\displaystyle _t^{T_{n1}}}𝑑T_{n2}..`$
$`..{\displaystyle _t^{T_{k+1}}}dT_k{\displaystyle _0^t}dT_{k1}{\displaystyle _0^{T_{k1}}}dT_{k2}..{\displaystyle _0^{T_2}}{\displaystyle \underset{k=1}{\overset{n}{}}}A_{z_1}\left(T_1\right)..{\displaystyle \frac{\tau _i}{2}}..A_{z_{n1}}\left(T_{n1}\right)`$ (3.30)
and writing out the path product identity for holonomies $`h_e(0,t)h_e(t,1)=h_e(0,1)=h_e`$ in the path ordered form (3.19) this collapses indeed to (3.27).
$`\mathrm{}`$
Let us now prove the second identity (3.25). Let us write $`P_i^{eϵ}`$ in the form $`E\left(f\right)`$ by choosing
$`\left(f_i^ϵ\right)_a^j\left(x\right)`$ $`=`$ $`{\displaystyle _1^1}ds{\displaystyle _V}d^{D1}u\delta (x,X_s\left(u\right)){\displaystyle \frac{1}{N}}g_ϵ\left(s\right)ϵ_{ab_1..b_{D1}}X_{s,u_1}^{b_1}\left(u\right)..X_{s,u_{D1}}^{b_{D1}}\left(u\right)\times `$ (3.31)
$`\times \text{tr}\left(\tau _ih_{e_s}(0,1/2)h_{\rho _{e_s}\left(X_s\left(u\right)\right)}\tau _jh_{\rho _{e_s}\left(X_s\left(u\right)\right)}^1h_{e_s}(0,1/2)^1\right)`$
From \[1.\] it is clear that, although $`f_a^j`$ depends on $`A`$, as far as (3.25) is concerned, we can treat it as a numerical function. By Lemma 3.1 we then have
$$\{P_i^e,h_e^{}\}_\gamma =\underset{ϵ0}{lim}\{E\left(f_i^ϵ\right),h_e^{}\}=\underset{ϵ0}{lim}_0^1𝑑t\dot{e}^a\left(t\right)\left(f_i^ϵ\right)_a^j\left(e^{}\left(t\right)\right)h_e^{}(0,t)\frac{\tau _j}{2}h_e^{}(t,1)$$
(3.32)
Suppose first that $`ee^{}`$. Then, no matter how complicated $`\gamma ,P_\gamma `$ look, for sufficiently small $`ϵ`$ the edge $`e^{}`$ does not intersect the region $`R_ϵ^e`$ and thus (3.32) vanishes. If $`e=e^{}`$ then $`e`$ intersects $`R_ϵ^e`$ for any value of $`ϵ`$ and we find
$`\{P_i^e,h_e\}_\gamma =\underset{ϵ0}{lim}{\displaystyle \frac{1}{N}}{\displaystyle _1^1}dsg_ϵ\left(s\right){\displaystyle _0^1}dt{\displaystyle _V}d^{D1}u\delta (e\left(t\right),X_s\left(u\right))\dot{e}^a\left(t\right)\times `$
$`\times ϵ_{ab_1..b_{D1}}X_{s,u_1}^{b_1}\left(u\right)..X_{s,u_{D1}}^{b_{D1}}\left(u\right)\text{tr}(\tau _ih_{e_s}(0,1/2)h_{\rho _{e_s}\left(X_s\left(u\right)\right)}\times `$
$`\times \tau _jh_{\rho _{e_s}\left(X_s\left(u\right)\right)}^1h_{e_s}(0,1/2)^1\left)h_e\right(0,t\left){\displaystyle \frac{\tau _j}{2}}h_e\right(t,1)`$ (3.33)
At fixed $`s`$ the only contribution of the integral over $`(t,u)[0,1]\times V`$ comes from the value $`\left(t=t_s,u=0\right)`$ since $`S_e^s`$ and $`e`$ intersect in the only point $`p_e^s`$ which is an interior point of $`[0,1]\times V`$ for sufficiently small $`ϵ`$. Here $`t_s`$ is defined by $`e_s\left(1/2\right)=e\left(t_s\right)`$. By definition of the orientation of the $`X_s\left(V\right)`$ we know that $`\dot{e}^a\left(t\right)ϵ_{ab_1..b_{D1}}X_{s,u_1}^{b_1}\left(u\right)..X_{s,u_{D1}}^{b_{D1}}\left(u\right)>0`$ at $`\left(t=t_s,u=0\right)`$. Since $`\rho _{e_s}\left(p_e^s\right)=p_e^s`$, (3.2.2) collapses to
$$\{P_i^e,h_e^{}\}_\gamma =\underset{ϵ0}{lim}\frac{1}{N}_1^1𝑑sg_ϵ\left(s\right)\text{tr}\left(\tau _ih_{e_s}(0,1/2)\tau _jh_{e_s}(0,1/2)^1\right)h_{e_s}(0,1/2)\frac{\tau _j}{2}h_{e_s}(1/2,1)$$
(3.34)
Since the integrand depends continuously on $`s`$ for any smooth connection we see that pointwise
$$\{P_i^e,h_e^{}\}_\gamma =\frac{1}{N}\text{tr}\left(\tau _ih_e(0,1/2)\tau _jh_e(0,1/2)^1\right)h_e(0,1/2)\frac{\tau _j}{2}h_e(1/2,1)$$
(3.35)
Now consider the matrix $`C_i=h_e(0,1/2)^1\tau _ih_e(0,1/2)`$ which is an element of the Lie algebra of $`G`$ because it is simply the transform of $`\tau _i`$ under the action of $`h_e(0,1/2)^1`$ in the adjoint representation of $`G`$ on $`Lie\left(G\right)`$. Thus we can expand $`C_i`$ in the basis $`\tau _j`$, resulting in $`C_i=\text{tr}\left(C_i\tau _j\right)\tau _j/N`$. Inserting this identity into (3.35) gives the result claimed.
\[3.\]
Let us now turn to the final third identity (3.26). It is clear that for $`ee^{}`$ and $`ϵ_1,ϵ_2`$ sufficiently small the regions $`R_e^{ϵ_1},R_e^{}^{ϵ_2}`$ are disjoint in which case the brackets vanish. Thus we only need to be concerned with the case $`e=e^{}`$. Let us again use the convention (3.31) and let us introduce the notation that for a function $`f`$ depending on $`(A,E)`$, the function $`\stackrel{~}{f}`$ is numerically equal to $`f`$ but $`\stackrel{~}{f}`$ is considered to be independent of $`(A,E)`$. In other words, $`\stackrel{~}{f}`$ drops out of Poisson brackets on $`M`$ but $`f`$ does not necessarily. Then we have by the Leibniz rule for $`\mathrm{\Omega }`$
$`\{P_i^{eϵ_1},P_j^{eϵ_2}\}=\{E\left(f_i^{ϵ_1}\right),E\left(f_j^{ϵ_2}\right)\}`$ (3.36)
$`=`$ $`\{E\left(\stackrel{~}{f}_i^{ϵ_1}\right),E\left(\stackrel{~}{f}_j^{ϵ_2}\right)\}+\{E\left(\stackrel{~}{f}_i^{ϵ_1}\right),\stackrel{~}{E}\left(f_j^{ϵ_2}\right)\}+\{\stackrel{~}{E}\left(f_i^{ϵ_1}\right),E\left(\stackrel{~}{f}_j^{ϵ_2}\right)\}+\{\stackrel{~}{E}\left(f_i^{ϵ_1}\right),\stackrel{~}{E}\left(f_j^{ϵ_2}\right)\}`$
The first term drops out by definition of $`\mathrm{\Omega }`$ (recall (3.13)) and the fourth by \[1.\] so that only the second and third term survive. More explicitly
$$\{P_i^{eϵ_1},P_j^{eϵ_2}\}=_\mathrm{\Sigma }d^3xE_k^a\left(x\right)\left[\{E\left(\stackrel{~}{f}_i^{ϵ_1}\right),\left(f_j^{ϵ_2}\right)_a^k\left(x\right))\right\}\{E\left(\stackrel{~}{f}_j^{ϵ_2}\right),\left(f_i^{ϵ_1}\right)_a^k\left(x\right))\}]$$
(3.37)
In order to compute (3.37) in the limit $`ϵ_1,ϵ_20`$ we have to consider two types of terms, namely $`\{E\left(\stackrel{~}{f}_i^ϵ\right),h_{\rho _{e_s}\left(X_s\left(u\right)\right)}\}`$ and $`\{E\left(\stackrel{~}{f}_i^ϵ\right),h_{e_s}(0,1/2)\}`$ The first term is given, according to Lemma 3.1, by
$`{\displaystyle _0^1}dt\dot{\rho }_{e_s}^a(X_s\left(u\right),t)\left(\stackrel{~}{f}_i^ϵ\right)_a^j(\rho _{e_s}\left(X_s\left(u\right)\right)h_{\rho _{e_s}\left(X_s\left(u\right)\right)}(0,t){\displaystyle \frac{\tau _j}{2}}h_{\rho _{e_s}\left(X_s\left(u\right)\right)}(t,1)`$ (3.38)
$`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle _0^1}dth_{\rho _{e_s}\left(X_s\left(u\right)\right)}(0,t){\displaystyle \frac{\tau _j}{2}}h_{\rho _{e_s}\left(X_s\left(u\right)\right)}(t,1){\displaystyle _1^1}drg_ϵ\left(r\right){\displaystyle _V}d^{D1}v\times `$
$`\times \delta (\rho _{e_s}(X_s\left(u\right),t),X_r\left(v\right))\dot{\rho }_{e_s}^a(X_s\left(u\right),t)ϵ_{ab_1..b_{D1}}X_{r,v_1}^{b_1}\left(v\right)..X_{r,v_{D1}}^{b_{D1}}\left(v\right)\times `$
$`\times \text{tr}\left(\tau _ih_{e_r}(0,1/2)h_{\rho _{e_r}\left(X_r\left(v\right)\right)}\tau _jh_{\rho _{e_r}\left(X_r\left(v\right)\right)}^1h_{e_r}(0,1/2)^1\right)`$
Using the fact that $`rX_r\left(V\right)`$ defines a foliation $`R_e`$ of surfaces diffeomorphic to $`S_e`$ we see that the integral over $`r,v`$ is supported at the interior point $`r=(s,v(s,t,u))`$ of $`[1,1]\times V`$ where $`X_s\left(v(s,t,u)\right)=\rho _{e_s}(X_s\left(u\right),t)`$. The integral can be performed with the result
$`{\displaystyle \frac{1}{N}}{\displaystyle _0^1}dth_{\rho _{e_s}\left(X_s\left(u\right)\right)}(0,t){\displaystyle \frac{\tau _j}{2}}h_{\rho _{e_s}\left(X_s\left(u\right)\right)}(t,1)g_ϵ\left(s\right)\times `$
$`\times [{\displaystyle \frac{\dot{\rho }_{e_s}^a(X_s\left(u\right),t)ϵ_{ab_1..b_{D1}}X_{r,v_1}^{b_1}\left(v\right)..X_{r,v_{D1}}^{b_{D1}}\left(v\right)}{|ϵ_{ab_1..b_{D1}}X_{r,r}^aX_{r,v_1}^{b_1}\left(v\right)..X_{r,v_{D1}}^{b_{D1}}\left(v\right)|}}\times `$
$`\times \text{tr}\left(\tau _ih_{e_r}(0,1/2)h_{\rho _{e_r}\left(X_r\left(v\right)\right)}\tau _jh_{\rho _{e_r}\left(X_r\left(v\right)\right)}^1h_{e_r}(0,1/2)^1\right)]_{r=s,v=v(s,t,u)}`$ (3.39)
Notice that the denominator in (3.2.2) is bounded away from zero as the curve $`sX_s\left(v\right)`$ for any fixed $`v`$ is transversal to $`X_s\left(V\right)`$. Now $`\rho _{e_s}(X_s\left(u\right),t)=\left(X_r\left(v\right)\right)_{r=s,v=v(s,t,u)}`$ for any $`t[0,1]`$, thus $`\dot{\rho }_{e_s}(X_s\left(u\right),t)=_{k=1}^{D1}\left(X_{r,v_k}\left(v\right)\right)_{r=s,v=v(s,t,u)}\frac{dv_k(s,t,u)}{dt}`$, thus the integrand of (3.2.2) vanishes for any finite $`ϵ`$.
Thus there will be no contribution from the holonomies along the $`\rho _{e_s}\left(X_s\left(u\right)\right)`$ and we can focus on the second term $`\{E\left(\stackrel{~}{f}_i^ϵ\right),h_{e_s}(0,1/2)\}`$ mentioned. Again, according to Lemma 3.1, it is given by
$`{\displaystyle _0^{1/2}}𝑑t\dot{e}_s^a\left(t\right)\left(\stackrel{~}{f}_i^ϵ\right)_a^j\left(e_s\left(t\right)\right)h_{e_s}(0,t){\displaystyle \frac{\tau _j}{2}}h_{e_s}(t,1/2)`$ (3.40)
$`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle _0^{1/2}}dth_{e_s}(0,t){\displaystyle \frac{\tau _j}{2}}h_{e_s}(t,1/2){\displaystyle _1^1}drg_ϵ\left(r\right){\displaystyle _V}d^{D1}v\delta (e_s\left(t\right),X_r\left(v\right))\dot{e}_s^a\left(t\right)\times `$
$`\times ϵ_{ab_1..b_{D1}}X_{r,v_1}^{b_1}\left(v\right)..X_{r,v_{D1}}^{b_{D1}}\left(v\right)\text{tr}\left(\tau _ih_{e_r}(0,1/2)h_{\rho _{e_r}\left(X_r\left(v\right)\right)}\tau _jh_{\rho _{e_r}\left(X_r\left(v\right)\right)}^1h_{e_r}(0,1/2)^1\right)`$
This time we will perform the $`t,v`$ integral to cancel the $`\delta `$-distribution. Thus we may as well perform the limits $`ϵ_1,ϵ_20`$ and cancel the $`s`$ and $`r`$ integrals first with the result that we are left with
$`{\displaystyle \frac{1}{N}}{\displaystyle _0^{1/2}}dth_e(0,t){\displaystyle \frac{\tau _j}{2}}h_e(t,1/2){\displaystyle _V}d^{D1}v\delta (e\left(t\right),X\left(v\right))\dot{e}^a\left(t\right)\times `$
$`\times ϵ_{ab_1..b_{D1}}X_{,v_1}^{b_1}\left(v\right)..X_{,v_{D1}}^{b_{D1}}\left(v\right)\text{tr}\left(\tau _ih_e(0,1/2)h_{\rho _e\left(X\left(v\right)\right)}\tau _jh_{\rho _e\left(X\left(v\right)\right)}^1h_e(0,1/2)^1\right)`$ (3.41)
The integrand is supported at $`t=1/2,v=0`$ which is now a boundary point of the interval $`[0,1/2]`$ which is why the $`\delta `$ distribution picks up a factor of $`1/2`$ as compared to the analogous argumentation in \[2.\]. This leads to the result
$$\frac{1}{2N}h_e(0,1/2)\frac{\tau _j}{2}\text{tr}\left(\tau _ih_e(0,1/2)\tau _jh_e(0,1/2)^1\right)=\frac{1}{4}\tau _ih_e(0,1/2)$$
Putting things together and using $`\delta g^1=g^1\delta gg^1`$ we find
$`\{P_i^e,P_j^e\}_\gamma `$ $`=`$ $`{\displaystyle _\mathrm{\Sigma }}d^3xE_k^a\left(x\right)[{\displaystyle \frac{1}{N}}{\displaystyle _V}d^{D1}u\delta (x,X\left(u\right))ϵ_{ab_1..b_{D1}}X_{u_1}^{b_1}\left(u\right)..X_{u_{D1}}^{b_{D1}}\left(u\right)]\times `$ (3.42)
$`\times `$ $`[\text{tr}\left(\tau _j\left\{{\displaystyle \frac{1}{4}}\tau _ih_e(0,1/2)\right\}h_{\rho _e\left(X\left(u\right)\right)}\tau _kh_{\rho _e\left(X\left(u\right)\right)}^1h_e(0,1/2)^1\right)`$
$`\text{tr}\left(\tau _jh_e(0,1/2)h_{\rho _e\left(X\left(u\right)\right)}\tau _kh_{\rho _e\left(X\left(u\right)\right)}^1h_e(0,1/2)^1\left\{{\displaystyle \frac{1}{4}}\tau _ih_e(0,1/2)\right\}h_e(0,1/2)^1\right)`$
$`ij]`$
$`=`$ $`{\displaystyle \frac{1}{4N}}{\displaystyle _{S_e}}(E)_k\left(x\right)\times `$
$`\times `$ $`[\text{tr}\left(\tau _j\tau _ih_e(0,1/2)h_{\rho _e\left(x\right)}\tau _kh_{\rho _e\left(x\right)}^1h_e(0,1/2)^1\right)`$
$`\text{tr}\left(\tau _jh_e(0,1/2)h_{\rho _e\left(x\right)}\tau _kh_{\rho _e\left(x\right)}^1h_e(0,1/2)^1\tau _i\right)`$
$`ij]`$
$`=`$ $`{\displaystyle \frac{1}{4N}}{\displaystyle _{S_e}}\left[2\text{tr}\left([\tau _j,\tau _i]h_e(0,1/2)h_{\rho _e\left(x\right)}E\left(x\right)h_{\rho _e\left(x\right)}^1h_e(0,1/2)^1\right)\right]`$
$`=`$ $`{\displaystyle \frac{f_{ij}^k}{N}}{\displaystyle _{S_e}}\text{tr}\left(\tau _kh_e(0,1/2)h_{\rho _e\left(x\right)}E\left(x\right)h_{\rho _e\left(x\right)}^1h_e(0,1/2)^1\right)`$
$`=`$ $`f_{ij}^kP_k^e`$
as claimed.
$`\mathrm{}`$
Notice that the factor of $`1/2`$ that appeared in (3.42) is also required if the bracket $`\{.,.\}_\gamma `$ is to satisfy the Leibniz rule : $`\{P_\gamma ^e,h_e\}_\gamma =\{P_\gamma ^e,h_e(0,1/2)\}h_e(1/2,1)+h_e(0,1/2)\{P_\gamma ^e,h_e(1/2,1)\}`$ which is consistent with (3.25) if indeed $`\{P_\gamma ^e,h_e(0,1/2)\}=\tau _ih_e(0,1/2)/4,\{P_\gamma ^e,h_e(1/2,1)\}=h_e(0,1/2)^1\tau _ih_e/4`$.
###### Theorem 3.3
The bracket $`\{.,.\}_\gamma `$ satisfies the Jacobi identity, moreover, it defines a non-degenerate two-form on $`M_\gamma `$, that is, it is a symplectic structure.
Proof of Theorem 3.3 :
i) Jacobi identity :
There are four kinds of double-brackets to check corresponding to $`n`$ momenta and $`3n`$ holonomies appearing with $`n=0,1,2,3`$.
$`n=0)`$ :
This case is trivial
$$\{\left(h_e\right)_{AB},\{\left(h_e^{}\right)_{CD},\left(h_{e^{\prime \prime }}\right)_{EF}\}_\gamma \}_\gamma +\text{cyclic}=0$$
(3.43)
since already the inner bracket vanishes by (3.24).
$`n=1)`$ :
Also this case is trivial
$$\{\left(h_e\right)_{AB},\{\left(h_e^{}\right)_{CD},P_i^{e^{\prime \prime }}\}_\gamma \}_\gamma +\text{cyclic}=0$$
(3.44)
because either the inner bracket already vanishes or the inner bracket gives a function depending only on holonomies by (3.25) and so the outer bracket is of the type (3.24) and vanishes.
$`n=2)`$ :
This is the first non-trivial case and it is quite remarkable that the signs and numerical factors in (3.24)-(3.26) come in precisely the right way out of the regularized derivation of theorem 3.2 :
$`\{h_e,\{P_i^e^{},P_j^{e^{\prime \prime }}\}_\gamma \}_\gamma +\{P_i^e^{},\{P_j^{e^{\prime \prime }},h_e\}_\gamma \}_\gamma +\{P_j^{e^{\prime \prime }},\{h_e,P_i^e^{}\}_\gamma \}_\gamma `$ (3.45)
$`=`$ $`\delta ^{e^{}e^{\prime \prime }}f_{ij}^k\{h_e,P_k^e^{}\}_\gamma +\delta _e^{e^{\prime \prime }}\{P_i^e^{},{\displaystyle \frac{\tau _j}{2}}h_e\}_\gamma \delta _e^e^{}\{P_j^{e^{\prime \prime }},{\displaystyle \frac{\tau _i}{2}}h_e\}_\gamma `$
$`=`$ $`\delta ^{e^{}e^{\prime \prime }}\delta _e^e^{}\left(f_{ij}^k{\displaystyle \frac{\tau _k}{2}}h_e+{\displaystyle \frac{\tau _j\tau _i}{4}}h_e{\displaystyle \frac{\tau _i\tau _j}{4}}h_e\right)`$
$`=`$ $`{\displaystyle \frac{1}{4}}\delta ^{e^{}e^{\prime \prime }}\delta ^{e^{}e}\left(2f_{ij}^k\tau _k+[\tau _j,\tau _i]\right)h_e=0`$
by definition of the structure constants.
$`n=3`$ :
This case is again easy, it just relies on the Jacobi identity for the generators of the Lie algebra of $`G`$ or, equivalently, for its structure constants :
$$\{P_i^e,\{P_j^e^{},P_k^{e^{\prime \prime }}\}_\gamma \}_\gamma +\text{cyclic}=\delta ^{e^{}e^{\prime \prime }}\delta _e^e^{}\left(f_{jk}^lf_{il}^m+\text{cyclic}\right)P_m^e=0$$
(3.46)
ii) Non-degeneracy
Obviously, by (3.24)-(3.26) the symplectic structure $`\mathrm{\Omega }_\gamma `$, if it exists, is diagonal with respect to the edge label,
$$\mathrm{\Omega }_\gamma =\underset{eE\left(\gamma \right)}{}\mathrm{\Omega }_e$$
(3.47)
where each $`\mathrm{\Omega }_e`$ is isomorphic with a Poisson structure $`\mathrm{\Omega }_G`$ on the cotangent bundle $`T^{}G`$ given by $`\{h_{AB},h_{CD}\}_G=0,\{P_i,h\}_G=\tau _ih/2,\{P_i,P_j\}_G=f_{ij}^kP_k`$. Thus in order to show regularity of $`\mathrm{\Omega }_\gamma `$ it will be necessary to show regularity of $`\mathrm{\Omega }_G`$, that is, $`\mathrm{\Omega }_G`$ is not only a Poisson structure but actually a symplectic structure.
To see that $`\mathrm{\Omega }_G`$ is everywhere nondegenerate on $`G`$ consider the atlas $`(U_\alpha ,\varphi _\alpha )_{\alpha I}`$ of $`G`$ given by open neighbourhoods $`U_\alpha `$ containing some point $`h_\alpha G`$ and charts defined by $`\varphi _\alpha ^1:=\mathrm{exp}:V_\alpha \text{ }\mathrm{R}^{dim\left(G\right)}U_\alpha ;\left(\theta _\alpha ^j\right)\mathrm{exp}(\theta _\alpha ^j\tau _j/2)h_\alpha `$. Obviously, the $`\theta _\alpha ^j`$ serve as local coordinates of $`G`$. Let now $`hG`$ be given, then there exists $`\alpha I,\left(\theta _\alpha ^j\right)V_\alpha `$ such that $`h=\mathrm{exp}\left(\theta _\alpha ^j\tau _j/2\right)h_\alpha `$. By choosing the size of the index set $`I`$ high enough we can assume that the range of each $`\theta _\alpha ^j`$ is contained in an open interval containing zero as small as we wish. Let us now expand $`h`$ in powers of $`\theta _\alpha ^j`$ in the relation $`\{P_j,h\}_G=\tau _jh/2`$ then we find by comparing powers that $`\{P_j,\theta ^k\}=\delta _j^k+O\left(\theta \right)`$ where $`O\left(\theta \right)`$ is a bounded function vanishing linearly in $`\theta `$. We conclude that the bracket when expressed in terms of the coordinates $`P_j,\theta _j`$ has locally the form of a block matrix consisting of four blocks of $`dim\left(G\right)\times dim\left(G\right)`$ matrices with the off-diagonal blocks given by plus/minus the identity matrix plus corrections of order $`\theta _\alpha ^j`$ and the diagonal blocks consist of a zero matrix and of the matrix with $`ij`$-entries given by $`f_{ij}^kP_k`$. This matrix is obviously invertible on $`U_\alpha `$ for any $`\alpha `$, proving our claim.
Finally we must show that $`\mathrm{\Omega }_\gamma `$ is a surjection, that is, for any $`(x_e,y_e)_{eE\left(\gamma \right)}_\gamma ^{}`$ there exists $`(x_e^{},y_e^{})_{eE\left(\gamma \right)}_\gamma `$ such that $`(x_e^{},y_e^{})\mathrm{\Omega }_e=(x_e,y_e)`$. Since $`_\gamma `$ is actually a Hilbert-manifold, we find $`(\stackrel{~}{x}_e,\stackrel{~}{y}_e)_\gamma `$ such that $`(x_e,y_e)=(\stackrel{~}{x}_e\rho _e,\stackrel{~}{y}_e\sigma _e)`$ thus the solution of our problem is given by $`(x_e^{},y_e^{})=(\stackrel{~}{x}_e\rho _e,\stackrel{~}{y}_e\sigma _e)\left(\mathrm{\Omega }_e\right)^1`$ ($`\mathrm{\Omega }_e`$ interpreted as a matrix written pointwise in $`M_\gamma `$) provided we can show that this defines an element of $`_\gamma `$. This is somewhat non-trivial because $`\mathrm{\Omega }_e`$ depends on $`P_i^e`$ which, a priori, can take arbitrarily large values. However, the normalizability of this vector follows from (3.18) which implies that in fact $`P_i^e`$ must be uniformly bounded in $`e`$.
$`\mathrm{}`$
In the appendix we show that $`\mathrm{\Omega }_G`$ and therefore $`\mathrm{\Omega }_\gamma `$ is even exact in the case of $`G=U\left(1\right),SU\left(2\right)`$ which one can probably prove also for general $`G`$ and we leave this as a future project.
#### 3.2.3 Continuum Phase Space from the Graph Phase Space
In the previous subsection we established how the symplectic manifold $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ can be derived from the symplectic manifold $`(M,\mathrm{\Omega })`$ for every $`\gamma \mathrm{\Gamma }`$. In this subsection we wish to show the opposite : in the limit that $`\gamma `$ grows ad infinitum in a prescribed way we find that $`(M,\mathrm{\Omega })`$ can be derived from $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$. This lies at the heart of later constructions in which we use $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ as our starting point for quantization. Namely, as the formulation of the theory in terms of $`M_\gamma `$ is a certain kind of discretization, the result just stated means that the continuum limit exists and is the expected one on the classical level. On the other hand, the symplectic manifold $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ is straightforward to quantize on a Hilbert space $`(_\gamma ,<.,.>_\gamma )`$ and the classical limit of this quantization is easily shown to give back $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$. In other words, we can establish the chain of limits $`(_\gamma ,<.,.>_\gamma )_\mathrm{}0(M_\gamma ,\mathrm{\Omega }_\gamma )_\gamma \mathrm{}(M,\mathrm{\Omega })`$ and the interesting question of the existence of the opposite limit will be subject of our companion paper .
In order to show that we can recover the continuum theory from the discrete one in the limit of infinite graphs we consider a certain one-parameter family of cubic lattices $`ϵ\gamma _ϵ`$ where $`ϵ`$ is associated with the length of the edges or links of the graph with respect to a certain background metric and the limit $`ϵ0`$ corresponds to sending the graph to the continuum $`\mathrm{\Sigma }`$. More precisely, we have the following :
The manifold $`\mathrm{\Sigma }`$ is described by an atlas of charts $`(U_\iota ,X_\iota )_\iota `$ where $`U_\iota `$ is an open region in $`\mathrm{\Sigma }`$ and $`X_\iota ^a:V_\iota \text{ }\mathrm{R}^DU_\iota ,(t^1,..,t^D)X_\iota \left(t\right)`$ is a local trivialization of $`U_\iota `$, that is, a smooth orientation preserving diffeomorphism. Consider an arbitrary but fixed decomposition $``$ of $`\mathrm{\Sigma }`$ into mutually disjoint, except for common boundary points, compact regions $`R`$ which is fine enough such that every $`R`$ lies in the domain of a chart and choose $`\iota \left(R\right)`$ to be such that $`RU_{\iota \left(R\right)}`$ (at this point we do not even need the cover $`\left\{U_i\right\}`$ to be locally finite or $`\mathrm{\Sigma }`$ to be paracompact although this is an assumption which goes into the definition of $`\mathrm{\Gamma }_\sigma ^\omega `$). Without loss of generality we can assume that each $`R`$ is diffeomorphic to a polyhedron of $`\text{ }\mathrm{R}^D`$ and we fix for each of its faces an orientation. Also, if $`\mathrm{\Sigma }`$ is not compact we take a refinement of the atlas if necessary such that each $`R`$ has finite Lebesgue measure.
Now $`R`$ is the image under $`X_R:=X_{\iota \left(R\right)}`$ of a compact region $`V_R`$, in fact a polyhedron, in $`\text{ }\mathrm{R}^D`$ which can always be decomposed, for sufficiently small $`ϵ`$, into regular cubes of volume $`ϵ^D`$ with respect to the Euclidean metric of $`\text{ }\mathrm{R}^D`$ possibly up to a subset near the the boundary of $`X_R^1\left(R\right)`$. Let then $`C_R`$ be the union of those cubes which fit into $`V_R`$.
Now each region $`C_R`$ is filled exactly with cubes of volume $`ϵ^D`$ and these cubes define a regular, oriented cubic lattice $`\gamma _R^0`$ in $`C_R`$, the orientation of the edges is chosen to be such that each of them points in the positive coordinate axis direction of $`\text{ }\mathrm{R}^D`$. Let $`P_{\gamma _R^0}^0`$ be the dual decomposition of $`V_R`$ obtained as follows : choose for each edge $`e_R^0`$ of $`\gamma _R^0`$ the open face $`S_{e_R^0}^0`$ to be the regular cubic hyperplane orthogonal to $`e_R^0`$ of area $`ϵ^{D1}`$ which cuts $`e_R^0`$ in the middle, is pierced by $`e_R^0`$ in its center and carries the orientation defined by choosing the tangent direction of $`e_R^0`$ to be the direction of its unit normal vector. The collection of all these faces $`S_{e_R^0}^0`$ can be completed uniquely to the unique, minimal cubic polyhedronal complex $`C_R^{}`$ which contains all of them. We assume w.l.g. that $`C_R^{}`$ fits into $`V_R`$ (decrease the size of $`C_R`$ by deleting some cubes on its boundary if necessary). Notice that $`C_R^{}`$ necessarily covers $`C_R`$. We complete the polyhedronal decomposition of $`V_R`$ by choosing $`X_R^1\left(R\right)C_R^{}`$ as the final polyhedron to cover $`V_R`$. Since the boundary of $`V_R`$ is already oriented this obviously defines an oriented polyhedronal decomposition of $`X_R^1\left(R\right)`$ dual to $`\gamma _R^0`$.
Finally, we consider the images $`\gamma _R=X_{\iota \left(R\right)}\left(\gamma _R^0\right),S_e=X_{\iota \left(R\right)}\left(S_{e_R^0}^0\right)P_{\gamma _R}=X_{\iota \left(R\right)}\left(P_{\gamma _R^0}^0\right)`$ and the corresponding variables $`h_e^R,P_i^{e,R}`$ for each $`R`$. The unions $`\gamma _ϵ:=_R\gamma _R`$ and $`P_{\gamma _ϵ}:=_RP_{\gamma _R}^0`$ define an oriented graph in $`\mathrm{\Sigma }`$ and an oriented decomposition of $`\mathrm{\Sigma }`$. It is not yet a decomposition dual to $`\gamma `$ since it is not minimal : it becomes minimal if we remove all the boundaries $`R`$. Let the resulting dual decomposition also be denoted by $`P_{\gamma _ϵ}`$.
Notice that the region $`\mathrm{\Sigma }\left[_RX_R\left(C_R^{}\right)\right]`$ does not contain any piece of $`\gamma _ϵ`$, however, its Lebesgue measure vanishes in the limit $`ϵ0`$ because it tends to $`_RR`$. We could avoid this by adding edges to $`\gamma `$ connecting the $`\gamma _R`$ which are contained in neigbouring $`R`$ but the resulting lattice may not be of cubic topology any longer. Since we will not need those edges for the sake of our argument, we will leave things as they are.
Let us fix a specific $`R`$ and define $`v_R:=X_R\left(v_R^0\right),e_{RI}\left(v\right):=X_R\left(e_{RI}^0\left(v^0\right)\right),S_R^I\left(v\right):=X_R\left(S_{e_{RI}^0\left(v_R^0\right)}^0\right)`$ respectively to be the image under $`X^R`$ of a vertex, edge and face of $`\gamma _R^0`$ respectively. Here $`e_{RI}^0\left(v_R^0\right)`$ denotes the straight line into the positive $`I`$direction between the vertices $`v_R^0`$ and $`v_R^0+ϵb_I`$ where $`\left\{b_I\right\}_{I=1}^D`$ denotes the standard oriented orthonormal basis of $`\text{ }\mathrm{R}^D`$ (sometimes $`v_R^0+ϵb_I`$ is not a vertex of $`\gamma _R^0`$ in which case we set $`e_{RI}^0\left(v_R^0\right)=v_R^0=S_{e_{RI}^0\left(v_R^0\right)}^0`$). Consider also for any $`xR`$ the vector fields $`Y_{RI}^a\left(x\right):=X_{R,I}^a\left(t\right)_{x=X_R\left(t\right)}`$ and co-vector densities of weight minus one $`n_{Ra}^I\left(x\right):=\frac{1}{\left(D1\right)!}ϵ_{ab_1..b_{D1}}ϵ^{IJ_1..J_{D1}}Y_{RJ_1}^{b_1}\left(x\right)..Y_{RJ_{D1}}^{b_{D1}}\left(x\right)`$. Since $`X_R:V_RR`$ is an orientation preserving diffeomorphism it is clear that $`det\left(\left(Y\right)\right)=n_a^IY_I^a>0`$ everywhere in $`R`$.
We define now as in (3.15) for every vertex $`v`$ of $`\gamma _R`$ the functions (we drop the label $`R`$)
$`h_I\left(v\right)`$ $`:=`$ $`h_{e_I\left(v\right)}\left(A\right)`$ (3.48)
$`P_i^I\left(v\right)`$ $`:=`$ $`{\displaystyle \frac{1}{N}}\text{tr}\left(\tau _ih_{e_I\left(v\right)}(0,1/2)\left[{\displaystyle _{S^I\left(v\right)}}h_{\rho _{e_I\left(v\right)}\left(x\right)}E\left(x\right)h_{\rho _{e_I\left(v\right)}\left(x\right)}^1\right]h_{e_I\left(v\right)}(0,1/2)^1\right)(A,E)`$
which defines a map $`D_ϵ:MM_\gamma `$.
On the other hand, consider now the following functions on $`M_\gamma `$ ($`v`$ is again a vertex of $`\gamma _R`$ and we drop the label $`R`$)
$`A_a^{\left(ϵ\right)i}\left(v\right)`$ $`:=`$ $`2{\displaystyle \frac{n_a^I}{det\left(Y\right)Nϵ}}\text{tr}\left(\tau _ih_I\left(v\right)\right)`$
$`E_i^{\left(ϵ\right)a}\left(v\right)`$ $`:=`$ $`{\displaystyle \frac{Y_I^a}{det\left(Y\right)ϵ^{D1}}}P_i^I\left(v\right)`$ (3.49)
Suppose first that $`(h_I\left(v\right),P_i^I\left(v\right))G\times Lie\left(G\right)`$ are obtained via the map $`D_ϵ`$. Then, using the smoothness of the fields $`(A,E)`$ it is easy to see that $`A_a^{\left(ϵ\right)i}\left(v\right)A_a^i\left(v\right)`$ and $`E_i^{\left(ϵ\right)a}\left(v\right)E_i^a\left(v\right)`$ are both of order $`ϵ`$.
We now select from the one-parameter family of lattices $`\gamma _ϵ`$ and decompositions $`P_{\gamma _ϵ}`$ a sequence of lattices $`\gamma _n`$ with the property that $`\gamma _{n+1}`$ and $`P_{\gamma _{n+1}}`$ respectively are refinements of $`\gamma _n`$ and $`P_{\gamma _n}`$ respectively. It is easy to see that the following sequence does the job : Start with some $`ϵ_0>0`$ and call $`\gamma _0:=\gamma _{ϵ_0},P_0:=P_{\gamma _{ϵ_0}}`$ respectively. Now consider the sequence $`ϵ_n:=ϵ_0/3^n`$. The corresponding $`\gamma _n:=\gamma _{ϵ_n},P_n:=P_{\gamma _{ϵ_n}}`$ are obtained iteratively as follows :
Given $`\gamma _{Rn}^0`$, subdivide each of the cubes it defines into $`3^D`$ axis parallel cubes of equal volume $`\left(ϵ_n/3\right)^D`$ and similarly for $`P_{\gamma _{Rn}}^0`$. Both of these lattices obviously refine the previous ones respectively.
To see that the refinement of $`P_{\gamma _{Rn}^0}^0`$ has all the properties that we required it to have in the above construction of a decomposition of $`V_R`$ dual to the refinement of $`\gamma _{Rn}^0`$ we remark the following :
If we consider $`\gamma _{Rn}^0`$ as a sublattice of an infinite regular cubic lattice $`L_n`$ in $`\text{ }\mathrm{R}^D`$ then $`P_{\gamma _{Rn}^0}^0`$ satisfies the required proprties if and only if its restriction to $`C_R^{}`$ is defined by a sublattice of the lattice $`L_n^{}`$ obtained from $`L_n`$ by translating it by the vector $`ϵ_n_{I=1}^Db_I/2`$. In other words, if we label the vertices of $`L_n`$ by the n-tuples $`\left(ϵ_nn_I\right)_{I=1}^D,n_IZ`$ then the vertices of $`L_n^{}`$ are labelled by the n-tuples $`\left(ϵ_n\left[n_I^{}+1/2\right]\right)_{I=1}^D,n_I^{}Z`$. Now the points of the refinements of $`L_n`$ and $`L_n^{}`$ respectively are labelled by $`\left(ϵ_nn_I/3\right)_I=ϵ_{n+1}n_I`$ and $`\left(ϵ_n\left[n_I^{}/3+1/2\right]\right)_I=\left(ϵ_n\left[\left(n_I^{}+1\right)/3+1/6\right]\right)_I=\left(ϵ_{n+1}\left[n_I^{\prime \prime }+1/2\right]\right)_I`$, in other words, the refinements coincide with $`L_{n+1}`$ and $`L_{n+1}^{}`$ respectively. Notice that our procedure would work also if we would choose to refine by $`k`$ instead of $`3`$ where $`k>3`$ can be any odd integer.
To complete the decomposition we add cubes to these refinements as to fill $`V_R`$ as densely as possible according to the rules we specified above (also deleting cubes if necessary as discussed above) and thus obtain, after mapping with $`X_R`$, $`\gamma _{n+1}`$ and $`P_{\gamma _{n+1}}`$.
We remark without proof that only cubic lattices seem to have the property that there are refinements such that a dual decomposition of the refinement can be a refinement of the dual decomposition (consider a simplicial decomposition to see the arising problems).
Now that we know that if $`v`$ is a vertex of $`\gamma _n`$ for some $`n`$ then it is a vertex of all $`\gamma _m`$ for all $`mn`$, the limit $`ϵ0`$ is well defined and we have, provided that $`A_a^{\left(ϵ_n\right)i}\left(v\right)=:A_a^{\left(n\right)i}\left(v\right),E_i^{\left(ϵ_n\right)a}\left(v\right)=:E_i^{\left(n\right)a}\left(v\right)`$ are defined via $`D_n:=D_{ϵ_n}`$
$`\underset{ϵ0}{lim}\left[A_a^{\left(ϵ\right)i}\left(v\right)A_a^i\left(v\right)\right]:=\underset{n\mathrm{}}{lim}\left[A_a^{\left(n\right)i}\left(v\right)A_a^i\left(v\right)\right]`$
$`\underset{ϵ0}{lim}\left[E_i^{\left(ϵ\right)a}\left(v\right)E_i^a\left(v\right)\right]:=\underset{n\mathrm{}}{lim}\left[E_a^{\left(n\right)a}\left(v\right)E_i^a\left(v\right)\right]`$ (3.50)
where convergence is pointwise on $`M`$. In other words, the map $`D_n:MM_nM_{\gamma _n}`$ is invertible in the limit $`n\mathrm{}`$. To see that also the symplectic structure $`\mathrm{\Omega }`$ of $`M`$ is recovered in this limit we notice first that for each $`f_a^i,F_i^a𝒮`$ we have $`F\left(A\right)=lim_n\mathrm{}FA^{\left(n\right)}\left(A\right),E\left(f\right)=lim_n\mathrm{}E^{\left(n\right)}(A,E)f`$ pointwise in $`M`$ where
$`FA^{\left(n\right)}`$ $`:=`$ $`{\displaystyle \underset{vV\left(\gamma _n\right)}{}}ϵ_n^D\left(det\left(Y\right)\right)\left(v\right)F_a^i\left(v\right)A_a^{\left(n\right)i}\left(v\right)`$
$`E^{\left(n\right)}f`$ $`:=`$ $`{\displaystyle \underset{vV\left(\gamma _n\right)}{}}ϵ_n^D\left(det\left(Y\right)\right)\left(v\right)E_i^{\left(n\right)a}\left(v\right)f_a^i\left(v\right)`$ (3.51)
###### Theorem 3.4
The bracket defined by
$`\{F\left(A\right),F^{}\left(A\right)\}^{}`$ $`:=`$ $`\underset{n\mathrm{}}{lim}\{FA^{\left(n\right)},F^{}A^{\left(n\right)}\}_{\gamma _n}`$ (3.52)
$`\{E\left(f\right),F\left(A\right)\}^{}`$ $`:=`$ $`\underset{n\mathrm{}}{lim}\{E^{\left(n\right)}f,FA^{\left(n\right)}\}_{\gamma _n}`$ (3.53)
$`\{E\left(f\right),E\left(f^{}\right)\}^{}`$ $`:=`$ $`\underset{n\mathrm{}}{lim}\{E^{\left(n\right)}f,E^{\left(n\right)}f^{}\}_{\gamma _n}`$ (3.54)
for all $`f_a^i,F_i^a,f_a^i,F_i^a𝒮`$ coincides with the symplectic structure $`\mathrm{\Omega }`$ on $`M`$ defined by (3.13).
Proof of Theorem 3.4 :
We simply have to use the symplectic structure of $`(M_{\gamma _n},\mathrm{\Omega }_{\gamma _n})`$ and take the limit using (3.2.3). In the notation of this subsection, the symplectic structure labelled by $`\gamma _n`$ can be written
$`\{\left(h_I\left(v\right)\right)_{AB},\left(h_J\left(v^{}\right)\right)_{CD}\}_{\gamma _n}=0`$
$`\{P_i^I\left(v\right),h_J\left(v^{}\right)\}_{\gamma _n}=\delta _J^I\delta _{v,v^{}}{\displaystyle \frac{\tau _i}{2}}h_J\left(v\right)`$
$`\{P_i^I\left(v\right),P_j^J\left(v^{}\right)\}_{\gamma _n}=\delta ^{IJ}\delta _{vv^{}}f_{ij}^kP_k^I\left(v\right)`$ (3.55)
\[1.\]
Then (3.52) is obvious since the right hand side vanishes already at any finite $`n`$.
\[2.\]
We have at fixed $`n`$ for the right hand side of (3.53)
$`\{E^{\left(n\right)}f,FA^{\left(n\right)}\}_{\gamma _n}={\displaystyle \underset{v,v^{}V\left(\gamma _n\right)}{}}ϵ_n^{2D}(det\left(Y\right))\left(v\right)(det\left(Y\right))\left(v^{}\right)f_a^i\left(v\right)F_j^b\left(v^{}\right)\times `$ (3.56)
$`\times \left[2{\displaystyle \frac{n_b^J\left(v^{}\right)}{det\left(Y\right)\left(v^{}\right)Nϵ_n}}{\displaystyle \frac{Y_I^a\left(v\right)}{det\left(Y\right)\left(v\right)ϵ_n^{D1}}}\right]\{P_i^I\left(v\right),\text{tr}\left(\tau _jh_J\left(v^{}\right)\right)\}_{\gamma _n}`$
$`=`$ $`{\displaystyle \frac{2}{N}}{\displaystyle \underset{vV\left(\gamma _n\right)}{}}ϵ_n^Df_a^i\left(v\right)F_j^b\left(v\right)\left[{\displaystyle \underset{I}{}}\left(n_b^IY_I^a\right)\left(v\right)\text{tr}\left(\tau _j{\displaystyle \frac{\tau _i}{2}}h_I\left(v\right)\right)\right]`$
$`=`$ $`{\displaystyle \frac{2}{N}}{\displaystyle \underset{vV\left(\gamma _n\right)}{}}ϵ_n^Df_a^i\left(v\right)F_j^b\left(v\right)\left[{\displaystyle \frac{N}{2}}\delta _{ij}\delta _b^adet\left(Y\right)\left(v\right)+{\displaystyle \underset{I}{}}\left(n_b^IY_I^a\right)\left(v\right)\text{tr}\left(\tau _j{\displaystyle \frac{\tau _i}{2}}\left[h_I\left(v\right)1\right]\right)\right]`$
$`=`$ $`\left\{{\displaystyle \underset{vV\left(\gamma _n\right)}{}}ϵ_n^D\left[det\left(Y\right)f_a^iF_i^a\right]\left(v\right)\right\}`$
$`{\displaystyle \frac{2}{N}}\left\{{\displaystyle \underset{vV\left(\gamma _n\right)}{}}ϵ^D\left[f_a^iF_j^b\right]\left(v\right){\displaystyle \underset{I}{}}\left(n_b^IY_I^a\right)\left(v\right)\text{tr}\left(\tau _j{\displaystyle \frac{\tau _i}{2}}\left[h_I\left(v\right)1\right]\right)\right\}`$
Consider the first term in the last equality of (3.56). We can write it as
$$\underset{R}{}\underset{v^0\gamma _R^0}{}ϵ_n^D|det(\left(\frac{X}{t}_{t\left(v^0\right)}\right|\left[f_a^iF_i^a\right]\left(X\left(t\left(v^0\right)\right)\right)$$
(3.57)
which defines a Riemann sum for the expression
$$\underset{R}{}_{C_R}d^Dt|(det\left(\left(\frac{X}{t}\right)_{|t}\right|\left[f_a^iF_i^a\right]\left(X\left(t\right)\right)=F\left(f\right)\underset{R}{}_{RX_R\left(C_R\right)}d^Dx\left[f_a^iF_i^a\right]\left(x\right)$$
(3.58)
and the second integral in (3.58) vanishes in the limit $`n\mathrm{}`$ (that is $`C_RR`$). To see this, notice that in the limit $`n\mathrm{}`$ the integral $`_{RX_R\left(C_R\right)}d^Dx\left[f_a^iF_i^a\right]\left(x\right)`$ becomes $`_Rd^Dx\left[f_a^iF_i^a\right]\left(x\right)`$ which vanishes for non-distributional spaces of test functions.
Turning to the second term in (3.56) we notice that there exists a positive constant $`k`$ such that $`\left|\text{tr}\left(\tau _j\frac{\tau _i}{2}[h_I\left(v\right)1]\right)\right]|kϵ_n`$ as $`n\mathrm{}`$ independent of the indices and $`v`$ because $`A`$ is a smooth and bounded function. Thus we see that the second term vanishes in the limit $`n\mathrm{}`$.
\[3.\]
Finally, at fixed $`n`$ the right hand side of (3.54) becomes
$`\{E^{\left(n\right)}f,E^{\left(n\right)}f^{}\}_{\gamma _n}`$ $`=`$ $`{\displaystyle \underset{v,v^{}V\left(\gamma _n\right)}{}}ϵ_n^2\left[f_a^iY_I^a\right]\left(v\right)\left[f_b^jY_J^b\right]\left(v^{}\right)\{P_i^I\left(v\right),P_j^J\left(v^{}\right)\}_{\gamma _n}`$
$`=`$ $`f_{ij}^k{\displaystyle \underset{vV\left(\gamma _n\right)}{}}ϵ_n^2\left[{\displaystyle \underset{I}{}}\left(f_a^iY_I^af_b^jY_I^bP_k^I\right)\left(v\right)\right]`$
$`=`$ $`f_{ij}^kϵ_n{\displaystyle \underset{vV\left(\gamma _n\right)}{}}ϵ_n^D\left[{\displaystyle \underset{I}{}}\left(f_a^iY_I^af_b^jY_I^b\left[E_i^cn_c^I+\left\{{\displaystyle \frac{P_k^I}{ϵ_n^{D1}}}E_i^cn_c^I\right\}\right]\right)\left(v\right)\right]`$
The whole sum is just an approximation for a Riemann integral times $`ϵ_n`$. The term in the curly bracket approaches zero as $`n0`$ and is therefore, together with the first term in the square bracket, integrable against the product of test functions displayed. Thus, the whole expression vanishes in the limit $`n\mathrm{}`$.
$`\mathrm{}`$
### 3.3 Structured Graphs as Labels for Generalized Projective Families
A natural question to ask is whether one can identify $`(M,\mathrm{\Omega })`$ with (the limit of) a generalized projective sequence of symplectic manifolds $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$. The answer is affirmative but somewhat involved because we first must introduce new labels for projective families :
First of all, the family $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ does not only depend on the graph $`\gamma `$ but actually on the set $``$ of structured graphs $`l=(\gamma ,P_\gamma ,\mathrm{\Pi }_\gamma )`$ consisting of a graph $`\gamma `$, a polyhedronal decomposition $`P_\gamma `$ dual to it and a choice of paths $`\rho _e\left(x\right)\mathrm{\Pi }_\gamma `$ adapted to $`\gamma ,P_\gamma `$ where $`\rho _e\left(x\right)S_e,eE\left(\gamma \right),xS_e`$. The family $``$ is partially ordered by inclusion but it is in general wrong that given two elements $`l,l^{}`$ there exists a common refinement, that is, an element $`\stackrel{~}{l}`$ such that $`l,l^{}\stackrel{~}{l}`$. In other words, the inclusion relation does not equip $``$ with the structure of a directed set on which the structure of a generalized projective limit crucially depends.
In order to proceed, we therefore must first modify the partial order. To motivate our choice we begin with the following observation :
Given a graph $`\gamma `$ the second and third entry of a structured graph $`l`$ such that $`\gamma \left(l\right)=\gamma `$ are largely arbitrary. On the other hand, if we consider structured graphs $`l,l^{}`$ with $`\gamma \left(l\right)=\gamma \left(l^{}\right)`$ then by construction we can always find a diffeomorphism that preserves $`\gamma \left(l\right)`$ and maps $`P_\gamma ,\mathrm{\Pi }_\gamma `$ to $`P_\gamma ^{},\mathrm{\Pi }_\gamma ^{}`$. This follows from the fact that all the $`S_e,\rho _e\left(x\right)S_e,xS_e`$ are obtained via a diffeomorphism from a universal object by definition 3.5. Now, while the actual values of the $`P^e`$ that we construct from $`l`$ or $`l^{}`$ respectively may differ (the $`h_e`$ are evidently the same), the Poisson algebras, that is to say the algebra of Hamiltonian vector fields, that we obtain are identical. Moreover, let us consider the $`h_e`$ as elements of the space of smooth functions $`C^{\mathrm{}}\left(𝒞_\gamma \right)`$ of the configuration space $`𝒞_\gamma `$ of $`M_\gamma `$ (in fact they are coordinate functions) and the $`P^e`$ as elements of the space of vector fields $`V\left(𝒞_\gamma \right)`$ on $`𝒞_\gamma `$ via the map $`(h_e,P_j^e)(h_e,\text{tr}(\left(\tau _jh_e\right)^T/h_e)`$. The space $`C^{\mathrm{}}\left(𝒞_\gamma \right)\times V\left(𝒞_\gamma \right)`$ is equipped with the Lie algebra structure $`[(f,u),(f^{},u^{})]=(u\left(f^{}\right)u^{}\left(f\right),[u,u^{}])`$ which is evidently closed and isomorphic with the Poisson bracket structure as obtained from both $`l,l^{}`$. In this form it is particularly obvious that both $`l,l^{}`$ give rise to the same Lie algebra.
What this means is that the information contained in $`l`$ beyond that of $`\gamma \left(l\right)`$ is irrelevant as far as the Poisson algebra is concerned. Since it is the Poisson structure which sets the correspondence with quantum theory we will obtain isomorphic quantum theories from both $`l,l^{}`$. The additional information contained in $`l`$ however comes in when we consider the classical limit of the theory. Namely, the coherent states constructed in are sensitive to the size and shape of the $`S_e`$ as well as the precise choice of the paths $`\rho _e`$.
These considerations shed light on the question why we have largely abused notation when writing $`\gamma `$ instead of $`l`$ and is also reflected in the subsequent definition.
###### Definition 3.7
We say that $`(\gamma ,P_\gamma ,\mathrm{\Pi }_\gamma )(\gamma ^{},P_\gamma ^{},\mathrm{\Pi }_\gamma ^{})`$ provided that $`\gamma (l)\gamma (l^{})`$ and that $`l,l^{}`$ are equivalent, $`ll^{}`$, if $`\gamma (l)=\gamma (l^{})`$.
In other words, there are no conditions at all on the second and third entry of a structured graph. In particular, we identify $`l`$ with $`l^{}`$ if one obtains $`l^{}`$ form $`l`$ by applying a diffemorphism that preserves $`\gamma \left(l\right)`$. The subsequent two lemmas are then almost trivial.
###### Lemma 3.2
The relation $``$ defined in definition 3.7 defines a partial order.
Proof of Lemma 3.2 :
1) Reflexivity : $`ll`$ since $`\gamma \left(l\right)=\gamma \left(l\right)`$.
2) Antisymmetry : $`ll^{}`$ and $`l^{}l`$ implies $`\gamma \left(l\right)=\gamma \left(l^{}\right)`$, that is, $`ll^{}`$.
3) Transitivity : $`ll^{}`$ and $`l^{}l^{\prime \prime }`$ implies $`\gamma \left(l\right)\gamma \left(l^{}\right)\gamma \left(l^{\prime \prime }\right)`$, thus $`ll^{\prime \prime }`$.
$`\mathrm{}`$
###### Lemma 3.3
The set $``$ is directed, that is, for any given $`l,l^{}`$ there exists $`\stackrel{~}{l}`$ such that $`l,l^{}\stackrel{~}{l}`$. Such an element $`\stackrel{~}{l}`$ is called a common refinement of $`l,l^{}`$.
Proof of Lemma 3.3 :
Given $`l,l^{}`$ consider the graph $`\stackrel{~}{\gamma }:=\gamma \left(l\right)\gamma \left(l^{}\right)`$. Choose any $`\stackrel{~}{l}`$ such that $`\stackrel{~}{\gamma }=\gamma \left(\stackrel{~}{l}\right)`$. Then $`l,l^{}\stackrel{~}{l}`$.
$`\mathrm{}`$
Next we need the notion of a projection of symplectic manifolds.
###### Definition 3.8
Let $`ll^{}`$, consider any edge $`eE(\gamma )`$ and find the edges $`e_1^{},..,e_n^{}E(\gamma ^{})`$ such that $`e=e_1^{}..e_n^{}`$. We then define the following projection
$$p_{l^{}l}:M_l^{}M_l;h_e:=h_{e_1^{}}..h_{e_n^{}}\text{ and }P^e:=P^{e_1^{}}$$
(3.60)
It is obvious that $`p_{l^{}l}`$ is onto for $`ll^{}`$ except in the presence of boundary conditions in which case the $`P_i^e`$ for sufficiently small $`S^e`$ would be bounded. As a map between the $`\overline{M}_l`$ it would be onto. Define $`m_l=\{h_e,P^e\}_{eE\left(\gamma \right)}`$ and consider an array of non-singular $`dim\left(G\right)\times dim\left(G\right)`$ matrices $`\lambda =\left\{\lambda _e\right\}`$ with an action on the points of $`M_l`$ given by $`\lambda m_l=\{h_e,\lambda _eP^e\}`$.
###### Definition 3.9
i) Consider the uncountable direct product $`:=\times _lM_l`$, then the following subset
$$_{\mathrm{}}:=\left\{\left(m_l\right)_l;\lambda ^{ll^{}}\lambda ^{ll^{}}m_l=p_{l^{}l}\left(m_l^{}\right)ll^{}\right\}$$
(3.61)
is called a generalized projective limit of the $`M_l`$.
ii) A family of symplectic structures $`(\mathrm{\Omega }_l)_l`$ is called a self-consistent or generalized projective family provided that the associated Poisson brackets project in the usual way
$$p_{l^{}l}^{}\{f,g\}_l:=\{p_{l^{}l}^{}f,p_{l^{}l}^{}g\}_l^{}f,gC^{\mathrm{}}\left(M_l\right)$$
(3.62)
that is, the $`p_{l^{}l}`$ are “non-invertible” symplectomorphisms.
It is easy to see that the symplectic structures $`\mathrm{\Omega }_l`$ (or $`\mathrm{\Omega }_\gamma `$ as we called them all the time) that we defined in section 3.2.3 form indeed a self-consistent family of symplectic structures on $`M_l`$ (or $`M_\gamma `$). This follows, as already said, from the astonishing fact that the $`\mathrm{\Omega }_\gamma `$ are completely insensitive to the size and shape of the faces of $`P_\gamma `$ and the choice of the paths of $`\mathrm{\Pi }_\gamma `$ as long as all the requirements of a dual decomposition are met. This is precisely the contents of the identities (3.24) – (3.26). This observation is tied to the fact that the smearing functions, edges and faces, are sufficiently singular and that the smearing process is background metric independent, so that only topological characteristics, such as intersection numbers of edges with faces, are the results of the calculation and are thus completely shape independent, they are locally diffeomorphism invariant (i.e. invariant under locally non-trivial diffeomorphisms). Once more, this observation is also the logic behind definition 3.7 and behind labelling $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ only by elements of $`\mathrm{\Gamma }`$ rather than by elements of $``$.
Let us then summarize :
We have shown in this subsection that the family of differentiable manifolds $`\left(M_l\right)`$ can be given the structure of a generalized projective limit $`_{\mathrm{}}`$ and the family of symplectic structures thereon can be given the structure of a self-consistent family of symplectic structures.
In subsection 3.2.3 on the other hand we showed that there is a sequence $`l_n`$ (there denoted $`\gamma _n`$) such that $`M=lim_n\mathrm{}M_{l_n}`$ and $`\mathrm{\Omega }=lim_n\mathrm{}\mathrm{\Omega }_{l_n}`$ (pointwise limits). Moreover, $`l_ml_n`$ for all $`mn`$, so the sequence is linearly ordered. The points $`m_{l_n}`$ defined in section 3.2.3 belong to $`M_{l_n}`$. By construction, we can extend every such sequence $`m_{l_n}`$ to a sequence $`\left(m_l\right)_l_{\mathrm{}}`$. (Explicitly, the array of matrices is for $`n>m\mathrm{}`$ approximately given by $`\lambda _{eij}^{l_ml_n}=\delta _{ij}\left(ϵ_n/ϵ_m\right)^{D1}`$). It follows that the sequence $`\left(m_{l_n}\right)`$ can be embedded into a generalized projective sequence which in turn defines a point of $`_{\mathrm{}}`$. Likewise, the standard symplectic manifold $`(M,\mathrm{\Omega })`$ can be identified with the sequence $`(M_{l_n},\mathrm{\Omega }_{l_n})`$ of symplectic manifolds which in turn can be extended to a subset of the generalized projective limit and self-consistent symplectic structures thereon respectively (with respect to the generalized projective limit $`_{\mathrm{}}`$).
Remark :
Of course, we have displayed $`(M,\mathrm{\Omega })`$ only as the union of a very special subset of all generalized projective sequences. An arbitrary generalized projective sequence wil not have any obvious interpretation in terms of connections and electric fields on any smooth manifold $`\mathrm{\Sigma }`$. This is possible because the set of graphs in – and dual decompositions of $`\mathrm{\Sigma }`$ have much more structure than the set of points of $`\mathrm{\Sigma }`$. In fact, the picture that emerges is completely combinatorical and only very special configurations of graphs and dual decomposition allow a manifold interpretation. In a sense, without specifying the embedding of abstract graphs and dual decompositions into a concrete $`\mathrm{\Sigma }`$ we are treating all manifolds $`\mathrm{\Sigma }`$ simultaneously. Thus, although in the canonical approach to quantum gravity one starts with a given differential manifold, the emerging classical and quantum theory does not depend any longer on the particular choice of $`\mathrm{\Sigma }`$. Only if one insists on a manifold interpretation there will be restrictions on possible graphs (they have to agree, for instance with the Euler characteristic of $`\mathrm{\Sigma }`$ and the dimension of $`\mathrm{\Sigma }`$) and on the spectra of operators . This opens the possibility to describe topology change within canonical quantum gravity.
## 4 The Gauss Constraint
In this section we implement the Gauss constraint into the theory. On $`(M,\mathrm{\Omega })`$ it is given by the function ($`\mathrm{\Lambda }^i𝒮`$)
$$G\left(\mathrm{\Lambda }\right):=_\mathrm{\Sigma }d^Dx\mathrm{\Lambda }^i\left(x\right)\left[_aE_i^a\left(x\right)+f_{ij}^kA_a^j\left(x\right)E_k^a\left(x\right)\right]$$
(4.1)
which generates infinitesimal gauge transformations
$`F\left(A\right)`$ $``$ $`F\left(A\right)+\{F\left(A\right),G\left(\mathrm{\Lambda }\right)\}=F\left(A+D\mathrm{\Lambda }\right)`$
$`E\left(f\right)`$ $``$ $`E\left(f\right)+\{E\left(f\right),G\left(\mathrm{\Lambda }\right)\}=E\left(f+[\mathrm{\Lambda },f]\right)`$ (4.2)
where $`\mathrm{\Lambda }=\mathrm{\Lambda }^i\tau _i/2,A=A^i\tau _i/2,E=E_i\tau _i/2,f=f^i\tau _i/2,F=F_i\tau _i/2`$. The maps (4) are the infinitesimal versions of the finite gauge transformations
$`F\left(A\right)`$ $``$ $`F\left(\text{Ad}_gAdgg^1\right)`$
$`E\left(f\right)`$ $``$ $`\left[\text{Ad}_gE\right]\left(f\right)`$ (4.3)
where $`\text{Ad}_gv:=gvg^1`$ is the adjoint representation of $`G`$ on $`Lie\left(G\right)`$. Indeed, for infinitesimal $`\mathrm{\Lambda }`$, (4) reproduces (4) to linear order provided we identify $`g\left(x\right)=\mathrm{exp}\left(\mathrm{\Lambda }^i\left(x\right)\tau _i/2\right)`$.
Our task is to write (4.1) in terms of $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$. That is, we must find a function $`G_\gamma \left(\mathrm{\Lambda }\right)`$ on $`M_\gamma `$ such that it converges pointwise on $`M`$ to $`G\left(\mathrm{\Lambda }\right)`$ and such that the limit of its Poisson brackets, that is, its Hamiltonian vector field on $`M_\gamma `$ converges to the Hamiltonian vector field of $`G\left(\mathrm{\Lambda }\right)`$ on $`M`$.
To do this we will proceed as follows :
1) Find the gauge transformations of the coordinates $`h_e,P^e`$ of $`M_\gamma `$ induced by (4).
2) Find a generator $`G_\gamma \left(\mathrm{\Lambda }\right)`$ on $`M_\gamma `$ of these infinitesimal transformations.
3) Study the generator and its Hamiltonian vector field on $`M_\gamma `$ and consider the limit $`\gamma \mathrm{\Sigma }`$.
Notice that this procedure works only because gauge transformations have the special feature to preserve $`M_\gamma `$ as we will see. This is not the case for more general gauge groups such as diffeomorphisms which map between different $`M_\gamma `$’s and which are relevant for quantum general relativity .
It is immediate from the definition of a principal fibre bundle with connection over $`\mathrm{\Sigma }`$ and an associated (under the adjoint representation of $`G`$) vector bundle that under finite gauge transformations $`x\mathrm{\Sigma }g\left(x\right)G`$
$`h_e\left(A\right)`$ $``$ $`g\left(e\left(0\right)\right)h_e\left(A\right)g\left(e\left(1\right)\right)^1`$
$`P^e(A,E)`$ $``$ $`g\left(e\left(0\right)\right)P^e(A,E)g\left(e\left(0\right)\right)^1=:\text{Ad}_{g\left(e\left(0\right)\right)}P^e(A,E)`$ (4.4)
where $`P^e=P_i^e\tau _i/2`$. This follows from the manifestly gauge covariant definition of the basic coordinates of $`M_\gamma `$ given in (3.15). The infinitesimal version of (4) is given by (with $`g\left(x\right)=\mathrm{exp}\left(\mathrm{\Lambda }^i\left(x\right)\tau _i/2\right)=\mathrm{exp}\left(\mathrm{\Lambda }\left(x\right)\right)`$)
$`h_e`$ $``$ $`h_e\mathrm{\Lambda }\left(e\left(0\right)\right)h_eh_e\mathrm{\Lambda }\left(e\left(1\right)\right)`$
$`P^e`$ $``$ $`P^e+[P^e,\mathrm{\Lambda }\left(e\left(0\right)\right)]`$ (4.5)
which should equal $`\{h_e,G_\gamma \left(\mathrm{\Lambda }\right\}_\gamma ,\{P^e,G_\gamma \left(\mathrm{\Lambda }\right)\}_\gamma `$ respectively.
It is immediately clear from (3.26) that the second line of (4) can be obtained by choosing
$$G_\gamma \left(\mathrm{\Lambda }\right)=\underset{vV\left(\gamma \right)}{}\mathrm{\Lambda }^i\left(v\right)\underset{eE\left(\gamma \right),e\left(0\right)=v}{}P_i^e+\text{ more}$$
(4.6)
where “more” should commute with all the $`P_i^e`$. Ansatz (4.6) already correctly reproduces also the $`\mathrm{\Lambda }\left(e\left(0\right)\right)`$ term of the first line of (4), the $`\mathrm{\Lambda }\left(e\left(1\right)\right)`$ term looks similar just that it corresponds to an insertion of $`\tau _i`$ from the right instead of from the left. Since the holonomies Poisson commute among themselves we are led to the following improved ansatz
$$G_\gamma \left(\mathrm{\Lambda }\right)=\underset{vV\left(\gamma \right)}{}\mathrm{\Lambda }^i\left(v\right)\left[\underset{eE\left(\gamma \right),e\left(0\right)=v}{}P_i^e+\underset{eE\left(\gamma \right),e\left(1\right)=v}{}M_{ij}\left(h_e\right)P_j^e\right]$$
(4.7)
where the matrix $`M_{ij}\left(h_e\right)`$ should satisfy $`M_{ij}\left(h_e\right)\tau _jh_e=h_e\tau _i`$ and $`\{P_i^e,M_{jk}\left(h_e\right)P_k^e\}_\gamma =0`$. The first requirement leads to the unique solution
$$M_{ij}\left(h\right)=\frac{1}{N}\text{tr}\left(h\tau _ih^1\tau _j\right)$$
(4.8)
while the second asks us to check the vanishing of (use $`\{.,h^1\}_G=h^1\{.,h\}_Gh^1`$)
$`\{P_i,M_{jk}\left(h\right)P_k\}_G`$ $`=`$ $`M_{jk}\left(h\right)f_{ik}^lP_l+\{P_i,M_{jl}\left(h\right)\}_GP_l`$ (4.9)
$`=`$ $`\left[M_{jk}\left(h\right)f_{ik}^l+{\displaystyle \frac{1}{N}}\text{tr}\left({\displaystyle \frac{\tau _i}{2}}h\tau _jh^1\tau _lh\tau _jh^1{\displaystyle \frac{\tau _i}{2}}\tau _l\right)\right]P_l`$
$`=`$ $`\left[f_{ik}^l+f_{li}^k\right]M_{jk}\left(h\right)P_l`$
which indeed vanishes for $`G`$ semisimple as we assume.
We notice that
$$P_i^e:=M_{ij}\left(h_e\right)P_j^e=P_i^{e^1}$$
(4.10)
which explains intuitively why it is possible that $`\{P_i^e,P_j^e^{}\}_\gamma `$ vanishes for any $`e^{}`$ : while $`P^e`$ depends only on the beginning half segment of the edge $`e`$, $`P^e`$ depends only on the ending half segment of the edge $`e`$ and given the symplectic structure (3.13) a non-vanishing bracket is therefore impossible (modulo the regularization procedure of section 3.2.2). Of course, in retrospect the result should have been guessed on general grounds as what we were trying to construct are the generators $`P,P^{}`$ respectively of left and right translations respectively on $`G`$ which, of course, commute.
In order to check the continuum limit of the function (4.7) we employ the sequence of graphs $`\gamma _n`$ of section 3.2.3. Using the notation of that section, in particular (3.48), we define for $`v\gamma _R,R`$ the quantity
$$E^I\left(v\right)=h_{e_I\left(v\right)}(0,1/2)^1\left[_{S^I\left(v\right)}h_{\rho _{e_I\left(v\right)}\left(x\right)}E\left(x\right)h_{\rho _{e_I\left(v\right)}\left(x\right)}^1\right]h_{e_I\left(v\right)}(0,1/2)$$
which to order $`ϵ_n^{D1}`$ equals $`ϵ_n^{D1}n_a^I\left(v\right)E_i^a\left(v\right)\tau _i`$ as $`n\mathrm{}`$. Then for fixed $`n`$
$`G_{\gamma _n}\left(\mathrm{\Lambda }\right)={\displaystyle \underset{R}{}}{\displaystyle \underset{vV\left(\gamma _R\right)}{}}\mathrm{\Lambda }^i\left(v\right)\left[{\displaystyle \underset{eE\left(\gamma _R\right),v=e\left(0\right)}{}}P_i^e+{\displaystyle \underset{eE\left(\gamma _R\right),v=e\left(1\right)}{}}M_{ij}\left(h_e\right)P_j^e\right]`$ (4.11)
$`=`$ $`{\displaystyle \underset{R}{}}{\displaystyle \underset{vV\left(\gamma _R\right)}{}}{\displaystyle \underset{I=1}{\overset{D}{}}}\left[P_i^I\left(v\right)\mathrm{\Lambda }^i\left(v\right)+M_{ij}\left(h_I\left(v\right)\right)P_j^I\left(v\right)\mathrm{\Lambda }^i\left(X_R\left(X_R^1\left(v\right)+ϵb_I\right)\right)\right]`$
$`=`$ $`{\displaystyle \underset{R}{}}{\displaystyle \underset{vV\left(\gamma _R\right)}{}}\mathrm{\Lambda }^i\left(v\right){\displaystyle \underset{I=1}{\overset{D}{}}}\left[P_i^I\left(v\right)+M_{ij}\left(h_I\left(X_R\left(X_R^1\left(v\right)ϵb_I\right)\right)\right)P_j^I\left(X_R\left(X_R^1\left(v\right)ϵb_I\right)\right)\right]`$
$`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{R}{}}{\displaystyle \underset{vV\left(\gamma _R\right)}{}}\mathrm{\Lambda }^i\left(v\right){\displaystyle \underset{I=1}{\overset{D}{}}}\text{tr}(\tau _i[h_I\left(v\right)E^I\left(v\right)h_I\left(v\right)^1E^I\left(X_R(X_R^1\left(v\right)ϵb_I)\right)]`$
$`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{R}{}}{\displaystyle \underset{vV\left(\gamma _R\right)}{}}\mathrm{\Lambda }^i\left(v\right){\displaystyle \underset{I=1}{\overset{D}{}}}\text{tr}(\tau _i[\{h_I\left(v\right)E^I\left(v\right)h_I\left(v\right)^1E^I\left(v\right)\}`$
$`+\{E^I\left(v\right)E^I\left(X_R(X_R^1\left(v\right)ϵb_I)\right)\}])`$
Consider the two curly brackets in the last line of (4.11). The first one is given to leading order $`ϵ^D`$ by $`Y_I^b\left(v\right)n_a^I\left(v\right)[A_b\left(v\right),E^a\left(v\right)]=det\left(Y_R\right)\left(v\right)[A_a\left(v\right),E^a\left(v\right)]`$. The second one is given to leading order $`ϵ_n^D`$ by
$$\left[\frac{\left(n_a^I\left(X_R\left(t\right)\right)E^a\left(X_R\left(t\right)\right)\right)}{t^I}\right]_{X_R\left(t\right)=v}=n_a^I\left(v\right)Y_I^b\left(v\right)_bE^a\left(v\right)=det\left(Y_R\right)\left(v\right)_aE^a\left(v\right)$$
since $`_I_In_a^I\left(X_R\left(t\right)\right)=0`$. Now the sum of the differences
$`h_I\left(v\right)E^I\left(v\right)h_I\left(v\right)^1E^I\left(v\right)ϵ_n^Ddet\left(Y_R\right)\left(v\right)[A_a\left(v\right),E^a\left(v\right)]`$ and
$`E^I\left(v\right)E^I\left(X_R\left(X_R^1\left(v\right)ϵb_I\right)\right)ϵ_n^Ddet\left(Y_R\right)\left(v\right)_aE^a\left(v\right)`$ can be written as $`ϵ_n^{D+1}det\left(Y_R\right)\left(v\right)K\left(v\right)`$ where $`K`$ is an integrable function. Thus, (4.11) becomes
$`G_{\gamma _n}\left(\mathrm{\Lambda }\right)`$ $`=`$ $`{\displaystyle \underset{R}{}}{\displaystyle \underset{vV\left(\gamma _R\right)}{}}ϵ_n^Ddet\left(Y_R\right)\left(v\right)\mathrm{\Lambda }^i\left(v\right)\left(_aE_i^a+f_{ij}^kA_a^jE_k^a\right)\left(v\right)`$ (4.12)
$`{\displaystyle \frac{ϵ_n}{N}}{\displaystyle \underset{R}{}}{\displaystyle \underset{vV\left(\gamma _R\right)}{}}ϵ_n^Ddet\left(Y_R\right)\left(v\right)\mathrm{\Lambda }^i\left(v\right){\displaystyle \underset{I=1}{\overset{D}{}}}\text{tr}\left(\tau _iK\left(v\right)\right)`$
and both sums are Riemann sum approximations of integrals. Recall that $`det\left(Y_R\right)\left(v\right)>0`$ for $`vR`$ and that $`ϵ_n^Ddet\left(Y\right)\left(v\right)`$ approximates the Lebesgue measure of the image under $`X_R`$ of a cube of volume $`ϵ`$ in $`V_R`$. It follows that
$$\underset{n\mathrm{}}{lim}G_{\gamma _n}\left(\mathrm{\Lambda }\right)=G\left(\mathrm{\Lambda }\right)\underset{n\mathrm{}}{lim}\frac{ϵ_n}{N}_\mathrm{\Sigma }d^Dx\mathrm{\Lambda }^i\left(x\right)\text{tr}\left(\tau _iK\left(x\right)\right)=G\left(\mathrm{\Lambda }\right)$$
(4.13)
as desired.
To check that also the Hamiltonian vector field of $`G_\gamma \left(\mathrm{\Lambda }\right)`$ converges to the one of $`G\left(\mathrm{\Lambda }\right)`$ we consider the brackets defined by
$`\{F\left(A\right),G\left(\mathrm{\Lambda }\right)\}^{}`$ $`:=`$ $`\underset{n\mathrm{}}{lim}\{FA^{\left(n\right)},G_{\gamma _n}\left(\mathrm{\Lambda }\right)\}_{\gamma _n}`$
$`\{E\left(f\right),G\left(\mathrm{\Lambda }\right)\}^{}`$ $`:=`$ $`\underset{n\mathrm{}}{lim}\{E^{\left(n\right)}f,G_{\gamma _n}\left(\mathrm{\Lambda }\right)\}_{\gamma _n}`$ (4.14)
where $`A^{\left(n\right)},E^{\left(n\right)}`$ are defined in (3.2.3), (3.2.3). Using the definitions and reasonings repeatedly outlined already in this paper we see that indeed $`\{F\left(A\right),G\left(\mathrm{\Lambda }\right)\}^{}=\{F\left(A\right),G\left(\mathrm{\Lambda }\right)\}`$ and $`\{E\left(f\right),G\left(\mathrm{\Lambda }\right)\}^{}=\{E\left(f\right),G\left(\mathrm{\Lambda }\right)\}`$. So, the Hamiltonian vector fields also coincide in the limit $`\gamma \mathrm{\Sigma }`$ and display (4.7) as a satisfactory discretization of $`G\left(\mathrm{\Lambda }\right)`$.
## 5 Quantization
In this section we quantize all the phase spaces $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$. Notice that in the literature so far one quantized either $`(M,\mathrm{\Omega })`$ or one quantized only one particular family of $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$’s that were defined through lattices in $`\mathrm{\Sigma }`$’s of the topology of $`\text{ }\mathrm{R}^3`$ . In the former case one took a classical function and tried to turn it into an operator after going through some regularization and renormalization steps. In the latter case one started directly with some operators and required that they have a certain continuum limit behaviour with respect to the lattice spacing, however, one did not establish a precise relation between these discrete operators and certin smeared objects of the continuum theory as we did in section 3. However, without such an analysis it is quite unclear what the operators so obtained do actually measure. In particular, one has to postulate the $`ϵ`$ expansion of the $`P^e,h_E`$ rather than being able to derive it from first principles.
By definition, quantization means to find an irreducible representation of an algebra of operators $`\widehat{h}_e,\widehat{P}_e`$ such that the symplectic and the reality structure of the classical theory is correctly implemented. More concretely, since $`M_\gamma `$ is isomorphic with the direct product of co-tangent bundles $`T^{}G`$, one copy for each edge of $`\gamma `$, it is suggested to choose the natural real polarization of the phase space in which wave functions depend only on holonomies. Thus we choose a Hilbert space $`_\gamma `$ of square integrable functions of the $`h_e,eE\left(\gamma \right)`$ with respect to a measure $`\mu _\gamma `$, that is, $`_\gamma =L_2(𝒞_\gamma ,d\mu _\gamma )`$ where $`𝒞_\gamma =G^{\left|E\left(\gamma \right)\right|}`$ is the complete quantum (and also classical in the absence of boundary conditions) configuration space and must represent the operators $`\widehat{h}_e^{AB},\widehat{P}_i^e`$ on $`_\gamma `$ in such a way that the following commutation relations hold :
$`[\widehat{h}_e^{AB},\widehat{h}_e^{}^{CD}]`$ $`=`$ $`0`$
$`[\widehat{P}_j^e,\widehat{h}_e^{}^{AB}]`$ $`=`$ $`i\mathrm{}\delta _e^{}^e\left({\displaystyle \frac{\tau _j}{2}}\widehat{h}_e\right)^{AB}`$
$`[\widehat{P}_j^e,\widehat{P}_k^e^{}]`$ $`=`$ $`i\mathrm{}\delta ^{ee^{}}\left(f_{jk}^l\right)\widehat{P}_l^e`$ (5.1)
More precisely, we must find a common dense domain $`𝒟_\gamma `$ of all the basic operators which they leave invariant so that it makes sense to compute commutators. Notice again that we allow the graph $`\gamma `$ to be infinite.
Furthermore, the reality structure of the classical theory is given by (let us choose $`G`$ to be a (subgroup of a) unitary group for definiteness)
$$\overline{h_e^{AB}}=\left(h_e^1\right)^{BA}\text{ and }\overline{P_i^e}=P_i^e$$
(5.2)
To see the latter, notice that from (3.15) $`P_i^e`$ is given by an integral of quantities of the form $`\text{tr}\left(g\tau _ig^1\tau _j\right)v^j`$ where $`v^j`$ is real and $`gG`$. From $`\overline{g}^Tg=1`$ it follows with $`g=\mathrm{exp}\left(\theta ^j\tau _j/2\right),\theta ^j`$ real that $`\overline{\tau }_j^T=\tau _j`$. Therefore the orthogonal matrix, using $`\text{tr}\left(M^T\right)=\text{tr}\left(M\right)`$,
$$O_{ij}\left(g\right):=\frac{1}{N}\text{tr}\left(g\tau _ig^1\tau _j\right)$$
(5.3)
is real.
In conclusion we must impose the following adjointness relations on $`\mu _\gamma `$
$$\left(\widehat{h}_e^{AB}\right)^{}=\widehat{\left(h_e^1\right)^{BA}}\text{ and }\left(\widehat{P}_i^e\right)^{}=\widehat{P}_i^e$$
(5.4)
where the first identity has to be understood in the sense that one should write the function $`h_e^1`$ in terms of $`h_e`$ and then replace it by $`\widehat{h}_e`$. No operator ordering problems arise since the $`h_e^{AB}`$ are mutually commuting. As advertized, we will choose the $`\widehat{h}_e^{AB}`$ as multiplication operators with values in $`G`$. As $`G`$ is compact, these operators are bounded and thus they are defined, together with $`\widehat{\left(\widehat{h}_e^1\right)^{AB}}`$, everywhere on $`_\gamma `$ so that there are no domain questions at all in the definition of $`\left(\widehat{h}_e^{AB}\right)^{}`$. The second identity in (5.4) says that $`\widehat{P}_j^e`$ is a self-adjoint operator and in order to settle the domain question we will determinine an explicit core $`𝒟_\gamma `$ of essential self-adjointness for all the $`\widehat{P}_i^e`$.
Let us choose as $`𝒟_\gamma :=C^{\mathrm{}}\left(𝒞_\gamma \right)`$ where we consider $`𝒞_\gamma `$ as a Banach manifold modelled on $`\text{ }\mathrm{R}^{dim\left(G\right)\left|E\left(\gamma \right)\right|}`$ similar as for $`M_\gamma `$. Then we choose the following action of the basic operators on $`f_\gamma 𝒟_\gamma `$
$`\left(\widehat{h}_e^{AB}f_\gamma \right)\left(\left\{h_e^{}\right\}\right)`$ $`:=`$ $`h_e^{AB}f_\gamma \left(\left\{h_e^{}\right\}\right)`$
$`\left(\widehat{P}_j^ef_\gamma \right)\left(\left\{h_e^{}\right\}\right)`$ $`:=`$ $`{\displaystyle \frac{i\mathrm{}}{2}}\left(X_j^ef_\gamma \right)\left(\left\{h_e^{}\right\}\right)`$ (5.5)
where $`X_j^e=X\left(h_e\right)_j,X\left(g\right)_i:=\text{tr}\left(\left(\tau _jg\right)^T/g\right)`$ denotes the right invariant vector field on $`G`$ (the generator of left translations). First of all, the operations (5) leave $`𝒟_\gamma `$ obviously invariant. Next we have the Lie algebra of vector fields on $`G`$ given by $`[X_j,X_k]=2f_{jk}^lX_l`$ and it is easy to see that with this choice the commutation relations (5) are identically satisfied.
The direct product structure of $`𝒞_\gamma `$ shows that we may choose
$$d\mu _\gamma \left(\left\{h_e\right\}_{eE\left(\gamma \right)}\right)=_{eE\left(\gamma \right)}d\mu _e\left(h_e\right)$$
and in order that the adjointness relations (5.4) be satisfied it will be sufficient to choose $`\mu _e=\mu _G`$ for all $`eE\left(\gamma \right)`$. Let $`_G=L_2(G,d\mu _G)`$, then we must establish the symmetry property
$$<f,iX_jf^{}>_G=i_Gd\mu _G\left(h\right)\overline{f\left(h\right)}\left(X_jf^{}\right)\left(h\right)=<iX_jf,f^{}>_G$$
(5.6)
for all $`f,f^{}C^{\mathrm{}}\left(G\right)`$. Notice that $`𝒟\left(X^{}\right)`$, the set of elements $`f_G`$ for which the map $`\psi <f,iX_j\psi >`$ defines a continuous linear functional on $`𝒟\left(X\right)=C^{\mathrm{}}\left(G\right)`$, certainly contains $`𝒟\left(X\right)`$ so that (5.6) implies symmetry.
Now $`X_j`$ generates left translations and $`G=\mathrm{}`$, thus, if we choose the measure $`\mu _G`$ to be left invariant we are done provided we can establish that the expression for $`X_j`$ is real valued. For compact groups the only solution is, up to a normalization, $`\mu _G=\mu _H`$, the Haar measure on $`G`$ which is simultaneously left and right invariant and normalized, $`<1,1>_G=1`$. To see that the expression for $`X_j`$ is real we perform the following calculation :
$`<f,X_jf^{}>_G`$ $`=`$ $`{\displaystyle _G}𝑑\mu _H\left(h\right)\overline{f\left(h\right)}{\displaystyle \frac{d}{dt}}_{t=0}f^{}\left(e^{t\tau _j/2}h\right)`$ (5.7)
$`=`$ $`{\displaystyle \frac{d}{dt}}_{t=0}{\displaystyle _G}𝑑\mu _H\left(h\right)\overline{f\left(h\right)}f^{}\left(e^{t\tau _j/2}h\right)`$
$`=`$ $`{\displaystyle \frac{d}{dt}}_{t=0}{\displaystyle _G}𝑑\mu _H\left(e^{t\tau _j/2}h\right)\overline{f\left(e^{t\tau _j/2}h\right)}f^{}\left(h\right)`$
$`=`$ $`{\displaystyle \frac{d}{dt}}_{t=0}{\displaystyle _G}𝑑\mu _H\left(h\right)\overline{f\left(e^{t\tau _j/2}h\right)}f^{}\left(h\right)`$
$`=`$ $`<X_jf,f^{}>_G`$
as claimed.
Finally, to see that $`iX_j`$ is essentially self-adjoint with core $`𝒟\left(X\right)=C^{\mathrm{}}\left(G\right)`$ we show that $`\left[X_j\pm \text{id}__G\right]𝒟_G`$ is dense in $`_G`$ (basic criterion of essential self-adjointness). The proof is simplified through an appeal to the Peter&Weyl theorem : the Hilbert space $`_G`$ is the completion of a countable orthogonal sum of finite dimensional Hilbert spaces $`_\pi `$ where $`\pi `$ runs through the set of equivalence classes of irreducible representations of $`G`$. A complete orthonormal basis of $`_\pi `$ is given by the functions $`\sqrt{d_\pi }\pi _{mn}\left(h\right)`$ where $`d_\pi `$ is the dimension of the representation and $`\pi _{mm^{}}`$ denotes the matrix elements of an abrbitary but fixed representant of that equivalence class. These functions obviously belong to $`𝒟\left(X\right)`$ and finite linear combinations of such functions are still in $`𝒟\left(X\right)`$. Thus, the finite linear combinations of such functions belong to the domain, $`_\pi _\pi 𝒟\left(X\right)`$.
Next, it is easy to see that $`X_j`$ preserves $`_\pi `$. Denote by $`X_j^\pi `$ the restriction of $`X_j`$ to $`_\pi `$ then $`iX_j^\pi `$ is a symmetric operator on the finite dimensional Hilbert space $`_\pi `$ and therefore self-adjoint on $`_\pi `$ with domain given by all of $`_\pi `$. By the basic criterion for self-adjointness, $`\left[X_j^\pi \pm 1\right]_\pi =_\pi `$. The proof is then complete with the observation that
$`\left[X_j\pm 1\right]_\pi _\pi =_\pi \left(\left[X_j^\pi \pm 1\right]_\pi \right)=_\pi _\pi \left[X_j\pm 1\right]𝒟\left(X\right)`$ (5.8)
$``$ $`_G=\overline{_\pi _\pi }\overline{\left[X_j\pm 1\right]𝒟\left(X\right)}_G`$
Remark :
To see that $`X_j^\pi `$ does not have real eigenvectors in a more elementary way, recall that $`\left(X_j^\pi \right)^2=\lambda _\pi <0`$ is the Laplacian on $`G`$.
In conclusion the (possibly infinite) tensor product of Hilbert spaces
$$_\gamma :=_{eE\left(\gamma \right)}_e=L_2(𝒞_\gamma ,d\mu _{0\gamma }=_{eE\left(\gamma \right)}d\mu _e)$$
(5.9)
where each of the $`_e`$ is isomorphic with $`L_2(G,d\mu _H)`$ is a faithful representation of the canonical commutation relations (5) and of the adjointness relations (5.4). Moreover, given the action (5), it is easy to see that the product Haar measure $`\mu _{0\gamma }`$ is the unique solution, that is, any other measure $`\mu _\gamma `$ which is regular with respect to it must be a constant multiple of it. Notice that infinite products of probability measures are well-defined and $`\sigma `$additive probability measures by the Kolmogorov theorem . Much more will be said about infinite tensor products of Hilbert spaces in the first reference of .
We now must quantize various functions on $`M_\gamma `$. It is at this point where our detailed analysis becomes crucial : while in the continuum theory important functions such as the Gauss constraint (4.1) are written as integrals over polynomials of the field variables $`A\left(x\right),E\left(x\right)`$ at the same point, that is, not as polynomials of the smeared functions $`F\left(A\right),E\left(f\right)`$, the functions on $`M_\gamma `$ are polynomials of the basic observables $`h_e,P^e`$ which are already smeared. Thus, while the quantization of, say, $`G\left(\mathrm{\Lambda }\right)`$ on $`(M,\mathrm{\Omega })`$ can possibly produce UV divergent objects, the quantization of $`G_\gamma \left(\mathrm{\Lambda }\right)`$ on $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ cannot suffer from such problems. (Of course, in both cases factor ordering problems might appear but in the latter case this is only an ambiguity and not the source of a divergence). One might think that problems occur when taking the limit $`\gamma \mathrm{\Sigma }`$, but as we will show, this does not happen.
On the other hand, in both cases we can still have IR divergencies. However, again, from our point of view this is not a problem at all ! Namely, our operator is densely defined, that is, it is an unbounded operator defined on a dense domain. This dense subset of the Hilbert space, however, does not contain states with infinite volume. Nevertheless it is possible to deal with this situation appropriately . In contrast, the perturbative quantization of general relativity is based on a cyclic vector with infinite volume and therefore IR divergencies necessarily occur.
Let us then quantize the Gauss constraint $`G_\gamma \left(\mathrm{\Lambda }\right)`$. We choose to order the momentum variables to the right of the configuration variables and obtain
$$\widehat{G}_\gamma \left(\mathrm{\Lambda }\right)=\underset{eE\left(\gamma \right)}{}\left[\mathrm{\Lambda }^i\left(e\left(0\right)\right)\delta _{ij}\mathrm{\Lambda }^i\left(e\left(1\right)\right)O_{ij}\left(\widehat{h}_e\right)\right]\widehat{P}_j^e$$
(5.10)
Let us check that there are no quantum anomalies. First of all we compute the classical constraint algebra. We have
$`\{O_{ij}\left(h\right)P_j,O_{kl}\left(h\right)P_l\}_G`$ $`=`$ $`\{O_{ij}\left(h\right),O_{kl}\left(h\right)P_l\}_GP_j=O_{kl}\left(h\right)\{O_{ij}\left(h\right),P_l\}_GP_j`$
$`=`$ $`{\displaystyle \frac{1}{N}}O_{kl}\left(h\right)\text{tr}\left({\displaystyle \frac{\tau _l}{2}}h\tau _ih^1\tau _jh\tau _ih^1{\displaystyle \frac{\tau _l}{2}}\tau _j\right)P_j=f_{jl}^mO_{kl}\left(h\right)O_{im}\left(h\right)P_j`$
Now since
$$h\tau _ih^1=\text{Ad}_h\tau _i=O_{ij}\left(h\right)\tau _j$$
(5.12)
we have the identity
$`[\text{Ad}_h\tau _i,\text{Ad}_h\tau _j]=O_{ik}\left(h\right)O_{jl}\left(h\right)[\tau _k,\tau _l]=\text{Ad}_h[\tau _i,\tau _j]`$ (5.13)
$``$ $`O_{ik}\left(h\right)O_{jl}\left(h\right)f_{kl}^m=f_{ij}^kO_{km}\left(h\right)`$
which, when inserted into (5), gives
$$\{O_{ij}\left(h\right)P_j,O_{kl}\left(h\right)P_l\}_G=f_{lm}^jO_{kl}\left(h\right)O_{im}\left(h\right)P_j=f_{ik}^mO_{mj}\left(h\right)P_j$$
(5.14)
or, recalling (4.10),
$$\{P_i^e,P_j^e^{}\}_\gamma =\delta ^{ee^{}}f_{ij}^kP_k^e$$
(5.15)
which is the algebra of left invariant vector fields on $`G`$ (notice the relative minus sign as compared to (3.26).
Thus, since the $`P_i^e,P_j^e^{}`$ Poisson commute by definition of the matrix $`O_{ij}=M_{ij}`$ we immediately get with
$$G_\gamma \left(\mathrm{\Lambda }\right)=\underset{vV\left(\gamma \right)}{}\mathrm{\Lambda }^i\left(v\right)\left[\underset{e\left(0\right)=v}{}P_i^e\underset{e\left(1\right)=v}{}P_i^e\right]$$
(5.16)
that
$$\{G_\gamma \left(\mathrm{\Lambda }\right),G_\gamma \left(\mathrm{\Lambda }^{}\right)\}_\gamma =\underset{vV\left(\gamma \right)}{}\mathrm{\Lambda }^i\left(v\right)\mathrm{\Lambda }^j\left(v\right)f_{ij}^k\left[\underset{e\left(0\right)=v}{}P_k^e+\underset{e\left(1\right)=v}{}P_k^e\right]=G_\gamma \left([\mathrm{\Lambda },\mathrm{\Lambda }^{}]\right)$$
(5.17)
as desired because we infer from (4) that the continuum Poisson algebra is given by (use the Jacobi identity)
$`\{\{E\left(f\right),G\left(\mathrm{\Lambda }\right)\},G\left(\mathrm{\Lambda }^{}\right)\}\{\{E\left(f\right),G\left(\mathrm{\Lambda }^{}\right)\},G\left(\mathrm{\Lambda }\right)\}=\{E\left(f\right),\{G\left(\mathrm{\Lambda }\right),G\left(\mathrm{\Lambda }^{}\right)\}\}`$ (5.18)
$`=`$ $`E\left([\mathrm{\Lambda }^{},[\mathrm{\Lambda },f]][\mathrm{\Lambda },[\mathrm{\Lambda }^{},f]]\right)=E\left([f,[\mathrm{\Lambda }^{},\mathrm{\Lambda }]]\right)=\{E\left(f\right),G\left([\mathrm{\Lambda },\mathrm{\Lambda }^{}]\right)\}`$
Thus, the algebra (5.17) converges to (5.18) by (4.13).
Now, it follows trivially from (5), (5.10) that
$$[\widehat{G}_\gamma \left(\mathrm{\Lambda }\right),\widehat{G}_\gamma \left(\mathrm{\Lambda }^{}\right)]=i\mathrm{}\left(\widehat{G}_\gamma \left([\mathrm{\Lambda },\mathrm{\Lambda }^{}]\right)\right)$$
(5.19)
as required.
Let us summarize :
We started with a continuum phase space $`(M,\mathrm{\Omega })`$ and derived from it a discrete phase space $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ for every graph $`\gamma `$. We also showed that $`(M,\mathrm{\Omega })`$ is the pointwise limit of $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ as $`\gamma \mathrm{\Sigma }`$. Next, we took the Gauss constraint $`G\left(\mathrm{\Lambda }\right)`$ which is a function on $`M`$ and derived from it a function $`G_\gamma \left(\mathrm{\Lambda }\right)`$ on $`M_\gamma `$ which again converges pointwise to $`G\left(\mathrm{\Lambda }\right)`$. Moreover, the Poisson algebra of the $`G_\gamma \left(\mathrm{\Lambda }\right)`$ with respect to $`\mathrm{\Omega }_\gamma `$ closes for every fixed $`\gamma `$ and converges pointwise to the Poisson algebra of the $`G\left(\mathrm{\Lambda }\right)`$. Finally, we quantized the $`G_\gamma \left(\mathrm{\Lambda }\right)`$ and obtained an anomaly free algebra of quantum constraints $`\widehat{G}_\gamma \left(\mathrm{\Lambda }\right)`$. Then two questions remain :
1.) Does this structure provide us with a quantization of $`G\left(\mathrm{\Lambda }\right)`$ as well ? That is, can we find an operator $`\widehat{G}\left(\mathrm{\Lambda }\right)`$ densely defined on all of $``$ and not only on $`_\gamma `$ such that
$$\widehat{G}\left(\mathrm{\Lambda }\right)f_\gamma =\widehat{G}_\gamma \left(\mathrm{\Lambda }\right)f_\gamma $$
(5.20)
for every function $`f_\gamma `$ cylindrical over a graph $`\gamma `$ ?
2.) If $`\widehat{G}\left(\mathrm{\Lambda }\right)`$ exists, does its classical limit coincide with the classical function $`G\left(\mathrm{\Lambda }\right)`$ ?
\[1.\]
It is easy to see that the first question can be answered affirmatively :
Namely, in order that (5.20) holds it is sufficient to show that the family of operators $`\widehat{G}_\gamma \left(\mathrm{\Lambda }\right)`$ is consistently defined. But this is trivially the case because we defined a function to be cylindrical over $`\gamma `$ if and only if it is a finite linear combination of spin-network functions which by definition depend non-trivially on the holnomy along each of its edges (that is, each edge is labelled with a non-trivial irreducible representation of $`G`$). Thus, if we superpose $`f_\gamma ,f_\gamma ^{}^{}`$ with $`\gamma \gamma ^{}`$ then $`G\left(\mathrm{\Lambda }\right)\left[f_\gamma +f_\gamma ^{}^{}\right]:=G_\gamma \left(\mathrm{\Lambda }\right)f_\gamma +G_\gamma ^{}\left(\mathrm{\Lambda }\right)f_\gamma ^{}^{}`$. It is also easy to see that this definition leads to the constraint algebra
$$[\widehat{G}\left(\mathrm{\Lambda }\right),\widehat{G}\left(\mathrm{\Lambda }^{}\right)]=i\mathrm{}(\widehat{G}(\left[\mathrm{\Lambda }\right),\mathrm{\Lambda }^{}])$$
(5.21)
by (5.19) since $`\widehat{G}_\gamma \left(\mathrm{\Lambda }\right)`$ preserves $`_\gamma `$. Thus, $`\widehat{G}\left(\mathrm{\Lambda }\right)`$ exists and defines a consistent quantum constraint algebra.
\[2.\]
To address the second question we first of all notice that we have shown that
$$G\left(\mathrm{\Lambda }\right)=\underset{\gamma \mathrm{\Sigma }}{lim}\left[\underset{\mathrm{}0}{lim}\widehat{G}_\gamma \left(\mathrm{\Lambda }\right)\right]$$
(5.22)
where the inner bracket has been demonstrated actually only by the usual “quantization rule”. A rigorous proof will be given elesewhere , see also below for a sketch. The outer limit is to be understood pointwise on $`M`$.
What we would like to establish now is the existence of the opposite limiting procedure, that is
$$G\left(\mathrm{\Lambda }\right)=\underset{\mathrm{}0}{lim}\left[\underset{\gamma \mathrm{\Sigma }}{lim}\widehat{G}_\gamma \left(\mathrm{\Lambda }\right)\right]$$
(5.23)
We will understand the inner bracket to be the operator $`\widehat{G}\left(\mathrm{\Lambda }\right)`$ defined in (5.19) through the self-consistent family of projections $`\widehat{G}_\gamma \left(\mathrm{\Lambda }\right)`$.
We can then rigorously define the limits (5.22) and (5.23) as follows :
Let $`\psi _{\gamma ,m}^{\mathrm{}}`$ be a coherent state, explicitly dependent on Planck’s constant, peaked at the point $`mM`$ (a smooth field configuration) in the following sense : For each graph $`\gamma `$ and its associated dual decomposition $`P_\gamma `$ consider the values of the holonomies and momenta $`h_e\left(m\right),P^e\left(m\right)`$ respectively. Then the operators $`\widehat{h}_e,\widehat{P}^e`$ have expectation values $`h_e\left(m\right),P^e\left(m\right)`$ in the state $`\psi _\gamma ^m`$ respectively and satisfy a minimal uncertainty condition.
Let now $`\gamma _n`$ be the family of graphs defined in section 3.2.3. We then consider the expectation values
$$G_n^{\mathrm{}}(\mathrm{\Lambda },m):=<\psi _{\gamma _n,m}^{\mathrm{}},\widehat{G}_{\gamma _n}\left(\mathrm{\Lambda }\right)\psi _{\gamma _n,m}^{\mathrm{}}>_{\gamma _n}$$
(5.24)
Notice that by definition of the Hilbert space $``$ and the operator $`\widehat{G}\left(\mathrm{\Lambda }\right)`$ also
$$G_n^{\mathrm{}}(\mathrm{\Lambda },m)=<\psi _{\gamma _n,m}^{\mathrm{}},\widehat{G}_{\gamma _n}\left(\mathrm{\Lambda }\right)\psi _{\gamma _n,m}^{\mathrm{}}>$$
(5.25)
Then the limit (5.22) means that
$$G(\mathrm{\Lambda },m)=\underset{n\mathrm{}}{lim}\left[\underset{\mathrm{}\mathrm{}}{lim}G_n^{\mathrm{}}(\mathrm{\Lambda },m)\right]$$
(5.26)
where now the inner limit is taken at fixed $`m,n`$ and is meant in the sense of complex numbers. The limit (5.23) on the other hand means that
$$G(\mathrm{\Lambda },m)=\underset{\mathrm{}\mathrm{}}{lim}\left[\underset{n\mathrm{}}{lim}G_n^{\mathrm{}}(\mathrm{\Lambda },m)\right]$$
(5.27)
and will be much more difficult to check for a more general operator because the $`\mathrm{}`$ corrections of the inner bracket might not converge. In our case, however, both limits are immediate and in fact reproduce $`G\left(\mathrm{\Lambda }\right)`$ as we will show as an example in the first publication of .
To conclude, we have shown that there is an anomaly-free quantization of the Gauss constraint on the continuum Hilbert space $``$ with the corrrect classical limit. The limit (5.22) says that the regularization procedure is meaningful while the limit (5.23) shows that the regulator can be removed without picking up divergencies and such that we obtain the correct classical limit.
## 6 Non-Commutativity Issues
The authors of considered the following classical functions on $`(M,\mathrm{\Omega })`$
$$E(S,f):=_S(E_if^i)\left(x\right)$$
(6.1)
where $`S`$ is an oriented smooth (D-1)-dimensional submanifold of $`\mathrm{\Sigma }`$ and $`f^i𝒮`$. Notice that $`E(S,f)`$ in contrast to our $`P\left(S\right)`$ of (3.15) is not gauge covariant for any choice of $`f`$ and that $`P\left(S\right)E(S,f)`$ since $`P\left(S\right)`$ depends explicitly on both $`A`$ and $`E`$ while (6.1) depends only on $`E`$.
In order to compute the Poisson brackets among the $`E(S,f)`$ induced by the symplectic structure $`\mathrm{\Omega }`$ one should introduce, as in section 3.2.3, a one parameter family of surfaces $`tS_t,t[1,1],S_0=S`$ and smooth regulator functions $`g_ϵ\left(t\right),lim_{ϵ0}g_ϵ\left(t\right)=\delta \left(t\right)`$. One obtains regulated quantities
$$E_ϵ(S,f):=_1^1𝑑tg_ϵ\left(t\right)_{X^1\left(S_t\right)}d^2x\left(X^{}f_a^iE_i^a\right)\left(x\right)$$
(6.2)
at the aid of which we compute the Poisson brackets
$$\{E(S,f),E(S^{},f^{})\}:=\underset{ϵ0}{lim}\{E_ϵ(S,f),E_ϵ(S^{},f^{})\}_\mathrm{\Omega }=0$$
(6.3)
by (3.13).
The authors of now proceeded as follows :
Since, according to the symplectic structure of $`\mathrm{\Omega }`$, we formally have $`\{E_i^a\left(x\right),A_b^j\left(y\right)\}=\delta _b^a\delta _i^j\delta ^{\left(D\right)}(x,y)`$, they represented the operator $`\widehat{E}_i^a\left(x\right)`$ by the functional derivative $`i\delta /\delta A_a^i\left(x\right)`$ defined on functions of smooth connections, substituted this derivative into (6.1), applied it to functions $`f_\gamma `$ of holonomies of smooth connections over a graph $`\gamma `$ and extended the final operator to distributional connections. The result is the following :
Without loss of generality we can subdivide the graph sufficiently and orient all the edges of $`\gamma `$ in such a way that any edge of $`\gamma `$ belongs to one of the following four categories : i) $`eS=\mathrm{}`$, ii) $`eS=e`$, iii) $`eS=e\left(0\right)`$ and $`e`$ lies on the “up” side of $`S`$ or iv) $`eS=e\left(0\right)`$ and $`e`$ lies on the “down” side of $`S`$. Notice that in case iii),iv) the edge is allowed to be tangent at $`e\left(0\right)`$. Denote the subset of edges belonging to category iii) and iv) respectively by $`E_S^u\left(\gamma \right)`$ and $`E_S^d\left(\gamma \right)`$ respecively and the subset of vertices in $`S\left(E_S^u\left(\gamma \right)E_S^d\left(\gamma \right)\right)`$ by $`V_S\left(\gamma \right)`$. Then
$$\widehat{E}(S,f)f_\gamma =i\mathrm{}\underset{pV_S\left(\gamma \right)}{}\frac{f^i\left(p\right)}{2}\left[\underset{e\left(0\right)=p,eE_S^u\left(\gamma \right)}{}X_i^e\underset{e\left(0\right)=p,eE_S^d\left(\gamma \right)}{}X_i^e\right]f_\gamma $$
(6.4)
where again $`X_i^e`$ denotes the right invariant vector field on the $`e`$’th copy of $`G`$. The expression (6.4) defines a self-consistent family of operators $`\widehat{E}_\gamma (S,f)`$ defined on (a dense subset of) $`_\gamma `$.
The non-commutativity becomes now obvious by choosing for instance $`S=S^{}`$ so that
$$[\widehat{E}(S,f),\widehat{E}(S,f^{})]f_\gamma =\mathrm{}^2\underset{pV_S\left(\gamma \right)}{}\frac{f^i\left(p\right)f^j\left(p\right)f_{ij}^k}{2}\underset{e\left(0\right)=p,eE_S^u\left(\gamma \right)E_S^d\left(\gamma \right)}{}X_k^e]f_\gamma $$
(6.5)
and even worse, one cannot write the right hand side as $`\widehat{E}(S,[f,f^{}])`$ !
These problems can be overcome as follows :
Partition the surface $`S`$ into disjoint open pieces $`S_p`$ carrying the same orientation as $`S`$ such that $`p`$ is the only point of $`V_S\left(\gamma \right)`$ lying in $`S_p`$ and $`_{pV_S\left(\gamma \right)}S_p=S`$ modulo boundary points. For each $`eE_S^u\left(\gamma \right)`$ or $`eE_S^d\left(\gamma \right)`$ respectively, deform $`S_p`$ in an arbitrarily small neighbourhood of $`p`$ into the direction of $`e`$ to a surface $`S_e`$ which intersects $`e`$ transversally in an interior point of $`e`$ but no other edge of $`\gamma `$ and which carries the same or opposite orientation as $`S_p`$. Obviously, these surfaces qualify as part of a dual decomposition of $`\gamma `$. We can now construct from these data the following function on $`M`$ which can also be considered as a function on $`M_\gamma `$
$$E_\gamma (S,f):=\underset{pV_S\left(\gamma \right)}{}f^i\left(p\right)\left[\underset{e\left(0\right)=p,eE_S^u\left(\gamma \right)}{}P_i^e\underset{e\left(0\right)=p,eE_S^d\left(\gamma \right)}{}P_i^e\right]$$
(6.6)
which obviously has classically nothing to do with $`E(S,f)`$. Nevertheless, the results of section 5 tell us that its quantization exactly agrees with (6.4), moreover, the algebra of operators of this kind reflects precisely the symplectic structure $`\mathrm{\Omega }_\gamma `$ which is derived from $`\mathrm{\Omega }`$.
In conclusion, we have demonstrated that the family of operators (6.4) can be considered as bona fide quantizations of a family of classical functions which do not Poisson commute with respect to $`\mathrm{\Omega }`$ and therefore the apparent contradiction between classical Poisson bracket algebra and quantum commutator algebra pointed out in evaporates.
The discussion of this section seems to reveal that not only there is ambiguity in quantizing a given classical functions due to the always existing possibility to add $`\mathrm{}`$ corrections, but also vice versa that there is an ambiguity in taking the classical limit, in the sense that one and the same operator can be considered as a quantization of two different classical functions. However, this is not the case if we insist that we begin with a classical phase space and operators have to have a commutator algebra reflecting the classical Poisson bracket algebra. From this point of view, the functions (6.1) must not be considered as classical limit of the operators (6.4) ! One can still argue that the $`S^e`$ are quite arbitrary and that the classical limit is therefore not really well defined, but as already said before, a well-defined classical limit can only be expected in the limit of $`\gamma \mathrm{\Sigma }`$ in a definite way in which the arbitrariness of the $`S^e`$ is lost.
Besides, the functions (6.4) are unphysical already from the point of view of the Gauss constraint : it is impossible to build from them gauge invariant observables except in the limit of infinitesimal faces where they have been used to build geometrical operators . However, in that limit we get anyway $`E(S,f)P\left(S^e\right)_if^i\left(e\left(0\right)\right)`$ so that one can equally well construct these operators from the $`P\left(S^e\right)_i`$ and so there is finally complete agreement between all the results previously obtained in the literature and our approach.
Acknowledgements
We thank O. Winkler for a careful reading of the manuscript.
## Appendix A The Symplectic Structure for $`G=SU(2),U(1)`$ as a Two-Form
Let us first fix our conventions :
For a $`p`$form $`\omega =\omega _{a_1..a_p}dx^{a_1}..dx^{a_p}=\omega _{a_1..a_p}dx^{a_1}..dx^{a_p}`$ on a finite dimensional manifold $`M`$ with $`\omega _{a_1..a_p}=\omega _{[a_1..a_p]}`$ we define exterior differential, interior products with vector fields $`v`$ and Lie derivatives respectively by
$`d\omega `$ $`=`$ $`_a\omega _{a_1..a_p}dx^adx^{a_1}..dx^{a_p}`$ (A.1)
$`i_v\omega `$ $`=`$ $`pv^a\omega _{aa_1..a_{p1}}dx^{a_1}..dx^{a_{p1}}`$ (A.2)
$`_v\omega `$ $`=`$ $`\left[i_vd+di_v\right]\omega `$ (A.3)
Let now $`(M,\mathrm{\Omega })`$ be a finite dimensional symplectic manifold and $`fC^{\mathrm{}}\left(M\right)`$. We define the Hamiltonian vector field $`\chi _f`$ of $`f`$ by
$$i_{\chi _f}\mathrm{\Omega }+df=0$$
(A.4)
and the Poisson bracket of $`f,gC^{\mathrm{}}\left(M\right)`$ with respect to $`\mathrm{\Omega }`$ by
$$\{f,g\}:=i_{\chi _f}i_{\chi _g}\mathrm{\Omega }=\chi _f\left(g\right)=i_{\chi _f}dg$$
(A.5)
If $`\mathrm{\Omega }=d\mathrm{\Theta }`$ is exact then $`\mathrm{\Theta }`$ is called a symplectic potential for $`\mathrm{\Omega }`$. Here is a quick method of how to go backwards from $`\{.,.\}`$ to $`\mathrm{\Omega }`$ :
Introduce local coordinates $`z^\alpha `$ on $`M`$ and corresponding tensor components $`\mathrm{\Omega }=\frac{1}{2}\mathrm{\Omega }_{\alpha \beta }dz^\alpha dz^\beta `$ and define by $`\mathrm{\Omega }^{\alpha \gamma }\mathrm{\Omega }_{\gamma \beta }=\delta _\beta ^\alpha `$ the inverse tensor. Then the Hamiltonian vecor field of any function $`f`$ is given by $`\mathrm{\Omega }^{\alpha \beta }_\beta f`$. Thus $`\{z^\alpha ,z^\beta \}=\mathrm{\Omega }^{\gamma \delta }\left(_\delta z^\alpha \right)\left(_\gamma z^\beta \right)=\mathrm{\Omega }^{\alpha \beta }`$ and so we just have to invert the matrix of Poisson brackets to obtain $`\mathrm{\Omega }`$.
The reader may verify that with our conventions the symplectic potential $`\mathrm{\Theta }=pdq`$ of $`M=T^{}\text{ }\mathrm{R}`$ leads to $`\{p,q\}=1`$.
### A.1 $`U(1)`$
The Lie algebra of $`U\left(1\right)`$ is spanned by $`i`$ (imaginary unit) and is therefore Abelian. Let $`hU\left(1\right)`$ be a complex number of modulus one. We want to compute the symplectic structure $`\mathrm{\Omega }`$ on $`M=T^{}U\left(1\right)`$ corresponding to the brackets $`\{h,h\}=0=\{p,p\},\{p,h\}=\frac{i}{2}h`$. Let $`z^1=p,z^2=h`$, then $`\mathrm{\Omega }^{\alpha \beta }=ihϵ^{\alpha \beta }/2`$ where $`ϵ^{\alpha \beta }`$ is the completely skew tensor density of weight one. Thus $`\mathrm{\Omega }_{\alpha \beta }=2ϵ_{\alpha \beta }/\left(ih\right)`$ and $`\mathrm{\Omega }=2dpdh/\left(ih\right)=2idpd\mathrm{ln}\left(h\right)`$ which equals $`2dpd\phi `$ locally if we write $`g=\mathrm{exp}\left(i\phi \right)`$. $`\mathrm{\Omega }`$ is real and exact with symplectic potential $`\mathrm{\Theta }=2ipd\mathrm{ln}\left(h\right)`$.
### A.2 $`SU(2)`$
This time there is considerably more work involved and we will only sketch the main steps.
Recall the following normalization conditions for our generators $`\text{tr}\left(\tau _i\tau _j\right)=2\delta _{ij},[\tau _i,\tau _j]=2ϵ_{ijk}\tau _k`$ (for instance $`\tau _j=i\sigma _j`$ the later being the standard Pauli matrices). Let us introduce the following global group coordinates
$$S:=\text{tr}\left(h\right),T^i:=\text{tr}\left(\tau _ih\right)$$
(A.6)
then $`h=\left(ST^i\tau _i\right)/2`$ and we have the following relation $`4S^2=\left(T^i\right)^2`$. Thus, instead of working with $`S,T^i`$ we can work with $`ϵ,T^i`$ where $`ϵ=S/\left|S\right|=0,\pm 1`$ is a discrete parameter. From (3.24) – (3.26) we compute the fundamental Poisson brackets
$`\left(\mathrm{\Omega }^1\right)^{jk}`$ $`:=`$ $`\{T^j,T^k\}=\{ϵ,ϵ\}=\{T^i,ϵ\}=0`$
$`\left(\mathrm{\Omega }^1\right)_j^k`$ $`:=`$ $`\{p_j,T^k\}={\displaystyle \frac{1}{2}}\left[ϵ\delta _{jk}\sqrt{4\left(T^m\right)^2}ϵ_{jkl}T^l\right]`$
$`\{p_j,ϵ\}`$ $`=`$ $`0`$
$`\left(\mathrm{\Omega }^1\right)_{jk}`$ $`:=`$ $`\{p_j,p_k\}=ϵ_{jkl}p_l`$ (A.7)
and certainly $`\left(\mathrm{\Omega }^1\right)_j^k=\left(\mathrm{\Omega }^1\right)_j^k`$. Thus, our task is to invert the $`6`$ x $`6`$ matrix $`\mathrm{\Omega }^1`$ defined in (A.2). We do not worry about the discrete parameter $`ϵ`$ which Poisson commutes with everything in what follows.
Let us introduce the $`3`$ x $`3`$ matrix $`\mathrm{\Lambda }\left(v\right)`$ defined for every vector $`v`$ by $`\mathrm{\Lambda }\left(v\right)_{ij}:=ϵ_{ijk}v^k`$. Let also $`z^\alpha =T^\alpha ,\alpha =1,2,3;z^\alpha =p_{\alpha 3},\alpha =4,5,6`$ and $`\mathrm{\Omega }^{\alpha \beta }:=\left(\mathrm{\Omega }^1\right)^{\alpha \beta }=\{z^\alpha ,z^\beta \}`$. Then $`\mathrm{\Omega }^1`$ is explicitly given by
$$\mathrm{\Omega }^1=\frac{1}{2}\left(\begin{array}{cc}0& S1_3+\mathrm{\Lambda }\left(T\right)\\ S1_3\mathrm{\Lambda }\left(T\right)& 2\mathrm{\Lambda }\left(p\right)\end{array}\right)$$
(A.8)
Thus, the $`6`$ x $`6`$ matrix decomposes into four blocks of $`3`$ x $`3`$ matrices. For the matrix $`\mathrm{\Omega }`$ we now make a similar block matrix ansatz
$$\mathrm{\Omega }=2\left(\begin{array}{cc}\mathrm{\Lambda }\left(a\right)& B\\ B^T& 2\mathrm{\Lambda }\left(c\right)\end{array}\right)$$
(A.9)
and study the relations that we obtain from $`\mathrm{\Omega }\mathrm{\Omega }^1=1_6`$ for the vectors $`a,c`$ and the $`3`$ x $`3`$ matrix $`B`$. Using the relations $`\mathrm{\Lambda }\left(U\right)\mathrm{\Lambda }\left(v\right)=vu(u,v)1_3`$ one finds after very lengthy calculations the result
$$\mathrm{\Omega }=\frac{1}{2}\left(\begin{array}{cc}2\left[\mathrm{\Lambda }\left(p\right)+\frac{TppT}{S}\right]& S1_3\mathrm{\Lambda }\left(T\right)+\frac{TT}{S}\\ S1_3\mathrm{\Lambda }\left(T\right)\frac{TT}{S}& 0\end{array}\right)$$
(A.10)
The matrix (A.10) is singular at $`S=0`$ but we will see that this is merely a coordinate singularity by simply working out $`\mathrm{\Omega }=\frac{1}{2}\mathrm{\Omega }_{\alpha \beta }dz^\alpha dz^\beta `$. The result is
$$\mathrm{\Omega }=\frac{1}{2}[(\mathrm{\Lambda }\left(p\right)+\frac{TppT}{S})_{ij}dT^idT^j+(S1_3\mathrm{\Lambda }\left(T\right)+\frac{TT}{S})_i^jdT^idp_j$$
(A.11)
and we find the following global symplectic potential (so $`\mathrm{\Omega }=d\mathrm{\Theta }`$ is exact)
$$\mathrm{\Theta }=\frac{1}{2}p_j\left(ϵ_{ijk}T_kdT_iSdT^j+T^jdS\right)$$
(A.12)
from which regularity is obvious. We also find the following locally defined momentum conjugate to $`T^i`$
$$\pi _i=\frac{1}{2}\left(ϵ_{ijk}T_kS\delta _{ij}+\frac{T^iT^j}{S}\right)p_j$$
(A.13)
and indeed after lengthy calculations using (A.2) we find that $`\{T^i,T^j\}=\{\pi _i,\pi _j\}=0,\{p_i,T^j\}=\delta _i^j`$. Clearly, $`\pi _i,T^i`$ can only be local Darboux coordinates otherwise we would have displayed $`T^{}SU\left(2\right)T^{}S^3`$ as $`T^{}B^2`$ where $`B^2`$ is a solid ball in $`\text{ }\mathrm{R}^3`$ (we are missing the discrete information coming from $`ϵ`$).
The form (A.12) could be the starting point of symplectic reduction of $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ by the Gauss constraint $`G_\gamma \left(\mathrm{\Lambda }\right)`$ at the classical level already using methods from geometric quantization which has not been done so far in the literature to the best of our knowledge. However, the manifold $`M_\gamma `$ reduced by $`G_\gamma `$ is rather singular except in the case of only one copy of $`G`$ and therefore is unattractive.
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# Relation between Raman spectra and Structure of Amorphous Silicon
## I Introduction
Many structural properties of a-Si, such as defect concentration and variation in mean bond angle, are difficult to determine experimentally. This is because it is impossible to measure directly the coordinates of the atoms in a-Si. However, important information on the structure of a-Si can be obtained indirectly, through a number of experimental techniques. These techniques include neutron, x-ray and Raman scattering, electron-spin resonance, and x-ray photo-absorption. Compared to other methods, Raman scattering is more sensitive to small changes in the short-range order of a-Si. For this reason, Raman measurements on a-Si are frequently used to obtain structural information .
The experimental Raman spectra of a-Si show two distinct peaks, at about 150 cm<sup>-1</sup> and 480 cm<sup>-1</sup>, associated with the transverse acoustic (TA) and the transverse optic (TO) vibrational modes, respectively. Certain features in the Raman spectrum are highly sensitive to the structural properties of the a-Si sample. For example, the width of the TO peak is related to the root-mean-square bond-angle variation $`\mathrm{\Delta }\theta `$ in the amorphous network.
In several computational studies, the relation between $`\mathrm{\Gamma }`$ and $`\mathrm{\Delta }\theta `$ was quantified. All studies indicate a broadening of the TO peak with increasing $`\mathrm{\Delta }\theta `$, but there is no quantitative agreement. Beeman’s linear relation, $`\mathrm{\Gamma }=15+6\mathrm{\Delta }\theta `$, which dates back to 1985, is often used by experimentalists to determine $`\mathrm{\Delta }\theta `$ from Raman measurements. Here, $`\mathrm{\Gamma }`$ is in cm<sup>-1</sup> and $`\mathrm{\Delta }\theta `$ in degrees.
Beeman derived his relation using nine structural models of a-Si. Of these models, five were generated from the same 238-atom, hand-built model by Connell and Temkin, which contains only even-membered rings. In contrast, all simulations on a-Si find an abundance of five- and seven-fold rings. Moreover, these five Connell-Temkin models are statistically dependent and not periodic, consequently containing a large fraction of surface atoms. Experimental values of $`\mathrm{\Delta }\theta `$, based on the radial distribution function of a-Si obtained in neutron-diffraction studies, range from 9.9 to 11.0 degrees. Of the nine structural models used by Beeman, only three exhibit values of $`\mathrm{\Delta }\theta `$ in this range. New techniques to generate a-Si structures, such as ART , as well as more powerful computers, have made it possible to generate larger and more realistic a-Si systems via computer simulation.
Also the description of the Raman scattering process has improved. Beeman used the bond polarizibility model proposed by Alben et al., which dates back to 1975. Characteristic for this model is the inclusion of three weighting parameters, whose values must be set somewhat arbitrarily. Several studies have indicated that the values originally proposed by Alben yield an incorrect value for the depolarization ratio . These studies therefore propose different weights. Since then, other polarizibility models have been proposed, for example by Marinov and Zotov .
In this manuscript, we will re-investigate the relation between $`\mathrm{\Gamma }`$ and $`\mathrm{\Delta }\theta `$ by computer simulation. This simulation is based on a large number of 1000-atom, periodic configurations, with structural properties (radial distribution function, spread in mean bond angle) that are in excellent agreement with experiment. Furthermore, recent advances in neutron scattering techniques have made it possible to directly compare the bond polarizibility models to experiment . We therefore also include a detailed comparison of the model of Alben and the model of Marinov and Zotov to experiment.
Additionally, we present two other methods to obtain structural information from the Raman spectrum. The TA/TO intensity ratio and the location of the TO-peak are believed to be directly related to $`\mathrm{\Delta }\theta `$; these relations will be quantified.
The outline of this paper is as follows. In section II, we explain the generation of the a-Si configurations used in this study. We then discuss how the Raman spectrum is obtained from these configurations. The results and conclusions are presented in sections III and IV, respectively.
## II Method
To calculate Raman spectra, three ingredients are required: (1) a potential describing the atomic interactions in the sample, (2) a continuous random network representing a realistic sample of a-Si, and (3) a model assigning Raman activities to the vibrational eigenmodes of the sample.
In the present work, we use a modified version of the Stillinger-Weber potential for all calculations. This potential has the same functional form as the original SW potential , but with different parameters. The parameters were chosen specifically to describe a-Si, see Ref. .
### A Sample generation
To study the effect of $`\mathrm{\Delta }\theta `$ on the Raman spectrum, we require a number of a-Si configurations with varying values of $`\mathrm{\Delta }\theta `$. These configurations were generated using the activation-relaxation technique (ART) . As was shown in previous studies , ART yields structures in good agreement with experiment. They display a low density of coordination defects, a narrow bond-angle distribution and an excellent overlap with the experimental radial distribution function (RDF). The method is outlined below:
1. Initially, 1000 atoms are placed at random in a periodic cubic cell; the configuration is then relaxed at zero pressure.
2. The configuration is annealed using ART. One ART move consists of two steps: (1) the sample is brought from a local energy minimum to a nearby saddle-point (activation), and (2) then relaxed to a new minimum with a local energy minimization scheme including volume optimization, at zero pressure. The new minimum is accepted with a Metropolis probability at temperature $`T=0.25`$ eV.
3. Every 50 ART moves, up to approximately five ART moves per atom when the energy has reached a plateau, the configuration is stored. For 1000-atom samples, this procedure yields a set containing 100 samples.
This procedure is repeated nine times, generating nine statistically independent sets of 100 correlated configurations each. For each set, we found that $`\mathrm{\Delta }\theta `$ ranges from approximately 10 for the well-annealed configurations, to approximately 14 for the poorly annealed configurations.
### B Calculation of Raman spectra
We focus on the reduced Raman spectrum $`I(\omega )`$, with thermal and harmonic oscillator factors removed. This spectrum is a function of frequency $`\omega `$ and of the form
$$I(\omega )=C(\omega )g(\omega ),$$
(1)
where $`g(\omega )`$ is the vibrational density of states (VDOS) and $`C(\omega )`$ a coupling parameter, which depends on frequency and on the polarization (HH or HV) of the incident light used in the Raman experiment.
To calculate the VDOS, the hessian is calculated. Diagonalization of the hessian gives the frequencies of the vibrational modes, from which the VDOS is obtained.
The function $`C(\omega )`$ is obtained from the polarizibility tensor $`\alpha (\omega )`$. In terms of $`\alpha (\omega )`$, the coupling parameter for HH and HV Raman scattering becomes $`C_{HH}(\omega _p)=7G^2+45A^2`$ and $`C_{HV}(\omega _p)=6G^2`$, respectively, with $`A`$ and $`G^2`$ given by the tensor invariants
$$A=\frac{1}{3}\left[\alpha _{11}+\alpha _{22}+\alpha _{33}\right]$$
(2)
and
$`G^2`$ $`=`$ $`3\left[\alpha _{12}^2+\alpha _{23}^2+\alpha _{31}^2\right]+`$ (3)
$`{\displaystyle \frac{1}{2}}`$ $`\left[(\alpha _{11}\alpha _{22})^2+(\alpha _{22}\alpha _{33})^2+(\alpha _{33}\alpha _{11})^2\right];`$ (4)
see for instance Ref. .
The form of the polarizibility tensor $`\alpha (\omega )`$ still needs to be specified; this is the most uncertain part of the calculation. Several models have been proposed, amongst which the commonly used model of Alben et al. and the more recent model of Marinov and Zotov .
In the model of Alben, a cylindrical symmetry of the individual bonds is assumed and each bond is treated independently as a homopolar, diatomic molecule. Three different forms for the bond polarizibility tensor are introduced:
$`\alpha _1(\omega _p)`$ $`=`$ $`{\displaystyle \underset{l,\mathrm{\Delta }}{}}\stackrel{}{u}_l\stackrel{}{r}_\mathrm{\Delta }\left[\stackrel{}{r}_\mathrm{\Delta }\stackrel{}{r}_\mathrm{\Delta }{\displaystyle \frac{1}{3}}𝐈\right],`$ (5)
$`\alpha _2(\omega _p)`$ $`=`$ $`{\displaystyle \underset{l,\mathrm{\Delta }}{}}\stackrel{}{u}_l\stackrel{}{r}_\mathrm{\Delta }\left[{\displaystyle \frac{1}{2}}\left(\stackrel{}{r}_\mathrm{\Delta }\stackrel{}{u}_l+\stackrel{}{u}_l\stackrel{}{r}_\mathrm{\Delta }\right){\displaystyle \frac{1}{3}}𝐈\right],`$ (6)
$`\alpha _3(\omega _p)`$ $`=`$ $`{\displaystyle \underset{l,\mathrm{\Delta }}{}}(\stackrel{}{u}_l\stackrel{}{r}_\mathrm{\Delta })𝐈.`$ (7)
Here, the summation runs over all atoms $`l`$ in the sample and their nearest neighbors $`\mathrm{\Delta }`$, $`\stackrel{}{r}_\mathrm{\Delta }`$ is the unit vector from the equilibrium position of atom $`l`$ to the nearest neighbor $`\mathrm{\Delta }`$, $`\stackrel{}{u}_l`$ is the displacement vector of atom $`l`$ when it is vibrating in mode $`p`$ and $`𝐈`$ is the unit dyadic. The total polarizibility tensor $`\alpha `$ is a weighted sum of the three terms, i.e. $`\alpha =B_1\alpha _1+B_2\alpha _2+B_3\alpha _3`$. As was stated in the introduction, the precise choice of the weights $`B_i`$ is somewhat arbitrary and this is the major shortcoming of the model. Several studies have indicated that mechanisms 1 and 3 provide the main contribution to the Raman scattering process; these propose to use $`B_1:B_2:B_3`$ proportional to $`2:0:1`$, respectively . In this paper, we will use this set of weights.
The model of Marinov and Zotov (MZ) has no free parameters. In this model, the bond polarizibility is expressed as a sum of three components; two components parallel to the bond arising from bonding and non-bonding electrons and a third component perpendicular to the bond, see Ref . Under these assumptions, the polarizibility tensor takes the form:
$`\alpha (\omega _p)`$ $`=`$ $`{\displaystyle \underset{m}{}}r_m^3[(\stackrel{}{b}_m\stackrel{}{r}_m)\stackrel{}{r}_m\stackrel{}{r}_m`$ (9)
$`+{\displaystyle \frac{1}{2}}(\stackrel{}{b}_m\stackrel{}{r}_m+\stackrel{}{r}_m\stackrel{}{b}_m)].`$
Here, the summation runs over all bonds $`m`$ in the sample, $`\stackrel{}{r}_m`$ is a unit vector parallel to the bond, $`r_m`$ is the bond length and $`\stackrel{}{b}_m`$ is defined as $`\stackrel{}{u}_j\stackrel{}{u}_i`$; where $`\stackrel{}{u}_i`$ and $`\stackrel{}{u}_j`$ are the displacement vectors of atoms $`i`$ and $`j`$ constituting the $`m`$th bond, when vibrating in mode $`p`$.
The coupling parameter for a-Si has also been determined experimentally through neutron scattering methods. According to this experiment, the coupling parameter is a slowly increasing function of frequency. We will use the experimental result of Ref. to test the validity of both the model of Alben and the MZ model.
## III Results
First, in section III A, we compare the polarizibility models of Alben and Marinov and Zotov to experiment. In the subsequent sections, we investigate the relation between the spread in the bond angle $`\mathrm{\Delta }\theta `$ and:
1. Raman TO peak width,
2. Raman TO peak position, and
3. Raman TO/TA intensity ratio.
We show results for HV polarized light only; this is the usual experimental situation. Results for HH polarized light have also been obtained and are available upon request.
### A Raman coupling parameter
The solid curves in Fig. 1 show the HV Raman coupling parameter for a-Si calculated using the model of Alben (top) and the MZ model (bottom) for the bond polarizibility. The experimental result of Ref. is also shown (dashed). For this calculation, we used a well-annealed, 1000-atom configuration with $`\mathrm{\Delta }\theta =10.0^{}`$, since this will most closely resemble the experimental sample. We have checked that the general features of the curves in Fig. 1 do not depend on the details of the configuration used: a number of other, well-annealed, configurations, with $`\mathrm{\Delta }\theta `$ ranging from 10.0 to 11.0, gave similar results.
Fig. 1 shows that both models yield an increasing coupling parameter for frequencies up to around 500 cm<sup>-1</sup>. This is in qualitative agreement with experiment. For higher frequencies, the model calculations predict a sharp decrease in the coupling parameter. This is not confirmed by experiment.
The quantitative agreement with experiment is rather poor, especially in the low-frequency regime; both models provide substantially less activity in this regime than observed in experiment. For this reason, Raman spectra calculated using either of the two models yield TA peak amplitudes far below experimental values. This, in our opinion, is their major shortcoming.
This point is further illustrated in the top graph of Fig. 2, where we show the Raman spectrum calculated using the MZ model (solid) and an experimental spectrum (dashed) taken from Ref. . The experimental spectrum was obtained from ion-implanted a-Si which had been annealed at 500C for two hours. Agreement between model and experiment, particularly in the low-frequency regime, is poor.
Given the overall poor performance of both the model of Alben and the MZ model, it may be feasible to follow a semi-experimental approach in which a computer generated a-Si sample is used to calculate the VDOS and experimental data is used to describe the coupling parameter. This approach is justified because several studies have indicated that changes in the Raman spectrum are due to changes in the VDOS and not to changes in the coupling parameter . The results of this approach are illustrated in the bottom graph of Fig. 2. Here, we show the Raman spectrum obtained by multiplying a computer-generated VDOS with experimental coupling parameter data (solid). To calculate the VDOS, we used a well-annealed a-Si sample with $`\mathrm{\Delta }\theta =10.0^{}`$; the experimental coupling parameter was taken from Ref. . The dashed line shows again the experimental Raman spectrum taken from Ref. . Agreement with experiment has improved substantially.
Of the two models considered here, the MZ model provides slightly more activity in the low-frequency regime than the model of Alben; comparison with experiment would therefore favor the MZ model. However, given the overall poor performance of both models, we will also show in the following sections results obtained using the semi-experimental approach.
### B Raman TO peak width vs. $`\mathrm{\Delta }\theta `$
The intensity of the Raman TO peak decreases abruptly on the high frequency side, but not on the low frequency side. Beeman therefore defines $`\mathrm{\Gamma }`$ as twice the half-width at half the maximum height on the high frequency side of the TO peak, as a meaningful parameter to specify the TO peak width . In this paper, we use the same definition.
The solid lines in Fig. 3 show the relation between $`\mathrm{\Gamma }`$ and $`\mathrm{\Delta }\theta `$ for HV polarized light, derived using the model of Alben (top), the MZ model (middle) and the semi-experimental approach (bottom). Also shown is the result obtained by Beeman (dashed).
Both the model of Alben and the MZ model produce similar results; linear least square fits yield the equations $`\mathrm{\Gamma }/2=3.0\mathrm{\Delta }\theta 6.5`$ and $`\mathrm{\Gamma }/2=3.7\mathrm{\Delta }\theta 13.3`$, respectively. Compared to the result of Beeman, $`\mathrm{\Gamma }/2=3.0\mathrm{\Delta }\theta +7.5`$, we see agreement on the sensitivity (i.e. slope of the lines) of $`\mathrm{\Gamma }`$ to $`\mathrm{\Delta }\theta `$, but not on the overall offset (i.e. intercepts of the lines). The same holds for the result obtained in the semi-experimental approach; least square fitting yields $`\mathrm{\Gamma }/2=3.3\mathrm{\Delta }\theta +9.2`$ in that case.
### C Raman TO peak position vs. $`\mathrm{\Delta }\theta `$
As another way to obtain structural information on a-Si from its Raman spectrum, we investigate the relation between the TO peak frequency ($`\omega _{TO}`$) and $`\mathrm{\Delta }\theta `$. Fig. 4 shows the relation between $`\omega _{TO}`$ and $`\mathrm{\Delta }\theta `$ for HV polarized light, derived using the model of Alben (top), the MZ model (middle) and the semi-experimental approach (bottom).
According to Fig. 4, $`\omega _{TO}`$ shifts to higher frequency as $`\mathrm{\Delta }\theta `$ decreases. However, agreement with experiment for both the model of Alben and the MZ model is poor. Experimental Raman spectra of well-annealed a-Si samples yield $`\omega _{TO}`$ in the order of 480 cm<sup>-1</sup> . For well-annealed configurations, for which $`\mathrm{\Delta }\theta `$ ranges from 9.9 to 11.0 degrees, the models exceed the experimental value by around 20 cm<sup>-1</sup>.
The semi-experimental approach is in much better agreement with experiment; a linear fit yields the equation $`\omega _{TO}=2.5\mathrm{\Delta }\theta +505.5`$ which for $`\mathrm{\Delta }\theta =10.0^{}`$ leads to $`\omega _{TO}=480.5`$ cm<sup>-1</sup>.
The cause of this is that in the model of Alben and the MZ model, $`\omega _{TO}`$ is determined by the peak in the coupling parameter, whereas in the semi-experimental approach, it is determined by the VDOS.
### D Raman TA/TO intensity ratio vs. $`\mathrm{\Delta }\theta `$
Next, we confirm that the TA/TO intensity ratio is directly related to $`\mathrm{\Delta }\theta `$. Fig. 5 shows the relation between reduced Raman TA/TO intensity and $`\mathrm{\Delta }\theta `$ for HV polarized light for the model of Alben (top), the MZ model (middle) and the semi-experimental approach (bottom)
From Fig. 5, we see that the TA/TO intensity ratio increases with increasing $`\mathrm{\Delta }\theta `$. The increase is approximately linear.
HV Raman experiments on well-annealed a-Si samples yield a reduced TA/TO intensity ratio around 0.11 . For both the model of Alben and the MZ model, the HV TA/TO ratio is of order $`10^2`$, i.e. one order of magnitude below the experimental value. This is consistent with the earlier finding that both models underestimate the Raman activity in the low-frequency regime of the spectrum, see section III A.
The semi-experimental approach, on the other hand, yields a TA/TO ratio of 0.14 for $`\mathrm{\Delta }\theta =10.0^{}`$, in better agreement with experiment.
## IV Conclusions
We have generated nine independent sets of 1000-atom samples of a-Si that display a variety of short-range order; the spread in nearest-neighbor bond angles ranges from 10 to 14 degrees. For these samples, the HV Raman spectra are calculated. To describe the Raman scattering process, we have used the earlier bond polarizibility model of Alben et al., the more recent model of Marinov and Zotov as well as experimental data taken from Ref. .
Comparison to experiment shows that both the model of Alben and the MZ model greatly underestimate the Raman activity in the low-frequency regime of the spectrum. This makes these models less suitable to describe low-frequency features of the Raman spectrum, for instance the TA peak. Of the two models considered here, the MZ model is closer to experiment. However, for a more accurate calculation of Raman spectra, we propose a semi-experimental approach. In this approach, the VDOS is obtained in computer simulation and experimental data is used to describe the coupling parameter.
As ways to obtain structural information on a-Si from its Raman spectrum, we have investigated the relation between the TO peak-width $`\mathrm{\Gamma }`$ and $`\mathrm{\Delta }\theta `$, as well as the relations between the TO peak frequency and the TA/TO intensity ratio as functions of $`\mathrm{\Delta }\theta `$.
According to our results, where we used the semi-experimental approach, $`\mathrm{\Gamma }`$ and $`\mathrm{\Delta }\theta `$ are related by $`\mathrm{\Gamma }/2=3.3\mathrm{\Delta }\theta +9.2`$ for HV polarized light. Here, $`\mathrm{\Gamma }`$ is in cm<sup>-1</sup> and $`\mathrm{\Delta }\theta `$ in degrees. Comparing this to the result of Beeman ($`\mathrm{\Gamma }/2=3\mathrm{\Delta }\theta +7.5`$), we find that our result is similar.
Our results also show a shift of the Raman TO peak frequency ($`\omega _{TO}`$) towards higher frequency, as $`\mathrm{\Delta }\theta `$ decreases. In the semi-experimental approach, a linear least-square fit yields $`\omega _{TO}=2.5\mathrm{\Delta }\theta +505.5`$ for HV polarized light. Here, $`\omega _{TO}`$ is in cm<sup>-1</sup> and $`\mathrm{\Delta }\theta `$ in degrees. According to this equation, the shift of $`\omega _{TO}`$ is approximately 7.5 cm<sup>-1</sup>, going from unannealed a-Si $`(\mathrm{\Delta }\theta 13^{})`$ to annealed a-Si $`(\mathrm{\Delta }\theta 10^{})`$. This is in quantitative agreement with experiment .
Finally, we have shown that the reduced Raman TA/TO intensity ratio $`(I)`$ is directly related to $`\mathrm{\Delta }\theta `$; $`I`$ decreases linearly with decreasing $`\mathrm{\Delta }\theta `$. Using the semi-experimental approach, we obtain the relation $`I=0.0078\mathrm{\Delta }\theta +0.0606`$, where $`\mathrm{\Delta }\theta `$ is in degrees.
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# Gauge Field Theory Coherent States (GCS) : III. Ehrenfest Theorems
## 1 Introduction
Quantum General Relativity (QGR) has matured over the past decade to a mathematically well-defined theory of quantum gravity. In contrast to string theory, by definition GQR is a manifestly background independent, diffeomorphism invariant and non-perturbative theory. The obvious advantage is that one will never have to postulate the existence of a non-perturbative extension of the theory, which in string theory has been called the still unknown M(ystery)-Theory.
The disadvantage of a non-perturbative and background independent formulation is, of course, that one is faced with new and interesting mathematical problems so that one cannot just go ahead and “start calculating scattering amplitudes”: As there is no background around which one could perturb, rather the full metric is fluctuating, one is not doing quantum field theory on a spacetime but only on a differential manifold. Once there is no (Minkowski) metric at our disposal, one loses familiar notions such as causality, locality, Poincaré group and so forth, in other words, the theory is not a theory to which the Wightman axioms apply. Therefore, one must build an entirely new mathematical apparatus to treat the resulting quantum field theory which is drastically different from the Fock space picture to which particle physicists are used to.
As a consequence, the mathematical formulation of the theory was the main focus of research in the field over the past decade. The main achievements to date are the following (more or less in chronological order) :
* Kinematical Framework
The starting point was the introduction of new field variables for the gravitational field which are better suited to a background independent formulation of the quantum theory than the ones employed until that time. In its original version these variables were complex valued, however, currently their real valued version, considered first in for classical Euclidean gravity and later in for classical Lorentzian gravity, is preferred because to date it seems that it is only with these variables that one can rigorously define the kinematics and dynamics of Euclidean or Lorentzian quantum gravity .
These variables are coordinates for the infinite dimensional phase space of an $`SU\left(2\right)`$ gauge theory subject to further constraints besides the Gauss law, that is, a connection and a canonically conjugate electric field. As such, it is very natural to introduce smeared functions of these variables, specifically Wilson loop and electric flux functions. (Notice that one does not need a metric to define these functions, that is, they are background independent). This had been done for ordinary gauge fields already before in and was then reconsidered for gravity (see e.g. ).
The next step was the choice of a representation of the canonical commutation relations between the electric and magnetic degrees of freedom. This involves the choice of a suitable space of distributional connections and a faithful measure thereon which, as one can show , is $`\sigma `$-additive. The proof that the resulting Hilbert space indeed solves the adjointness relations induced by the reality structure of the classical theory as well as the canonical commutation relations induced by the symplectic structure of the classical theory can be found in . Independently, a second representation, called the loop representation, of the canonical commutation relations had been advocated (see e.g. and especially and references therein) but both representations were shown to be unitarily equivalent in (see also for a different method of proof).
This is then the first major achievement : The theory is based on a rigorously defined kinematical framework.
* Geometrical Operators
The second major achievement concerns the spectra of positive semi-definite, self-adjoint geometrical operators measuring lengths , areas and volumes of curves, surfaces and regions in spacetime. These spectra are pure point (discete) and imply a discrete Planck scale structure. It should be pointed out that the discreteness is, in contrast to other approaches to quantum gravity, not put in by hand but it is a prediction !
* Regularization- and Renormalization Techniques
The third major achievement is that there is a new regularization and renormalization technique for diffeomorphism covariant, density-one-valued operators at our disposal which was successfully tested in model theories . This technique can be applied, in particular, to the standard model coupled to gravity and to the Poincaré generators at spatial infinity . In particular, it works for Lorentzian gravity while all earlier proposals could at best work in the Euclidean context only (see, e.g. and references therein). The algebra of important operators of the resulting quantum field theories was shown to be consistent . Most surprisingly, these operators are UV and IR finite ! Notice that this result, at least as far as these operators are concerned, is stronger than the believed but unproved finiteness of scattering amplitudes order by order in perturbation theory of the five critical string theories, in a sense we claim that the perturbation series converges. The absence of the divergences that usually plague interacting quantum fields propagating on a Minkowski background can be understood intuitively from the diffeomorphism invariance of the theory : “short and long distances are gauge equivalent”. We will elaborate more on this point in future publications.
* Spin Foam Models
After the construction of the densely defined Hamiltonian constraint operator of , a formal, Euclidean functional integral was constructed in and gave rise to the so-called spin foam models (a spin foam is a history of a graph with faces as the history of edges) . Spin foam models are in close connection with causal spin-network evolutions , state sum models and topological quantum field theory, in particular BF theory . To date most results are at a formal level and for the Euclidean version of the theory only but the programme is exciting since it may restore manifest four-dimensional diffeomorphism invariance which in the Hamiltonian formulation is somewhat hidden.
* Finally, the fifth major achievement is the existence of a rigorous and satisfactory framework for the quantum statistical description of black holes which reproduces the Bekenstein-Hawking Entropy-Area relation and applies, in particular, to physical Schwarzschild black holes while stringy black holes so far are under control only for extremal charged black holes.
Summarizing, the work of the past decade has now culminated in a promising starting point for a quantum theory of the gravitational field plus matter and the stage is set to pose and answer physical questions.
The most basic and most important question that one should ask is : Does the theory have classical general relativity as its classical limit ? Notice that even if the answer is negative, the existence of a consistent, interacting, diffeomorphism invariant quantum field theory in four dimensions is already a quite non-trivial result. However, we can claim to have a satisfactory quantum theory of Einstein’s theory only if the answer is positive.
To settle this issue we have launched an attack based on coherent states which has culminated in a series of papers called “Gauge Field Theory Coherent States” and this paper is the third one of this collection (to be continued). It is closely connected with the companion paper . In we established peakedness properties of the coherent states of the heat kernel family introduced by Hall for arbitrary compact gauge groups which were later applied to gauge field theories in . The results of rest on the explicit determination of the configuration space complexification via the complexifier framework . They reveal that the heat kernel family more or less has all the properties that one would like coherent states to have and that one is used to from the harmonic oscillator coherent states. In particular, these states $`\psi _m^t`$ are labelled by a point $`m=(q,p)`$ in the classical phase space and 1) are eigenstates of certain annihilation operators, 2) are overcomplete, 3) saturate the unquenched Heisenberg uncertainty bound and 4) are peaked in the configuration representation at $`x=q`$, in the momentum representation at $`k=p`$ and in the Bargmann-Segal representation at $`z=qipm`$. Here $`t`$ is a classicality parameter proportional to Planck’s constant $`\mathrm{}`$ and the peak is the sharper (the decay width $`\sqrt{t}`$ the smaller) the smaller $`t`$ and resembles almost a Gaussian.
The properties listed ensure that normal ordered products of creation and annihilation operators have exactly the expectation value, with respect to $`\psi _m^t`$, given by the product of the associated classical functions evaluated at the phase space point $`m`$ without any quantum corrections. However, to establish that this expectation value property also holds with respect to the elementary operators in terms of which important operators of quantum gauge field theory, such as Hamiltonians, are formulated is not granted a priori. The problem arises because the creation and annihilation operators of are not polynomial functions of the elementary operators which in turn is directly related to the kinematically non-linear nature of the theory.. Therefore, the framework of to prove Ehrenfest theorems and the determination of the classical limit by using the harmonic oscillator coherent states does not extend to our case since the methods of crucially rest on the assumption that the basic operators are linear combinations of creation and annihilation operators.
The present paper is devoted to filling this gap. As in all the proofs will be carried out for the case of rank one gauge groups, that is $`G=SU\left(2\right),U\left(1\right)`$, only but by the arguments given in they should readily extend to the case of an arbitrary compact gauge group which we leave for future work . With this restriction, the main result of the present article is that the Ehrenfest property, to zeroth and first order, indeed holds for our coherent states. In other words, the expectation values of polynomials of the elementary operators as well as of an important operator, associated with the volume operator of quantum general relativity mentioned above, which is not a polynomial (not even analytical !) function of the elementary operators, reproduce, to zeroth order in $`t`$, the values of the correponding classical functions at the phase space point given by the coherent state. Moreover, the expectation values of commutators divided by $`it`$ reproduces the corresponding Poisson bracket, to zeroth order in $`t`$, at the given phase space point. These results imply that the quantum dynamics of the operators constructed in , as expected, reproduce the infinitesimal classical dynamics of general relativity , putting the worries raised in ad acta.
The architecture of the present article is as follows :
Section two summarizes the classical and quantum kinematical framework for diffeomorphism invariant quantum gauge field theories.
In section three, after a brief review of the heat kernel coherent states, we prove the above mentioned Ehrenfest theorems for the case of the gauge-variant coherent states for the gauge group $`G=SU\left(2\right)`$. As stated already in , we are mostly interested in gauge-variant coherent states because a) only those manage to verify the satisfaction of a consistent quantum constraint algebra and b) the expectation values of gauge – and diffeomorphism invariant operators are gauge – and diffeomorphism invariant since both gauge groups are represented unitarily on the Hilbert space.
Finally, in appendix A we repeat the analysis of section three for the gauge group $`G=U\left(1\right)`$. As in , the Abelian nature of $`U\left(1\right)`$ shortens all the proofs given for $`SU\left(2\right)`$ by an order of magnitude and the reader is urged to study first the appendix before delving into the technically much harder section three.
## 2 Kinematical Structure of Diffeomorphism Invariant Quantum Gauge Theories
In this section we will recall the main ingredients of the mathematical formulation of (Lorentzian) diffeomorphism invariant classical and quantum field theories of connections with local degrees of freedom in any dimension and for any compact gauge group. See and references therein for more details. In this section we will take all quantities to be dimensionless. More about dimensionful constants will be said in section 3.
### 2.1 Classical Theory
Let $`G`$ be a compact gauge group, $`\mathrm{\Sigma }`$ a $`D`$dimensional manifold admitting a principal $`G`$bundle with connection over $`\mathrm{\Sigma }`$. Let us denote the pull-back to $`\mathrm{\Sigma }`$ of the connection by local sections by $`A_a^i`$ where $`a,b,c,..=1,..,D`$ denote tensorial indices and $`i,j,k,..=1,..,dim\left(G\right)`$ denote indices for the Lie algebra of $`G`$. Likewise, consider a density-one vector bundle of electric fields, whose pull-back to $`\mathrm{\Sigma }`$ by local sections (their Hodge dual is a $`D1`$ form) is a Lie algebra valued vector density of weight one. We will denote the set of generators of the rank $`N1`$ Lie algebra of $`G`$ by $`\tau _i`$ which are normalized according to $`\text{tr}\left(\tau _i\tau _j\right)=N\delta _{ij}`$ and $`[\tau _i,\tau _j]=2f_{ij}^k\tau _k`$ defines the structure constants of $`Lie\left(G\right)`$.
Let $`F_i^a`$ be a Lie algebra valued vector density test field of weight one and let $`f_a^i`$ be a Lie algebra valued covector test field. We consider the smeared quantities
$$F\left(A\right):=_\mathrm{\Sigma }d^DxF_i^aA_a^i\text{ and }E\left(f\right):=_\mathrm{\Sigma }d^DxE_i^af_a^i$$
(2.1)
While both of them are diffeomorphism covariant, only the latter is gauge covariant, one reason to introduce the singular smearing functions discussed below. The choice of the space of pairs of test fields $`(F,f)𝒮`$ depends on the boundary conditions on the space of connections and electric fields which in turn depends on the topology of $`\mathrm{\Sigma }`$ and will not be specified in what follows.
The set of all pairs of smooth functions $`(A,E)`$ on $`\mathrm{\Sigma }`$ such that (2.1) is well defined for any $`(F,f)𝒮`$ defines an infinite dimensional set $`M`$. We define a topology on $`M`$ through the following globally defined metric :
$`d_{\rho ,\sigma }[(A,E),(A^{},E^{})]`$
$`:=`$ $`\sqrt{{\displaystyle \frac{1}{N}}{\displaystyle _\mathrm{\Sigma }}d^Dx\left[\sqrt{det\left(\rho \right)}\rho ^{ab}\text{tr}\left(\left[A_aA_a^{}\right]\left[A_bA_b^{}\right]\right)+{\displaystyle \frac{[\sigma _{ab}\text{tr}\left([E^aE^a][E^bE^b]\right)}{\sqrt{det\left(\sigma \right)}}}\right]}`$
where $`\rho _{ab},\sigma _{ab}`$ are fiducial metrics on $`\mathrm{\Sigma }`$ of everywhere Euclidean signature. Their fall-off behaviour has to be suited to the boundary conditions of the fields $`A,E`$ at spatial infinity. Notice that the metric (2.1) on $`M`$ is gauge invariant. It can be used in the usual way to equip $`M`$ with the structure of a smooth, infinite dimensional differential manifold modelled on a Banach (in fact Hilbert) space $``$ where $`𝒮\times 𝒮`$. (It is the weighted Sobolev space $`H_{0,\rho }^2\times H_{0,\sigma ^1}^2`$ in the notation of ).
Finally, we equip $`M`$ with the structure of an infinite dimensional symplectic manifold through the following strong (in the sense of ) symplectic structure
$$\mathrm{\Omega }((f,F),(f^{},F^{}))_m:=_\mathrm{\Sigma }d^Dx\left[F_i^af_a^iF_i^af_a^i\right]\left(x\right)$$
(2.3)
for any $`(f,F),(f^{},F^{})`$. We have abused the notation by identifying the tangent space to $`M`$ at $`m`$ with $``$. To prove that $`\mathrm{\Omega }`$ is a strong symplectic structure one uses standard Banach space techniques. Computing the Hamiltonian vector fields (with respect to $`\mathrm{\Omega }`$) of the functions $`E\left(f\right),F\left(A\right)`$ we obtain the following elementary Poisson brackets
$$\{E\left(f\right),E\left(f^{}\right)\}=\{F\left(A\right),F^{}\left(A\right)\}=0,\{E\left(f\right),A\left(F\right)\}=F\left(f\right)$$
(2.4)
As a first step towards quantization of the symplectic manifold $`(M,\mathrm{\Omega })`$ one must choose a polariztion. As usual in gauge theories, we will use a particular real polarization, specifically connections as the configuration variables and electric fields as canonically conjugate momenta. As a second step one must decide on a complete set of coordinates of $`M`$ which are to become the elementary quantum operators. The analysis just outlined suggests to use the coordinates $`E\left(f\right),F\left(A\right)`$. However, the well-known immediate problem is that these coordinates are not gauge covariant. Thus, we proceed as follows :
Let $`\mathrm{\Gamma }_0^\omega `$ be the set of all piecewise analytic, finite, oriented graphs $`\gamma `$ embedded into $`\mathrm{\Sigma }`$ and denote by $`E\left(\gamma \right)`$ and $`V\left(\gamma \right)`$ respectively its sets of oriented edges $`e`$ and vertices $`v`$ respectively. Here finite means that $`E\left(\gamma \right)`$ is a finite set. (One can extend the framework to $`\mathrm{\Gamma }_0^\omega `$, the restriction to webs of the set of piecewise smooth graphs but the description becomes more complicated and we refrain from doing this here). It is possible to consider the set $`\mathrm{\Gamma }_\sigma ^\omega `$ of piecewise analytic, infinite graphs with an additional regularity property but for the purpose of this paper it will be sufficient to stick to $`\mathrm{\Gamma }_0^\omega `$. The subscript <sub>0</sub> as usual denotes “of compact support” while <sub>σ</sub> denotes “$`\sigma `$-finite”.
We denote by $`h_e\left(A\right)`$ the holonomy of $`A`$ along $`e`$ and say that a function $`f`$ on $`𝒜`$ is cylindrical with respect to $`\gamma `$ if there exists a function $`f_\gamma `$ on $`G^{\left|E\left(\gamma \right)\right|}`$ such that $`f=p_\gamma ^{}f_\gamma =fp_\gamma `$ where $`p_\gamma \left(A\right)=\left\{h_e\left(A\right)\right\}_{eE\left(\gamma \right)}`$. Holonomies are invariant under reparameterizations of the edge and in this article we assume that the edges are always analyticity preserving diffeomorphic images from $`[0,1]`$ to a one-dimensional submanifold of $`\mathrm{\Sigma }`$. Gauge transformations are functions $`g:\mathrm{\Sigma }G;xg\left(x\right)`$ and they act on holonomies as $`h_eg\left(e\left(0\right)\right)h_eg\left(e\left(1\right)\right)`$.
Next, given a graph $`\gamma `$ we choose a polyhedronal decomposition $`P_\gamma `$ of $`\mathrm{\Sigma }`$ dual to $`\gamma `$. The precise definition of a dual polyhedronal decomposition can be found in but for the purposes of the present paper it is sufficient to know that $`P_\gamma `$ assigns to each edge $`e`$ of $`\gamma `$ an open “face” $`S_e`$ (a polyhedron of codimension one embedded into $`\mathrm{\Sigma }`$) with the following properties :
(1) the surfaces $`S_e`$ are mutually non-intersecting,
(2) only the edge $`e`$ intersects $`S_e`$, the intersection is transversal and consists only of one point which is an interiour point of both $`e`$ and $`S_e`$,
(3) $`S_e`$ carries the orientation which agrees with the orientation of $`e`$.
Furthermore, we choose a system $`\mathrm{\Pi }_\gamma `$ of paths $`\rho _e\left(x\right)S_e,xS_e,eE\left(\gamma \right)`$ connecting the intersection point $`p_e=eS_e`$ with $`x`$. The paths vary smoothly with $`x`$ and the triples $`(\gamma ,P_\gamma ,\mathrm{\Pi }_\gamma )`$ have the property that if $`\gamma ,\gamma ^{}`$ are diffeomorphic, so are $`P_\gamma ,P_\gamma ^{}`$ and $`\mathrm{\Pi }_\gamma ,\mathrm{\Pi }_\gamma ^{}`$, see for details.
With these structures we define the following function on $`(M,\mathrm{\Omega })`$
$$P_i^e(A,E):=\frac{1}{N}\text{tr}\left(\tau _ih_e(0,1/2)\left[_{S_e}h_{\rho _e\left(x\right)}E\left(x\right)h_{\rho _e\left(x\right)}^1\right]h_e(0,1/2)^1\right)$$
(2.5)
where $`h_e(s,t)`$ denotes the holonomy of $`A`$ along $`e`$ between the parameter values $`s<t`$, $``$ denotes the Hodge dual, that is, of $`E`$ is a $`\left(D1\right)`$form on $`\mathrm{\Sigma }`$, $`E^a:=E_i^a\tau _i`$ and we have chosen a parameterization of $`e`$ such that $`p_e=e\left(1/2\right)`$.
Notice that in contrast to similar variables used earlier in the literature the function $`P_i^e`$ is gauge covariant. Namely, under gauge transformations it transforms as $`P^eg\left(e\left(0\right)\right)P^eg\left(e\left(0\right)\right)^1`$, the price to pay being that $`P^e`$ depends on both $`A`$ and $`E`$ and not only on $`E`$. The idea is therefore to use the variables $`h_e,P_i^e`$ for all possible graphs $`\gamma `$ as the coordinates of $`M`$.
The problem with the functions $`h_e\left(A\right)`$ and $`P_i^e(A,E)`$ on $`M`$ is that they are not differentiable on $`M`$, that is, $`Dh_e,DP_i^e`$ are nowhere bounded operators on $``$ as one can easily see. The reason for this is, of course, that these are functions on $`M`$ which are not properly smeared with functions from $`𝒮`$, rather they are smeared with distributional test functions with support on $`e`$ or $`S_e`$ respectively. Nevertheless one would like to base the quantization of the theory on these functions as basic variables because of their gauge and diffeomorphism covariance. Indeed, under diffeomorphisms $`h_eh_{\phi ^1\left(e\right)},P_i^eP_i^{\phi ^1\left(e\right)}`$ where the latter notation means that $`P_e^{\phi ^1\left(e\right)}`$ is labelled by $`\phi ^1\left(S_e\right),\phi ^1\left(\mathrm{\Pi }_\gamma \right)`$. We proceed as follows.
###### Definition 2.1
By $`\overline{M}_\gamma `$ we denote the direct product $`[G\times Lie(G)]^{|E(\gamma )|}`$. The subset of $`\overline{M}_\gamma `$ of pairs $`(h_e(A),P_i^e(A,E))_{eE(\gamma )}`$ as $`(A,E)`$ varies over $`M`$ will be denoted by $`(\overline{M}_\gamma )_{|M}`$. We have a corresponding pull-back map $`p_\gamma :M\overline{M}_\gamma `$ which maps $`M`$ onto $`(\overline{M}_\gamma )_{|M}`$.
Notice that the set $`\left(\overline{M}_\gamma \right)_{|M}`$ is in general a proper subset of $`\overline{M}_\gamma `$, depending on the boundary conditions on $`(A,E)`$, the topology of $`\mathrm{\Sigma }`$ and the “size” of $`e,S_e`$. For instance, in the limit of $`e,S_eeS_e`$ but holding the number of edges fixed, $`\left(\overline{M}_\gamma \right)_{|M}`$ will consist of only one point in $`M_\gamma `$. This follows from the smoothness of the $`(A,E)`$.
We equip a subset $`M_\gamma `$ of $`\overline{M}_\gamma `$ with the structure of a differentiable manifold modelled on the Banach space $`_\gamma =\text{ }\mathrm{R}^{2dim\left(G\right)\left|E\left(\gamma \right)\right|}`$ by using the natural direct product manifold structure of $`\left[G\times Lie\left(G\right)\right]^{\left|E\left(\gamma \right)\right|}`$. While $`\overline{M}_\gamma `$ is a kind of distributional phase space, $`M_\gamma `$ satisfies suitable regularity properties similar to $`M`$.
In order to proceed and to give $`M_\gamma `$ a symplectic structure derived from $`(M,\mathrm{\Omega })`$ one must regularize the elementary functions $`h_e,P_i^e`$ by writing them as limits (in which the regulator vanishes) of functions which can be expressed in terms of the $`F\left(A\right),E\left(f\right)`$. Then one can compute their Poisson brackets with respect to the symplectic structure $`\mathrm{\Omega }`$ at finite regulator and then take the limit pointwise on $`M`$. The result is the following well-defined strong symplectic structure $`\mathrm{\Omega }_\gamma `$ on $`M_\gamma `$.
$`\{h_e,h_e^{}\}_\gamma `$ $`=`$ $`0`$
$`\{P_i^e,h_e^{}\}_\gamma `$ $`=`$ $`\delta _e^{}^e{\displaystyle \frac{\tau _i}{2}}h_e`$
$`\{P_i^e,P_j^e^{}\}_\gamma `$ $`=`$ $`\delta ^{ee^{}}f_{ij}^kP_k^e`$ (2.6)
Since $`\mathrm{\Omega }_\gamma `$ is obviously block diagonal, each block standing for one copy of $`G\times Lie\left(G\right)`$, to check that $`\mathrm{\Omega }_\gamma `$ is non-degenerate and closed reduces to doing it for each factor together with an appeal to well-known Hilbert space techniques to establish that $`\mathrm{\Omega }_\gamma `$ is a surjection of $`_\gamma `$. This is done in where it is shown that each copy is isomorphic with the cotangent bundle $`T^{}G`$ equipped with the symplectic structure (2.1) (choose $`e=e^{}`$ and delete the label $`e`$).
Now that we have managed to assign to each graph $`\gamma `$ a symplectic manifold $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ we can quantize it by using geometric quantization. This can be done in a well-defined way because the relations (2.1) show that the corresponding operators are non-distributional. This is therefore a clean starting point for the regularization of any operator of quantum gauge field theory which can always be written in terms of the $`\widehat{h}_e,\widehat{P}^e,eE\left(\gamma \right)`$ if we apply this operator to a function which depends only on the $`h_e,eE\left(\gamma \right)`$.
The question is what $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ has to do with $`(M,\mathrm{\Omega })`$. In it is shown that there exists a partial order $``$ on the set $``$ of triples $`l=(\gamma ,P_\gamma ,\mathrm{\Pi }_\gamma )`$. In particular, $`\gamma \gamma ^{}`$ means $`\gamma \gamma ^{}`$ and $``$ is a directed set so that one can form a generalized projective limit $`M_{\mathrm{}}`$ of the $`M_\gamma `$ (we abuse notation in displaying the dependence of $`M_\gamma `$ on $`\gamma `$ only rather than on $`l`$). For this one verifies that the family of symplectic structures $`\mathrm{\Omega }_\gamma `$ is self-consistent in the sense that if $`(\gamma ,P_\gamma ,\mathrm{\Pi }_\gamma )(\gamma ^{},P_\gamma ^{},\mathrm{\Pi }_\gamma ^{})`$ then $`p_{\gamma ^{}\gamma }^{}\{f,g\}_\gamma =\{p_{\gamma ^{}\gamma }^{}f,p_{\gamma ^{}\gamma }^{}g\}_\gamma ^{}`$ for any $`f,gC^{\mathrm{}}\left(M_\gamma \right)`$ and $`p_{\gamma ^{}\gamma }:M_\gamma ^{}M_\gamma `$ is a system of natural projections, more precisely, of (non-invertible) symplectomorphisms.
Now, via the maps $`p_\gamma `$ of definition 2.1 we can identify $`M`$ with a subset of $`M_{\mathrm{}}`$. Moreover, in it is shown that there is a generalized projective sequence $`(\gamma _n,P_{\gamma _n},\mathrm{\Pi }_{\gamma _n})`$ such that $`lim_n\mathrm{}p_{\gamma _n}^{}\mathrm{\Omega }_{\gamma _n}=\mathrm{\Omega }`$ pointwise in $`M`$. This displays $`(M,\mathrm{\Omega })`$ as embedded into a projective generalized limit of the $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$, intuitively speaking, as $`\gamma `$ fills all of $`\mathrm{\Sigma }`$, we recover $`(M,\mathrm{\Omega })`$ from the $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$. Of course, this works with $`\mathrm{\Gamma }_0^\omega `$ only if $`\mathrm{\Sigma }`$ is compact, otherwise we need the extension to $`\mathrm{\Gamma }_\sigma ^\omega `$.
It follows that quantization of $`(M,\mathrm{\Omega })`$, and conversely taking the classical limit, can be studied purely in terms of $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ for all $`\gamma `$. The quantum kinematical framework is given in the next subsection.
### 2.2 Quantum Theory
Let us denote the set of all smooth connections by $`𝒜`$. This is our classical configuration space and we will choose for its coordinates the holonomies $`h_e\left(A\right),e\gamma ,\gamma \mathrm{\Gamma }_0^\omega `$. $`𝒜`$ is naturally equipped with a metric topology induced by (2.1).
Recall the notion of a function cylindrical over a graph from the previous subsection. A particularly useful set of cylindrical functions are the so-called spin-netwok functions . A spin-network function is labelled by a graph $`\gamma `$, a set of non-trivial irreducible representations $`\stackrel{}{\pi }=\left\{\pi _e\right\}_{eE\left(\gamma \right)}`$ (choose from each equivalence class of equivalent representations once and for all a fixed representant), one for each edge of $`\gamma `$, and a set $`\stackrel{}{c}=\left\{c_v\right\}_{vV\left(\gamma \right)}`$ of contraction matrices, one for each vertex of $`\gamma `$, which contract the indices of the tensor product $`_{eE\left(\gamma \right)}\pi _e\left(h_e\right)`$ in such a way that the resulting function is gauge invariant. We denote spin-network functions as $`T_I`$ where $`I=\{\gamma ,\stackrel{}{\pi },\stackrel{}{c}\}`$ is a compound label. One can show that these functions are linearly independent. ¿From now on we denote by $`\stackrel{~}{\mathrm{\Phi }}_\gamma `$ finite linear combinations of spin-network functions over $`\gamma `$, by $`\mathrm{\Phi }_\gamma `$ the finite linear combinations of elements from any possible $`\stackrel{~}{\mathrm{\Phi }}_\gamma ^{},\gamma ^{}\gamma `$ a subgraph of $`\gamma `$ and by $`\mathrm{\Phi }`$ the finite linear combinations of spin-network functions over an arbitrary finite collection of graphs. Clearly $`\stackrel{~}{\mathrm{\Phi }}_\gamma `$ is a subspace of $`\mathrm{\Phi }_\gamma `$. To express this distinction we will say that functions in $`\stackrel{~}{\mathrm{\Phi }}_\gamma `$ are labelled by the “coloured graphs” $`\gamma `$ while functions in $`\mathrm{\Phi }_\gamma `$ are labelled simply by graphs $`\gamma `$ where we abuse notation by using the same symbol $`\gamma `$.
The set $`\mathrm{\Phi }`$ of finite linear combinations of spin-network functions forms an Abelian algebra of functions on $`𝒜`$. By completing it with respect to the sup-norm topology it becomes an Abelian C algebra (here the compactness of $`G`$ is crucial). The spectrum $`\overline{𝒜}`$ of this algebra, that is, the set of all algebraic homomorphisms $`\text{ }\mathrm{C}`$ is called the quantum configuration space. This space is equipped with the Gel’fand topology, that is, the space of continuous functions $`C^0\left(\overline{𝒜}\right)`$ on $`\overline{𝒜}`$ is given by the Gel’fand transforms of elements of $``$. Recall that the Gel’fand transform is given by $`\stackrel{~}{f}\left(\overline{A}\right):=\overline{A}\left(f\right)\overline{A}\overline{𝒜}`$. It is a general result that $`\overline{𝒜}`$ with this topology is a compact Hausdorff space. Obviously, the elements of $`𝒜`$ are contained in $`\overline{𝒜}`$ and one can show that $`𝒜`$ is even dense . Generic elements of $`\overline{𝒜}`$ are, however, distributional.
The idea is now to construct a Hilbert space consisting of square integrable functions on $`\overline{𝒜}`$ with respect to some measure $`\mu `$. Recall that one can define a measure on a locally compact Hausdorff space by prescribing a positive linear functional $`\chi _\mu `$ on the space of continuous functions thereon. The particular measure we choose is given by $`\chi _{\mu _0}\left(\stackrel{~}{T}_I\right)=1`$ if $`I=\{\left\{p\right\},\stackrel{}{0},\stackrel{}{1}\}`$ and $`\chi _{\mu _0}\left(\stackrel{~}{T}_I\right)=0`$ otherwise. Here $`p`$ is any point in $`\mathrm{\Sigma }`$, $`0`$ denotes the trivial representation and $`1`$ the trivial contraction matrix. In other words, (Gel’fand transforms of) spin-network functions play the same role for $`\mu _0`$ as Wick-polynomials do for Gaussian measures and like those they form an orthonormal basis in the Hilbert space $`:=L_2(\overline{𝒜},d\mu _0)`$ obtained by completing their finite linear span $`\mathrm{\Phi }`$.
An equivalent definition of $`\overline{𝒜},\mu _0`$ is as follows :
$`\overline{𝒜}`$ is in one to one correspondence, via the surjective map $`H`$ defined below, with the set $`\overline{𝒜}^{}:=\text{Hom}(𝒳,G)`$ of homomorphisms from the groupoid $`𝒳`$ of composable, holonomically independent, analytical paths into the gauge group. The correspondence is explicitly given by $`\overline{𝒜}\overline{A}H_{\overline{A}}\text{Hom}(𝒳,G)`$ where $`𝒳eH_{\overline{A}}\left(e\right):=\overline{A}\left(h_e\right)=\stackrel{~}{h}_e\left(\overline{A}\right)G`$ and $`\stackrel{~}{h}_e`$ is the Gel’fand transform of the function $`𝒜Ah_e\left(A\right)G`$. Consider now the restriction of $`𝒳`$ to $`𝒳_\gamma `$, the groupoid of composable edges of the graph $`\gamma `$. One can then show that the projective limit of the corresponding cylindrical sets $`\overline{𝒜}_\gamma ^{}:=\text{Hom}(𝒳_\gamma ,G)`$ coincides with $`\overline{𝒜}^{}`$. Moreover, we have $`\left\{\left\{H\left(e\right)\right\}_{eE\left(\gamma \right)};H\overline{𝒜}_\gamma ^{}\right\}=\left\{\left\{H_{\overline{A}}\left(e\right)\right\}_{eE\left(\gamma \right)};\overline{A}\overline{𝒜}\right\}=G^{\left|E\left(\gamma \right)\right|}`$. Let now $`f`$ be a function cylindrical over $`\gamma `$ then
$$\chi _{\mu _0}\left(\stackrel{~}{f}\right)=_{\overline{𝒜}}𝑑\mu _0\left(\overline{A}\right)\stackrel{~}{f}\left(\overline{A}\right)=_{G^{\left|E\left(\gamma \right)\right|}}_{eE\left(\gamma \right)}d\mu _H\left(h_e\right)f_\gamma \left(\left\{h_e\right\}_{eE\left(\gamma \right)}\right)$$
where $`\mu _H`$ is the Haar measure on $`G`$. As usual, $`𝒜`$ turns out to be contained in a measurable subset of $`\overline{𝒜}`$ which has measure zero with respect to $`\mu _0`$.
Let $`\mathrm{\Phi }_\gamma `$, as before, be the finite linear span of spin-network functions over $`\gamma `$ and $`_\gamma `$ its completion with respect to $`\mu _0`$. Clearly, $``$ itself is the completion of the finite linear span $`\mathrm{\Phi }`$ of vectors from the mutually orthogonal $`\stackrel{~}{\mathrm{\Phi }}_\gamma `$. Our basic coordinates of $`M_\gamma `$ are promoted to operators on $``$ with dense domain $`\mathrm{\Phi }`$. As $`h_e`$ is group-valued and $`P^e`$ is real-valued we must check that the adjointness relations coming from these reality conditions as well as the Poisson brackets (2.1) are implememted on our $``$. This turns out to be precisely the case if we choose $`\widehat{h}_e`$ to be a multiplication operator and $`\widehat{P}_j^e=i\mathrm{}\kappa X_j^e/2`$ where $`X_j^e=X_j\left(h_e\right)`$ and $`X^j\left(h\right),hG`$ is the vector field on $`G`$ generating left translations into the $`jth`$ coordinate direction of $`Lie\left(G\right)T_h\left(G\right)`$ (the tangent space of $`G`$ at $`h`$ can be identified with the Lie algebra of $`G`$) and $`\kappa `$ is the coupling constant of the theory. For details see .
## 3 Ehrenfest Theorems
Let us recall the most important facts from .
Instead of working with the quantities $`P_i^e`$ of section 2 we use the dimensionless objects $`p_i^e=P_i^e/a^{n_D}`$. If $`P_i^e`$ is already dimensionless then so is $`a`$ and we choose $`n_D=1`$. Otherwise $`a`$ is an arbitrary but fixed constant of the dimension $`\left[a\right]=\text{cm}^1`$ and the power $`n_D`$ is so chosen that $`p_i^e`$ is dimensionless. In both cases, the numerical value of $`a`$ is macroscopic, say $`a=1`$ or $`a=`$1cm respectively. The power $`n_D`$ depends on the dimensionality of $`\mathrm{\Sigma }`$ and the theory, e.g. $`n_D=2`$ for general relativity in $`D=3`$ spatial dimensions. Also, if $`\kappa `$ is the coupling constant of the theory (the coefficient $`1/\kappa `$ in front of the classical action) then $`P_i^e`$ in (2.1) has to be replaced by $`P_i^e/\kappa `$. It follows from the canonical commutation relations that if $`\widehat{h}_e`$ as before is a multiplication operator in the connection representation then $`\widehat{p}_j^e=itX_j^e/2`$ where
$$t:=\frac{\alpha }{a^{n_D}}\text{ and }\alpha =\kappa \mathrm{}$$
(3.1)
define the classicality parameter and the Feinstruktur constant respectively. For instance, in four-dimensional general relativity $`\alpha =\mathrm{}_p^2`$ is the Planck area and for $`a=1`$cm we have $`\sqrt{t}/a=\mathrm{}_p/cm10^{32}`$. All our estimates are based on the fact that $`t`$ is a tiny positive number.
Consider first only one edge $`e`$ of a graph $`\gamma `$, then we define the complexifier
$$\widehat{C}_e:=\frac{a^{n_D}}{2\kappa }\delta ^{ij}\widehat{p}_i^e\widehat{p}_j^e$$
(3.2)
and the coherent state, labelled by the phase space point $`g_e=e^{ip_j^e\tau _j/2}h_eG^{\text{ }\mathrm{C}}`$ and the classicality parameter $`t`$, by
$$\psi _{e,g_e}^t\left(h_e\right):=\left[e^{\frac{\widehat{C}_e}{\mathrm{}}}\delta _h^{}\left(h_e\right)\right]_{h^{}g_e}=\left[e^{t\mathrm{\Delta }_e/2}\delta _h^{}\left(h_e\right)\right]_{h^{}g_e}$$
(3.3)
where $`\delta _h^{}`$ denotes the $`\delta `$ distribution on $`G`$ with support at $`h^{}`$ and in (3.3) one is supposed to analytically continue $`h^{}`$. One can give the following explicit sum formula for $`\psi _{e,g_e}^t\left(h_e\right)`$,
$$\psi _{e,g_e}^t\left(h_e\right)=\underset{\pi }{}d_\pi e^{t\lambda _\pi /2}\chi _\pi \left(g_eh_e^1\right)$$
(3.4)
where the sum is over the equivalence classes of irreducible representations $`\pi `$ of $`G`$, $`\chi _\pi `$ is the character of $`\pi `$ and $`\lambda _\pi `$ is the eigenvalue of $`\mathrm{\Delta }_e`$ with eigenfunctions $`\chi _\pi \left(g_eh_e^1\right)`$. The operator $`e^{t\mathrm{\Delta }_e/2}`$ is sometimes called the heat kernel operator. Notice that the states $`\psi _{e,g_e}^t`$ are not normalized.
The generalization to the whole graph $`\gamma `$ is straightforward, it is simply given by the product over edges
$$\psi _{\gamma ,g_\gamma }^t\left(h_\gamma \right)=\underset{eE\left(\gamma \right)}{}\psi _{e,g_e}^t\left(h_e\right)$$
(3.5)
where $`g_\gamma =\left\{g_e\right\}_{eE\left(\gamma \right)},h_\gamma =\left\{h_e\right\}_{eE\left(\gamma \right)}`$. The product states (3.5) are obtained by applying the operator $`\mathrm{exp}\left(t\mathrm{\Delta }_\gamma /2\right),\mathrm{\Delta }_\gamma =_{eE\left(\gamma \right)}\mathrm{\Delta }_e`$ to the product delta – distribution
$$\delta _{\gamma ,h_\gamma ^{}}\left(h_\gamma \right)=\underset{eE\left(\gamma \right)}{}\delta _{h_e^{}}\left(h_e\right)$$
(3.6)
followed by analytical continuation. This formula is meaningful only if $`\gamma \mathrm{\Gamma }_0^\omega `$ is a finite graph, for truly infinite graphs in $`\mathrm{\Gamma }_\sigma ^\omega `$ we must work immediately with products of normalized coherent states as otherwise such states would not be normalizable. For the purpose of this paper it will be sufficient to stick with $`\mathrm{\Gamma }_0^\omega `$ for the following reason. Since the Poisson brackets (2.1) are promoted to the following commutation relations
$`[\widehat{h}_e,\widehat{h}_e^{}]`$ $`=`$ $`0`$
$`[\widehat{p}_j^e,\widehat{h}_e^{}]`$ $`=`$ $`it\delta _e^{}^e{\displaystyle \frac{\tau _j}{2}}\widehat{h}_e`$
$`[\widehat{p}_i^e,\widehat{p}_j^e^{}]`$ $`=`$ $`it\delta ^{ee^{}}f_{ij}^k\widehat{p}_k^e`$ (3.7)
on the Hilbert space $`_\gamma `$, the completion of $`\mathrm{\Phi }_\gamma `$, it follows that due to the commutativity of operators labelled by different edges every polynomial of the elementary operators $`\{\widehat{h}_e,\widehat{p}_j^e\}_{eE\left(\gamma \right)}`$ can in fact be reduced to sums of monomials of the form
$$\widehat{O}_\gamma =\underset{eE\left(\gamma \right)}{}\widehat{O}_e$$
(3.8)
where for each $`e`$ the operator $`\widehat{O}_e=\widehat{O}_e(\widehat{h}_e,\widehat{p}_j^e)`$ is a certain polynomial of the $`2dim\left(G\right)`$ independent operators $`\left(\widehat{h}_e\right)_{AB},\widehat{p}_j^e`$ for the same $`e`$ ($`A,B,C,\mathrm{}`$ are group indices). Obviously the order of the operators $`\widehat{O}_e`$ is irrelevant but not the order of the elementary operators appearing in $`\widehat{O}_e`$. It follows that the expectation value of (infinite) sums of monomials of the type (3.8) is given by the same sum over expectation values of monomials and the latter have the following simple product structure with respect to the state (3.5)
$$\frac{<\psi _{\gamma ,g_\gamma }^t,\widehat{O}_\gamma \psi _{\gamma ,g_\gamma }^t>}{\psi _{\gamma ,g_\gamma }^t^2}=\underset{eE\left(\gamma \right)}{}\frac{<\psi _{e,g_e}^t,\widehat{O}_e\psi _{e,g_e}^t>}{\psi _{e,g_e}^t^2}$$
(3.9)
Also, as far as commutators are concerned, notice the commutator formula for monomial operators $`\widehat{O}_\gamma ,\widehat{O}_\gamma ^{}`$ of the type (3.8) given by
$$[\widehat{O}_\gamma ,\widehat{O}_\gamma ^{}]=\underset{eE\left(\gamma \right)}{}[\widehat{O}_e,\widehat{O}_e^{}]\underset{e^{}E\left(\gamma \right)\left\{e\right\}}{}\left(\widehat{O}_e^{}\widehat{O}_e^{}^{}\right)$$
(3.10)
which can be proved by complete induction over $`\left|E\left(\gamma \right)\right|`$. Thus, commutators of monomials again reduce to sums over monomials. We summarize these simple observations by the following theorem.
###### Theorem 3.1
Let $`\gamma \mathrm{\Gamma }_0^\omega `$ be a graph, $`g_\gamma M_\gamma `$ a point in the phase space and $`\widehat{O}_\gamma ,\widehat{O}_\gamma ^{}`$ monomial operators. Suppose that for each $`eE(\gamma )`$ we have
$`\underset{t0}{lim}{\displaystyle \frac{<\psi _{e,g_e}^t,\widehat{O}_e\psi _{e,g_e}^t>}{\psi _{e,g_e}^t^2}}=O_e(h_e\left(g_e\right),p_j^e\left(g_e\right))`$
$`\underset{t0}{lim}{\displaystyle \frac{<\psi _{e,g_e}^t,\frac{[\widehat{O}_e,\widehat{O}_e^{}]}{it}\psi _{e,g_e}^t>}{\psi _{e,g_e}^t^2}}=\{O_e,O_e^{}\}_e((h_e\left(g_e\right),p_j^e\left(g_e\right))`$ (3.11)
where the polar decomposition $`g_e=H_e(g_e)h_e(g_e),H_e(g_e)=\mathrm{exp}(i\tau _jp_j^e(g_e)/2)`$ specifies $`h_e(g_e),p_j^e(g_e)`$ uniquely. Then
$`\underset{t0}{lim}{\displaystyle \frac{<\psi _{\gamma ,g_\gamma }^t,\widehat{O}_\gamma \psi _{\gamma ,g_\gamma }^t>}{\psi _{\gamma ,g_\gamma }^t^2}}=O_\gamma (h_\gamma \left(g_\gamma \right),p_j^\gamma \left(g_\gamma \right))`$
$`\underset{t0}{lim}{\displaystyle \frac{<\psi _{\gamma ,g_\gamma }^t,\frac{[\widehat{O}_\gamma ,\widehat{O}_\gamma ^{}]}{it}\psi _{\gamma ,g_\gamma }^t>}{\psi _{\gamma ,g_\gamma }^t^2}}=\{O_\gamma ,O_\gamma ^{}\}_\gamma (h_\gamma \left(g_\gamma \right),p_j^\gamma \left(g_\gamma \right))`$ (3.12)
where $`p_j^\gamma =\{p_j^e\}_{eE(\gamma )}`$.
This theorem shows that in order to establish Ehrenfest theorems it will be completely sufficient to consider the problem for one copy of the group only. This is even true in the extension from $`\mathrm{\Gamma }_0^\omega `$ to $`\mathrm{\Gamma }_\sigma ^\omega `$ because the operators that appear in applications can be written as infinite sums of monomials each of which depends on a finite subgraph of an infinite graph only. However, if $`\mathrm{\Gamma }_0^\omega \gamma \gamma ^{}\mathrm{\Gamma }_\sigma ^\omega `$ then we can write a given monomial operator $`\widehat{O}_\gamma `$ also as $`\widehat{O}_\gamma ^{}=\widehat{O}_\gamma _{e^{}E\left(\gamma ^{}\right)\gamma }1_e^{}`$ where $`1_e^{}`$ denotes the unit operator on $`_e^{}`$. Thus, we get for the expectation values
$$\frac{<\psi _{\gamma ^{},g_\gamma ^{}}^t,\widehat{O}_\gamma ^{}\psi _{\gamma ^{},g_\gamma ^{}}^t>}{\psi _{\gamma ^{},g_\gamma ^{}}^t^2}=\frac{<\psi _{\gamma ,g_\gamma }^t,\widehat{O}_\gamma \psi _{\gamma ,g_\gamma }^t>}{\psi _{\gamma ,g_\gamma }^t^2}$$
(3.13)
so the problem reduces again to one for $`\gamma \mathrm{\Gamma }_0^\omega `$.
Equation (3.13) seems to indicate that an extension of coherent states to infinite graphs is not really necessary. However, if $`\mathrm{\Sigma }`$ is non-compact then the only way to approximate, say a classical metric in general relativity which is everywhere non-degenerate, by coherent states over graphs $`\gamma \mathrm{\Gamma }_0^\omega `$ is by using a countably infinite superposition of them, say $`\stackrel{~}{\psi }_{\gamma ,g_\gamma }^t=_nz_n\psi _{\gamma _n,g_{\gamma _n}}^t,z_n\text{ }\mathrm{C},n\text{ }\mathrm{N}`$ where $`\gamma =_n\gamma _n`$ and the $`\psi _{\gamma _n,g_{\gamma _n}}^t`$ are the coherent states (3.5). Suppose now that we consider the following operator $`\widehat{O}_\gamma =_n\widehat{O}_{\gamma _n}`$ over $`\gamma `$. In applications it is usually true that the $`\gamma _n`$ are mutually disjoint but they fill $`\mathrm{\Sigma }`$ everywhere with respect to the metric to be approximated, that is, $`\gamma `$ fills $`\mathrm{\Sigma }`$. Now the $`\psi _{\gamma _n,g_{\gamma _n}}^t`$ are not mutually orthogonal, rather for $`nm`$ we have $`<\psi _{\gamma _m,g_{\gamma _m}}^t,\psi _{\gamma _n,g_{\gamma _n}}^t>=1`$ while $`\psi _{\gamma _m,g_{\gamma _m}}^t>1`$ for all $`g_\gamma ,t`$. Now in applications it turns out that
$$<\psi _{\gamma _m,g_{\gamma _m}}^t,\widehat{O}_{\gamma _n}\psi _{\gamma _p,g_{\gamma _p}}^t>=\delta _{m,n}\delta _{n,p}<\psi _{\gamma _n,g_{\gamma _n}}^t,\widehat{O}_{\gamma _n}\psi _{\gamma _n,g_{\gamma _n}}^t>$$
and thus indeed
$$<\stackrel{~}{\psi }_{\gamma ,g_\gamma }^t,\widehat{O}_\gamma \stackrel{~}{\psi }_{\gamma ,g_\gamma }^t>=\underset{n}{}\left|z_n\right|^2<\psi _{\gamma _n,g_{\gamma _n}}^t,\widehat{O}_{\gamma _n}\psi _{\gamma _n,g_{\gamma _n}}^t>$$
yields the correct expectation value provided that $`\left|z_n\right|=1/\psi _{\gamma _n,g_{\gamma _n}}^t`$. However, then the norm squared of $`\stackrel{~}{\psi }_{\gamma ,g_\gamma }^t`$ is given by
$$\stackrel{~}{\psi }_{\gamma ,g_\gamma }^t^2=\left[\underset{n}{}1\right]+2\underset{m<n}{}\mathrm{}\left(\overline{z}_mz_n\right)$$
which is divergent. Thus, the only way to deal with semiclassical physics in the case that $`\mathrm{\Sigma }`$ is non-compact is to use the extension to infinite graphs $`\mathrm{\Gamma }_\sigma ^\omega `$.
The remainder of this section then is subdivided into two parts. In the first one we prove the Ehrenfest Theorem for the polynomial algebra of operators for one copy of the group only for the case of $`G=SU\left(2\right)`$. In the second we use these results to extend the theorem to a certain class of operators which are not polynomial functions of the elementary operators and mix operators labelled by different edges by making an appeal to the moment problem by Hamburger for measures.
### 3.1 Polynomial Algebra of Operators over One Edge
As the problem is isomorphic for all the edges of a graph, we drop the label $`e`$ in this subsection and deal only with the operators $`\widehat{h}_{AB},\widehat{p}_j;A,B,C,..\{1/2,1/2\};j,k,l..=1,2,3`$ which obey the CCR algebra (3) with the label $`e=e^{}`$ dropped.
This subsection is divided into four parts. In the first we reduce the computation of expectation values of operator monomials to the computation of matrix elements of elementary operators between coherent states. In the second we estimate the matrix elements for the momentum operator and in the third for the holonomy operator. As expected the matrix elements are concentrated at $`g=g^{}`$ and simply given by the expectation value of the operator in question times the matrix element of the unit operator so that to leading order in $`t`$ the expectation value of the monomial is indeed the monomial of the expectation values.
The expectation values of operator monomials are computed in parts two and three by using the overcompleteness of the coherent states. This displays them as particularly useful in deriving a coherent states path integral formulation of the theory . In contrast, in the fourth part we use a different method based on an $`SL(2,\text{ }\mathrm{C})`$ operator identity and the so-called moment problem due to Hamburger which we deal with in detail in the second subsection.
#### 3.1.1 The Completeness Relation
Recall from for the case of a general compact Lie group or from for the special case of $`G=SU\left(2\right)`$ that the coherent states $`\psi _g^t`$ possess the following “reproducing property”
$$\left(\widehat{U}_t\psi \right)\left(g\right)=<\psi _g^{}^t,\psi >_{\mu _H}$$
(3.14)
where $`g^{}=\left(\overline{g}^T\right)^1`$ is the unique involution on $`G^{\text{ }\mathrm{C}}`$ that preserves $`G`$. Here $`\psi L_2(G,d\mu _H)=:`$ is an arbitrary state and $`\widehat{U}_t:^{\text{ }\mathrm{C}}=L_2(G^{\text{ }\mathrm{C}},d\nu _t)\text{Hol}\left(G^{\text{ }\mathrm{C}}\right)`$ the coherent state transform of , that is, the generalization of (3.3) to arbitrary states (heat kernel evolution followed by analytic continuation) mapping the $`\mu _H`$ square integrable functions on $`G`$ to holomorphic, $`\nu _t`$ square integrable functions on $`G^{\text{ }\mathrm{C}}`$. The measure $`\nu _t`$ of this Bargmann-Segal Hilbert space, defined generally in is computed explicitly in for the case of $`G=SU\left(2\right)`$ and is chosen such that $`\widehat{U}_t`$ is a unitary operator. Using this unitarity and the reproducing property we compute
$`<\psi ,\psi ^{}>_{\mu _H}=<\widehat{U}_t\psi ,\widehat{U}_t\psi ^{}>_{\nu _t}={\displaystyle _{G^{\text{ }\mathrm{C}}}}d\nu _t\left(g\right)\overline{\left(\widehat{U}_t\psi \right)\left(g\right)}\left(\widehat{U}_t\psi ^{}\right)\left(g\right)`$ (3.15)
$`=`$ $`{\displaystyle _{G^{\text{ }\mathrm{C}}}}d\nu _t\left(g\right)<\psi ,\psi _g^{}^t>_{\mu _H}<\psi _g^{}^t,\psi ^{}>_{\mu _H}`$
and using the involution invariance of the measure $`\nu _t`$ which is essentially a Gaussian in $`p_j`$ we find the completeness relation
$$_{G^{\text{ }\mathrm{C}}}𝑑\nu _t\left(g\right)|\psi _g^t><\psi _g^t|=1_{}$$
(3.16)
Suppose now that we are given an operator monomial $`\widehat{O}=\widehat{O}_1..\widehat{O}_n`$ where each of the $`\widehat{O}_k,k=1,..,n<\mathrm{}`$ stands for one of the elementary operators $`\widehat{h}_{AB},\widehat{p}_j`$. Then, using (3.16), we can write the expectation value of $`\widehat{O}`$ as
$`{\displaystyle \frac{<\psi _g^t,\widehat{O}\psi _g^t>}{\psi _g^t^2}}`$ (3.17)
$`=`$ $`{\displaystyle \frac{1}{\psi _g^t^2}}{\displaystyle _{G^{\text{ }\mathrm{C}}}}d\nu _t\left(g_1\right)..{\displaystyle _{G^{\text{ }\mathrm{C}}}}d\nu _t\left(g_{n1}\right){\displaystyle \underset{k=1}{\overset{n}{}}}[<\psi _{g_{k1}}^t,\widehat{O}_k\psi _{g_k}^t>]`$
$`=`$ $`{\displaystyle _{G^{\text{ }\mathrm{C}}}}d\mathrm{\Omega }\left(g_1\right)..{\displaystyle _{G^{\text{ }\mathrm{C}}}}d\mathrm{\Omega }\left(g_{n1}\right)\left({\displaystyle \underset{k=1}{\overset{n1}{}}}\left[\nu _t\left(g_k\right)\right|\left|\psi _{g_k}^t\right||^2]\right)\left({\displaystyle \underset{k=1}{\overset{n}{}}}\left[{\displaystyle \frac{<\psi _{g_{k1}}^t,\widehat{O}_k\psi _{g_k}^t>}{\psi _{g_{k1}}^t\psi _{g_k}^t}}\right]\right)`$
where $`\mathrm{\Omega }`$ is the Liouville measure on $`G^{\text{ }\mathrm{C}}T^{}G`$ and we have set $`g_0=g_n=g`$. Now we recall from that the quantity
$$j^t(g,g^{})=\frac{<\psi _g^t,\psi _g^{}^t>}{\psi _g^t\psi _g^{}^t}$$
(3.18)
is exponentially small (at least for $`G=SU\left(2\right)`$) in the sense of a Gaussian needle of width $`\sqrt{t}`$ unless $`g=g^{}`$ (where it equals unity of course) and that the quantity $`\nu _t\left(g\right)\psi _g^t^2`$ approaches exponentially fast the constant $`2/\left(\pi t^3\right)`$. Thus, it is conceivable that
$$\frac{<\psi _{g_{k1}}^t,\widehat{O}_k\psi _{g_k}^t>}{\psi _{g_{k1}}^t\psi _{g_k}^t}\frac{<\psi _{g_k}^t,\widehat{O}_k\psi _{g_k}^t>}{\psi _{g_k}^t^2}j^t(g_{k1},g_k)$$
(3.19)
If that would be the case then the product Liouville measure $`d\mathrm{\Omega }\left(g_1\right)..d\mathrm{\Omega }\left(g_{n1}\right)`$ would be essentially supported at $`g_1=..=g_{n1}=g`$ and we could pull the expectation values in (3.19) out of the integral (3.17) with $`g_k`$ replaced by $`g`$ and the remaining integral would then equal unity, of course. Thus we would have indeed shown that
$$\frac{<\psi _g^t,\widehat{O}\psi _g^t>}{\psi _g^t^2}\underset{k=1}{\overset{n}{}}\frac{<\psi _g^t,\widehat{O}_k\psi _g^t>}{\psi _g^t^2}$$
(3.20)
Thus, in order to prove the desired result (3.20) it is sufficient to prove (3.19) together with the precise meaning of $``$. The proof of (3.19) will also be the key ingredient for the derivation of path integrals based on the coherent states $`\psi _g^t`$ .
#### 3.1.2 Matrix Elements for the Momentum Operator
Recall from the beginning of this section that $`\widehat{p}_j=itX_j/2`$ where $`X_j\left(h\right)=\text{tr}\left(\left(\tau _jh\right)^T/h\right)`$ is a basis of right invariant vector fields on $`G`$ which generate left translations. Thus for any vector $`\psi `$ in the domain of $`\widehat{p}_j`$ (say smooth square integrable functions) we have
$$\left(\widehat{p}_j\psi \right)\left(h\right)=\frac{it}{2}\left(\frac{d}{ds}\right)_{s=0}\psi \left(e^{s\tau _j}h\right)$$
(3.21)
Since for our coherent states it holds that $`\psi _g^t\left(uh\right)=\psi _{u^1g}^t\left(h\right)`$ we have
$$\widehat{p}_j\psi _g^t=\frac{it}{2}\left(\frac{d}{ds}\right)_{s=0}\psi _{e^{s\tau _j}g}^t$$
(3.22)
It follows by explicit computation of the scalar product (see for details) that
$$<\psi _g^t,\widehat{p}_j\psi _g^{}^t>=\frac{it}{2}\left(\frac{d}{ds}\right)_{s=0}\psi _{e^{s\tau _j}g^{}\overline{g}^T}^{2t}\left(1_2\right)$$
(3.23)
Upon defining the complex number $`z`$ uniqely via (see for details)
$$\mathrm{cosh}\left(z\right)=\frac{1}{2}\text{tr}\left(e^{s\tau _j}g^{}\overline{g}^T\right)=\frac{1}{2}\left[\text{tr}\left(g^{}\overline{g}^T\right)s\text{tr}\left(\tau _jg^{}\overline{g}^T\right)\right]+O\left(s^2\right)$$
(3.24)
and using the Weyl character formula for $`G=SU\left(2\right)`$ we obtain with $`d_j=2j+1,j=0,1/2,1,3/2,..`$ the spin quantum numbers
$`<\psi _g^t,\widehat{p}_j\psi _g^{}^t>`$ $`=`$ $`{\displaystyle \frac{it}{2}}\left({\displaystyle \frac{d}{ds}}\right)_{s=0}{\displaystyle \underset{j}{}}d_je^{tj\left(j+1\right)}{\displaystyle \frac{\mathrm{sinh}\left(d_jz\right)}{\mathrm{sinh}\left(z\right)}}`$ (3.25)
$`=`$ $`{\displaystyle \frac{it}{2}}\left({\displaystyle \frac{d}{ds}}\right)_{s=0}{\displaystyle \frac{e^{t/4}}{\mathrm{sinh}\left(z\right)}}{\displaystyle \underset{j}{}}d_je^{td_j^2/4}\mathrm{sinh}\left(d_jz\right)`$
$`=`$ $`{\displaystyle \frac{it}{2}}\left({\displaystyle \frac{d}{ds}}\right)_{s=0}{\displaystyle \frac{e^{t/4}}{2\mathrm{sinh}\left(z\right)}}{\displaystyle \underset{n}{}}ne^{tn^2/4}e^{nz}`$
$`=`$ $`{\displaystyle \frac{it}{2}}\left({\displaystyle \frac{d}{ds}}\right)_{s=0}{\displaystyle \frac{e^{t/4}}{2\mathrm{sinh}\left(z\right)T}}{\displaystyle \underset{n}{}}\left(nT\right)e^{\left(nT\right)^2}e^{\left(nT\right)\xi }`$
where $`n\text{ }\mathrm{Z}`$ and we have defined $`T=\sqrt{t}/2,\xi =z/T`$. The Fourier transform $`\stackrel{~}{f}\left(k\right)=_{\text{ }\mathrm{R}}\frac{dx}{2\pi }e^{ikx}f\left(x\right)`$ of the function $`f\left(x\right):=xe^{x^2}e^{x\xi }`$ is given by
$$\stackrel{~}{f}\left(k\right)=\frac{\xi ik}{4\sqrt{\pi }}e^{\frac{\left(\xi ik\right)^2}{4}}$$
(3.26)
using a contour argument. An appeal to the Poisson summation formula (see, e.g., )
$$\underset{n}{}f\left(nT\right)=\frac{2\pi }{T}\underset{n}{}\stackrel{~}{f}\left(\frac{2\pi n}{T}\right)$$
(3.27)
then reveals
$$<\psi _g^t,\widehat{p}_j\psi _g^{}^t>=\frac{it}{2}\left(\frac{d}{ds}\right)_{s=0}\frac{\sqrt{\pi }e^{t/4}}{4\mathrm{sinh}\left(z\right)T^3}\underset{n}{}\left(z2\pi in\right)e^{\frac{\left(z2\pi in\right)^2}{t}}$$
(3.28)
By the very same methods we easily obtain
$`\psi _g^t^2`$ $`=`$ $`{\displaystyle \frac{\sqrt{\pi }e^{t/4}}{4\mathrm{sinh}\left(p\right)T^3}}{\displaystyle \underset{n}{}}\left(p2\pi in\right)e^{\frac{\left(p2\pi in\right)^2}{t}}`$
$`\psi _g^{}^t^2`$ $`=`$ $`{\displaystyle \frac{\sqrt{\pi }e^{t/4}}{4\mathrm{sinh}\left(p^{}\right)T^3}}{\displaystyle \underset{n}{}}\left(p^{}2\pi in\right)e^{\frac{\left(p^{}2\pi in\right)^2}{t}}`$ (3.29)
where $`\mathrm{cosh}\left(p\right)=\text{tr}\left(g\overline{g}^T\right)/2`$ with the polar decomposition $`g=e^{ip_j\tau _j/2}h,p=\sqrt{p_jp_j}`$ and similar for $`g^{}`$.
Let now $`\mathrm{cosh}\left(z_0\right):=\frac{1}{2}\text{tr}\left(g^{}\overline{g}^T\right)`$ and $`z=z_0+\delta `$ where $`\delta `$ is of first order in $`s`$. Then comparing
$$\mathrm{cosh}\left(z\right)=\mathrm{cosh}\left(z_0\right)\mathrm{cosh}\left(\delta \right)+\mathrm{sinh}\left(z_0\right)\mathrm{sinh}\left(\delta \right)=\mathrm{cosh}\left(z_0\right)+\delta \mathrm{sinh}\left(z_0\right)+O\left(\delta ^2\right)$$
(3.30)
with (3.24) we conclude that
$$\delta =s\frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{2\mathrm{sinh}\left(z_0\right)}$$
(3.31)
where the sign of $`\mathrm{sinh}\left(z_0\right)`$ follows from the formulas given in . It follows that
$$\left(\frac{d}{ds}\right)_{s=0}=\frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{2\mathrm{sinh}\left(z_0\right)}\left(\frac{d}{dz}\right)_{z=z_0}$$
(3.32)
and performing the derivative in (3.28) we end up with
$`<\psi _g^t,\widehat{p}_j\psi _g^{}^t>={\displaystyle \frac{it}{4}}{\displaystyle \frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{\mathrm{sinh}\left(z_0\right)}}{\displaystyle \frac{\sqrt{\pi }e^{t/4}}{4T^3}}\times `$ (3.33)
$`\times `$ $`{\displaystyle \underset{n}{}}e^{\frac{\left(z_02\pi in\right)^2}{t}}\left[2{\displaystyle \frac{\left(z_02\pi in\right)^2}{t\mathrm{sinh}\left(z_0\right)}}\left(z_02\pi in\right){\displaystyle \frac{\mathrm{cosh}\left(z_0\right)}{\mathrm{sinh}^2\left(z_0\right)}}+{\displaystyle \frac{1}{\mathrm{sinh}\left(z_0\right)}}\right]`$
Combination with (3.1.2) results in the exact expression for the momentum matrix element
$`<\widehat{p}_j>_{gg^{}}^t`$ $`:=`$ $`{\displaystyle \frac{<\psi _g^t,\widehat{p}_j\psi _g^{}^t>}{\psi _g^t\psi _g^{}^t}}`$
$`=`$ $`\left[{\displaystyle \frac{iz_0}{2}}{\displaystyle \frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{\mathrm{sinh}\left(z_0\right)}}\right]\times `$
$`\times `$ $`{\displaystyle \frac{\left\{\frac{t}{2z_0}_ne^{\frac{\left(z_02\pi in\right)^2p^2/2\left(p^{}\right)^2/2}{t}}\left[2\frac{\left(z_02\pi in\right)^2}{t\mathrm{sinh}\left(z_0\right)}\left(z_02\pi in\right)\frac{\mathrm{cosh}\left(z_0\right)}{\mathrm{sinh}^2\left(z_0\right)}+\frac{1}{\mathrm{sinh}\left(z_0\right)}\right]\right\}}{\left[_n\frac{p2\pi in}{\mathrm{sinh}\left(p\right)}e^{\frac{4\pi ^2n^2}{t}}e^{i\frac{4\pi np}{t}}\right]^{1/2}\left[_n\frac{p^{}2\pi in}{\mathrm{sinh}\left(p^{}\right)}e^{\frac{4\pi ^2n^2}{t}}e^{i\frac{4\pi np^{}}{t}}\right]^{1/2}}}`$
The arguments $`D^t\left(p\right),D^t\left(p^{}\right)`$ of the square roots in the denominator of (3.1.2) were already estimated in and we will only recall the result here without derivation
$$\frac{p}{\mathrm{sinh}\left(p\right)}\left(1K_t\right)D^t\left(p\right)\frac{p}{\mathrm{sinh}\left(p\right)}\left(1+K_t\right)\text{ where }K_t=2\underset{n=1}{\overset{\mathrm{}}{}}e^{\frac{4\pi ^2n^2}{t}}\left(1+\frac{8\pi ^2n}{t}\right)$$
(3.35)
vanishes exponentially fast with $`t0`$ and similar for $`D^t\left(p^{}\right)`$ with $`p`$ replaced by $`p^{}`$. The term in the curly brackets of the numerator in (3.1.2) can be more explicitly written as
$`N^t(g,g^{})=e^{\frac{z_0^2p^2/2\left(p^{}\right)^2}{t}}\{[{\displaystyle \frac{z_0}{\mathrm{sinh}\left(z_0\right)}}{\displaystyle \frac{t}{2}}{\displaystyle \frac{\mathrm{cosh}\left(z_0\right)}{\mathrm{sinh}^2\left(z_0\right)}}+{\displaystyle \frac{1}{z_0\mathrm{sinh}\left(z_0\right)}}]`$
$`+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{4\pi ^2n^2}{t}}[({\displaystyle \frac{2z_0^28\pi ^2n^2}{z_0\mathrm{sinh}\left(z_0\right)}}{\displaystyle \frac{t\mathrm{cosh}\left(z_0\right)}{\mathrm{sinh}^2\left(z_0\right)}}+{\displaystyle \frac{t}{z_0\mathrm{sinh}\left(z_0\right)}})\mathrm{cos}\left({\displaystyle \frac{4\pi nz_0}{t}}\right)`$
$`+2\pi nt({\displaystyle \frac{\mathrm{cosh}\left(z_0\right)}{z_0\mathrm{sinh}^2\left(z_0\right)}}{\displaystyle \frac{4}{t\mathrm{sinh}\left(z_0\right)}})\mathrm{sin}\left({\displaystyle \frac{4\pi nz_0}{t}}\right)]\}`$ (3.36)
which is superficially divergent at the points $`z_0=0,i\pi `$ which reminds us, of course, of the singularity structure of the overlap function in and we will proceed similarly to estimate (3.1.2). Thus, we will separate the discussion into cases A) and B) respectively, writing $`N^t(g,g^{})`$ in terms of $`z_0`$ and $`z_0^{}=z_0i\pi `$ respectively for $`0\varphi \left(1c\right)\pi `$ and $`\left(1c\right)\pi \varphi \pi `$ respectively where $`0<c<1`$ is a constant. As shown in , $`z_0=s+i\varphi `$ is always uniquely determined with $`s\text{ }\mathrm{R},\varphi [0,\pi ]`$.
Case A)
Let us pull out a factor of $`z_0/\mathrm{sinh}\left(z_0\right)`$ from (3.1.2) since at $`g=g^{}`$ it will cancel against the $`p/\mathrm{sinh}\left(p\right)`$ coming from $`D^t\left(p\right)`$ and separate terms into those which are regular and irregular respectively at $`z_0=0`$, resulting in
$`N^t(g,g^{})=e^{\frac{z_0^2p^2/2\left(p^{}\right)^2}{t}}{\displaystyle \frac{z_0}{\mathrm{sinh}\left(z_0\right)}}\{\left[1\right]+\left[{\displaystyle \frac{t}{2}}{\displaystyle \frac{\frac{\mathrm{sinh}\left(z_0\right)}{z_0}\mathrm{cosh}\left(z_0\right)}{z_0\mathrm{sinh}\left(z_0\right)}}\right]`$
$`+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{4\pi ^2n^2}{t}}(\left[2\mathrm{cos}\left({\displaystyle \frac{4\pi nz_0}{t}}\right)\right]+\left[t{\displaystyle \frac{\frac{\mathrm{sinh}\left(z_0\right)}{z_0}\mathrm{cosh}\left(z_0\right)}{z_0\mathrm{sinh}\left(z_0\right)}}\mathrm{cos}\left({\displaystyle \frac{4\pi nz_0}{t}}\right)\right][8\pi n\mathrm{sin}\left({\displaystyle \frac{4\pi nz_0}{t}}\right)/z_0]`$
$`+2\pi nt[{\displaystyle \frac{\mathrm{cosh}\left(z_0\right)}{z_0^2\mathrm{sinh}\left(z_0\right)}}\mathrm{sin}\left({\displaystyle \frac{4\pi nz_0}{t}}\right){\displaystyle \frac{4\pi n}{tz_0^2}}\mathrm{cos}\left({\displaystyle \frac{4\pi nz_0}{t}}\right)])\}`$ (3.37)
and obviously all terms in the square brackets, except for the last one, are regular at $`z_0=0`$. However, expanding numerator and denominator to second order in $`z_0`$ we see
$`{\displaystyle \frac{\mathrm{cosh}\left(z_0\right)}{z_0^2\mathrm{sinh}\left(z_0\right)}}\mathrm{sin}\left({\displaystyle \frac{4\pi nz_0}{t}}\right){\displaystyle \frac{4\pi n}{tz_0^2}}\mathrm{cos}\left({\displaystyle \frac{4\pi nz_0}{t}}\right)`$ (3.38)
$`=`$ $`{\displaystyle \frac{1}{z_0\mathrm{sinh}\left(z_0\right)}}\left[\mathrm{cosh}\left(z_0\right){\displaystyle \frac{\mathrm{sin}\left(\frac{4\pi nz_0}{t}\right)}{z_0}}{\displaystyle \frac{4\pi n}{t}}{\displaystyle \frac{\mathrm{sinh}\left(z_0\right)}{z_0}}\mathrm{cos}\left({\displaystyle \frac{4\pi nz_0}{t}}\right)\right]`$
$`=`$ $`{\displaystyle \frac{1}{z_0^2\left(1+O\left(z_0^2\right)\right)}}\left[{\displaystyle \frac{4\pi n}{t}}\left(1+O\left(z_0^2\right)\right){\displaystyle \frac{4\pi n}{t}}\left(1+O\left(z_0^2\right)\right)\right]=O\left(z_0^2\right)`$
so that (3.38) even vanishes at $`z_0=0`$.
We want to put bounds on all those terms in the square brackets of (3.1.2) for $`0\varphi \left(1c\right)\pi `$ except for the first one and, in particular, estimate the series. To that end we write (3.1.2) in a yet more suggestive form
$`N^t(g,g^{})`$ $`=`$ $`e^{\frac{z_0^2p^2/2\left(p^{}\right)^2/2}{t}}{\displaystyle \frac{z_0}{\mathrm{sinh}\left(z_0\right)}}\times `$ (3.39)
$`\times `$ $`\{\left[1\right]+\left[{\displaystyle \frac{t}{2}}{\displaystyle \frac{\frac{\mathrm{sinh}\left(z_0\right)}{z_0}\mathrm{cosh}\left(z_0\right)}{z_0\mathrm{sinh}\left(z_0\right)}}\right][1+2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{4\pi ^2n^2}{t}}\mathrm{cos}\left({\displaystyle \frac{4\pi nz_0}{t}}\right)]`$
$`+2\left[{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{4\pi ^2n^2}{t}}\left(\mathrm{cos}\left({\displaystyle \frac{4\pi nz_0}{t}}\right){\displaystyle \frac{\left(4\pi n\right)^2}{t}}{\displaystyle \frac{\mathrm{sin}\left(\frac{4\pi nz_0}{t}\right)}{\frac{4\pi nz_0}{t}}}\right)\right]`$
$`+2\pi t\left[{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{4\pi ^2n^2}{t}}n({\displaystyle \frac{\mathrm{cosh}\left(z_0\right)}{z_0\mathrm{sinh}\left(z_0\right)}}{\displaystyle \frac{\mathrm{sin}\left(\frac{4\pi nz_0}{t}\right)}{z_0}}{\displaystyle \frac{4\pi n}{tz_0^2}}\mathrm{cos}\left({\displaystyle \frac{4\pi nz_0}{t}}\right))\right]\}`$
$`=:`$ $`e^{\frac{z_0^2p^2/2\left(p^{}\right)^2/2}{t}}{\displaystyle \frac{z_0}{\mathrm{sinh}\left(z_0\right)}}\times `$
$`\times `$ $`\left\{1+{\displaystyle \frac{t}{2}}{\displaystyle \frac{\frac{\mathrm{sinh}\left(z_0\right)}{z_0}\mathrm{cosh}\left(z_0\right)}{z_0\mathrm{sinh}\left(z_0\right)}}I_1+2I_2+2\pi tI_3\right\}`$
and it remains to estimate the sums $`I_1,I_2,I_3`$ as well as the expression
$$I:=\frac{\frac{\mathrm{sinh}\left(z_0\right)}{z_0}\mathrm{cosh}\left(z_0\right)}{z_0\mathrm{sinh}\left(z_0\right)}$$
(3.40)
Notice that for the terms proportional to $`t`$ in (3.39) it will be sufficient to estimate them by a function integrable against the Gaussian prefactor of (3.39).
Focussing first on (3.40) we first prove two elementary lemmas.
###### Lemma 3.1
For any $`z=s+i\varphi \text{ }\mathrm{C}`$ such that $`0\varphi (1c)\pi `$ for some $`0<c<1`$ we have
$`{\displaystyle \frac{s^2}{\mathrm{sinh}^2\left(s\right)}}|{\displaystyle \frac{z}{\mathrm{sinh}\left(z\right)}}|^2{\displaystyle \frac{\varphi ^2}{\mathrm{sin}^2\left(\varphi \right)}}\left[{\displaystyle \frac{\pi \left(1c\right)}{\mathrm{sin}\left(\pi \left(1c\right)\right)}}\right]^2=:k_c^2`$ (3.41)
$`\left|{\displaystyle \frac{z}{\mathrm{sinh}\left(z\right)}}\right|^2{\displaystyle \frac{\varphi ^2}{\mathrm{sin}^2\left(\varphi \right)}}{\displaystyle \frac{s^2}{\mathrm{sinh}^2\left(s\right)}}k_c^2{\displaystyle \frac{s^2}{\mathrm{sinh}^2\left(s\right)}}`$ (3.42)
Proof of Lemma 3.1 :
Notice the identity
$$\left|\frac{z}{\mathrm{sinh}\left(z\right)}\right|^2=\frac{s^2+\varphi ^2}{\mathrm{sinh}^2\left(s\right)\mathrm{cos}^2\left(\varphi \right)+\mathrm{cosh}^2\left(s\right)\mathrm{sin}^2\left(\varphi \right)}=\frac{s^2+\varphi ^2}{\mathrm{sinh}^2\left(s\right)+\mathrm{sin}^2\left(\varphi \right)}$$
(3.43)
Both the lower and upper bounds in (3.41) turn out to be equivalent with the inequality
$$\left(\frac{\mathrm{sinh}\left(s\right)}{s}\right)^2\left(\frac{\mathrm{sin}\left(\varphi \right)}{\varphi }\right)^2$$
(3.44)
which is true as the left hand side and right hand side respectively both take its minimum and maximum respectively at $`s=\varphi =0`$, in fact, the left hand side and right hand side respectively are strictly increasing and decreasing functions respectively.
Inequality (3.42) in turn can be transformed into the equivalent form
$$\frac{1}{s^2}\frac{1}{\mathrm{sinh}^2\left(s\right)}\frac{1}{\mathrm{sin}^2\left(\varphi \right)}\frac{1}{\varphi ^2}$$
(3.45)
Both the left and the right hand side of this inequality are always positive in the range considered and in fact the left hand side approches $`0`$ as $`s\mathrm{}`$ while the right hand side approaches $`+\mathrm{}`$ as $`\varphi \pi `$. At $`\varphi =s=0`$ both sides equal $`1/3`$. We will in fact prove that
$$\frac{1}{s^2}\frac{1}{\mathrm{sinh}^2\left(s\right)}\frac{1}{3}\frac{1}{\mathrm{sin}^2\left(\varphi \right)}\frac{1}{\varphi ^2}$$
(3.46)
Consider first the left hand side. This can be written in the equivalent form
$$\left(3s^2\right)\mathrm{sinh}^2\left(s\right)3s^2$$
(3.47)
which is obviously true for $`s^23`$ so that we may restrict examination to $`s^2<3`$. In that case we may write (3.47) in the equivalent form
$$\frac{\mathrm{cosh}\left(2s\right)1}{2s^2}\frac{1}{1s^2/3}$$
(3.48)
As $`s^2/3<1`$ the right hand side can be expanded into a geometric series. Introducing $`x=2s`$ and employing the Taylor series for $`\mathrm{cosh}`$ we can write (3.48) as
$$\underset{n=0}{\overset{\mathrm{}}{}}x^{2n}\left[\frac{1}{12^n}\frac{2}{\left(2\left(n+1\right)\right)!}\right]0$$
(3.49)
and it will be sufficient to establish non-negativity of every coefficient. Using the basic estimate $`\mathrm{ln}\left(n!\right)n\left(\mathrm{ln}\left(n\right)1\right)+1`$ valid for $`n1`$ it is easy to see that this is indeed the case for $`n>4`$ while for the cases $`n=0,1,2,3,4`$ this can be checked by direct computation.
Turning to the right hand side of (3.46) we can write it in the equivalent form
$$3\varphi ^2\mathrm{sin}^2\left(\varphi \right)\left[3+\varphi ^2\right]$$
(3.50)
which due to $`\mathrm{sin}^2\left(\varphi \right)1`$ is certainly true for $`\varphi ^23/2`$ so that we can focus attention on the case $`\varphi ^2/3<3/2`$. Writing (3.50) in the equivalent form
$$\frac{1\mathrm{cos}\left(2\varphi \right)}{2\varphi ^2}\frac{1}{1\left(\varphi ^2/3\right)}$$
(3.51)
exploiting that $`0\varphi ^2/3<1/2`$ lies in the radius of convergence of the geometric series, introducing $`0y=\left(2\varphi \right)^2/6<1`$ and
$$b_n:=\frac{1}{2^n}\frac{26^n}{\left(2\left(n+1\right)\right)!}$$
(3.52)
we may write (3.51) in the form
$$\underset{n=0}{\overset{\mathrm{}}{}}\left(1\right)^ny^nb_n=\underset{n=0}{\overset{\mathrm{}}{}}\left[y^{2n}b_{2n}y^{2n+1}b_{2n+1}\right]0$$
(3.53)
Since $`0y<1`$ this will be the case if $`b_{2n}b_{2n+1}`$ for all $`n0`$. As one can check, $`b_0=b_1=0`$ and in (3.49) we have already seen that $`b_n>0`$ for $`n2`$. Thus, it is enough to prove $`b_{2n}b_{2n+1}`$ for $`n1`$. In fact, we will prove more, namely that $`b_n`$ is strictly decreasing for $`n2`$. This turns out to be equivalent with $`\left(2\left(n+1\right)\right)!212^n>1`$ for all $`n2`$ which in turn would follow from $`\left(2\left(n+1\right)\right)!>\left(2+\frac{1}{144}\right)12^n`$. The latter condition can be demonstrated to be true by methods similar to those outlined in (3.49).
$`\mathrm{}`$
###### Lemma 3.2
Let
$$S\left(z\right):=\frac{\mathrm{sinh}\left(z\right)z}{z^2\mathrm{sinh}\left(z\right)},C\left(z\right):=\frac{\mathrm{cosh}\left(z\right)1}{z\mathrm{sinh}\left(z\right)},k_c^{}:=\sqrt{1+k_c\mathrm{cosh}\left(\pi \left(1c\right)\right)}$$
(3.54)
Then for any $`z=s+i\varphi \text{ }\mathrm{C},\mathrm{\hspace{0.33em}0}\varphi (1c)\pi `$ it holds that
$$\left|S\left(z\right)\right|\sqrt{2}k_c^{}\text{ and }\left|C\left(z\right)\right|\sqrt{2}k_c^{}$$
(3.55)
Proof of Lemma 3.2 :
Using hyperbolic and trigonometric identities we derive
$`\left|S\left(z\right)\right|^2`$ $`=`$ $`{\displaystyle \frac{\left(\mathrm{sinh}\left(s\right)\mathrm{cos}\left(\varphi \right)s\right)^2+\left(\mathrm{cosh}\left(s\right)\mathrm{sin}\left(\varphi \right)\varphi \right)^2}{\left|z\right|^4\left(\mathrm{cosh}^2\left(s\right)\mathrm{cos}^2\left(\varphi \right)\right)}}`$ (3.56)
$`=`$ $`{\displaystyle \frac{\left(\frac{s^2}{\left|z\right|^2}\frac{\mathrm{sinh}\left(s\right)s}{s^2}\mathrm{cos}\left(\varphi \right)+\varphi \frac{s\varphi }{\left|z\right|^2}\frac{\mathrm{cos}\left(\varphi \right)1}{\varphi ^2}\right)^2+\left(\frac{s\varphi }{\left|z\right|^2}\frac{\mathrm{cosh}\left(s\right)1}{s}\frac{\mathrm{sin}\left(\varphi \right)}{\varphi }+\frac{\varphi ^2}{\left|z\right|^2}\varphi \frac{\mathrm{sin}\left(\varphi \right)\varphi }{\varphi ^3}\right)^2}{\mathrm{cosh}^2\left(s\right)\mathrm{cos}^2\left(\varphi \right)}}`$
$``$ $`{\displaystyle \frac{\left(\left|\frac{\mathrm{sinh}\left(s\right)s}{s^2}\right|\left|\mathrm{cos}\left(\varphi \right)\right|+\varphi \left|\frac{\mathrm{cos}\left(\varphi \right)1}{\varphi ^2}\right|\right)^2+\left(\left|\frac{\mathrm{cosh}\left(s\right)1}{s}\right|\left|\frac{\mathrm{sin}\left(\varphi \right)}{\varphi }\right|+\varphi \left|\frac{\mathrm{sin}\left(\varphi \right)\varphi }{\varphi ^3}\right|\right)^2}{\mathrm{cosh}^2\left(s\right)\mathrm{cos}^2\left(\varphi \right)}}`$
$``$ $`{\displaystyle \frac{\left(\left|\frac{\mathrm{sinh}\left(s\right)s}{s^2}\right|+\varphi \left|\frac{\mathrm{cos}\left(\varphi \right)1}{\varphi ^2}\right|\right)^2+\left(\left|\frac{\mathrm{cosh}\left(s\right)1}{s}\right|+\varphi \left|\frac{\mathrm{sin}\left(\varphi \right)\varphi }{\varphi ^3}\right|\right)^2}{\mathrm{cosh}^2\left(s\right)\mathrm{cos}^2\left(\varphi \right)}}`$
$``$ $`\left(\left|{\displaystyle \frac{\mathrm{sinh}\left(s\right)s}{s^2\mathrm{sinh}\left(s\right)}}\right|+\left|{\displaystyle \frac{\varphi }{\mathrm{sin}\left(\varphi \right)}}\right|\left|{\displaystyle \frac{\mathrm{cos}\left(\varphi \right)1}{\varphi ^2}}\right|\right)^2+\left(\left|{\displaystyle \frac{\mathrm{cosh}\left(s\right)1}{s\mathrm{sinh}\left(s\right)}}\right|+\left|{\displaystyle \frac{\varphi }{\mathrm{sin}\left(\varphi \right)}}\right|\left|{\displaystyle \frac{\mathrm{sin}\left(\varphi \right)\varphi }{\varphi ^3}}\right|\right)^2`$
$``$ $`\left(\left|S\left(s\right)\right|+k_c\left|{\displaystyle \frac{\mathrm{cos}\left(\varphi \right)1}{\varphi ^2}}\right|\right)^2+\left(\left|C\left(s\right)\right|+k_c\left|{\displaystyle \frac{\mathrm{sin}\left(\varphi \right)\varphi }{\varphi ^3}}\right|\right)^2`$
where in the first inequality we used $`\left|s\right|,\left|\varphi \right|\left|z\right|`$, in the second that $`\left|\mathrm{cos}\left(\varphi \right)\right|,\left|\mathrm{sin}\left(\varphi \right)/\varphi \right|1`$, in the third that $`\mathrm{cosh}^2\left(s\right)\mathrm{cos}^2\left(\varphi \right)\mathrm{sinh}^2\left(s\right),\mathrm{sin}^2\left(\varphi \right)`$ and in the fourth we used Lemma 3.1 and again the definition of $`S\left(z\right),C\left(z\right)`$.
Proceeding similarly with $`C\left(z\right)`$ we arrive at
$`\left|C\left(z\right)\right|^2`$ $`=`$ $`{\displaystyle \frac{\left(\mathrm{cosh}\left(s\right)\mathrm{cos}\left(\varphi \right)1\right)^2+\left(\mathrm{sinh}\left(s\right)\mathrm{sin}\left(\varphi \right)\right)^2}{\left|z\right|^2\left(\mathrm{cosh}^2\left(s\right)\mathrm{cos}^2\left(\varphi \right)\right)}}`$ (3.57)
$`=`$ $`{\displaystyle \frac{\left(s\frac{\mathrm{cosh}\left(s\right)1}{s}\mathrm{cos}\left(\varphi \right)+\varphi ^2\frac{\mathrm{cos}\left(\varphi \right)1}{\varphi ^2}\right)^2+\left(\varphi \mathrm{sinh}\left(s\right)\frac{\mathrm{sin}\left(\varphi \right)}{\varphi }\right)^2}{\left|z\right|^2\left(\mathrm{cosh}^2\left(s\right)\mathrm{cos}^2\left(\varphi \right)\right)}}`$
$``$ $`{\displaystyle \frac{\left(\left|\frac{\mathrm{cosh}\left(s\right)1}{s}\right|\left|\mathrm{cos}\left(\varphi \right)\right|+\varphi \left|\frac{\mathrm{cos}\left(\varphi \right)1}{\varphi ^2}\right|\right)^2+\left(\mathrm{sinh}\left(s\right)\left|\frac{\mathrm{sin}\left(\varphi \right)}{\varphi }\right|\right)^2}{\left(\mathrm{cosh}^2\left(s\right)\mathrm{cos}^2\left(\varphi \right)\right)}}`$
$``$ $`{\displaystyle \frac{\left(\left|\frac{\mathrm{cosh}\left(s\right)1}{s}\right|+\varphi \left|\frac{\mathrm{cos}\left(\varphi \right)1}{\varphi ^2}\right|\right)^2+\mathrm{sinh}^2\left(s\right)}{\left(\mathrm{cosh}^2\left(s\right)\mathrm{cos}^2\left(\varphi \right)\right)}}`$
$``$ $`\left(\left|{\displaystyle \frac{\mathrm{cosh}\left(s\right)1}{s\mathrm{sinh}\left(s\right)}}\right|+\left|{\displaystyle \frac{\varphi }{\mathrm{sin}\left(\varphi \right)}}\right|\left|{\displaystyle \frac{\mathrm{cos}\left(\varphi \right)1}{\varphi ^2}}\right|\right)^2+1`$
$``$ $`\left(\left|C\left(s\right)\right|+k_c\left|{\displaystyle \frac{\mathrm{cos}\left(\varphi \right)1}{\varphi ^2}}\right|\right)^2+1`$
Using the Taylor series expansion of the trigonometric functions we easily obtain
$`\left|{\displaystyle \frac{\mathrm{cos}\left(\varphi \right)1}{\varphi ^2}}\right|=\left|{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\left(1\right)^n\varphi ^{2\left(n1\right)}}{\left(2n\right)!}}\right|=\left|{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\left(1\right)^n\varphi ^{2n}}{\left(2n+2\right)!}}\right|{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\varphi ^{2n}}{\left(2\left(n+1\right)\right)!}}`$ (3.58)
$``$ $`\mathrm{cosh}\left(\varphi \right)\mathrm{cosh}\left(\pi \left(1c\right)\right)`$
$`\left|{\displaystyle \frac{\mathrm{sin}\left(\varphi \right)\varphi }{\varphi ^3}}\right|=\left|{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\left(1\right)^n\varphi ^{2\left(n1\right)}}{\left(2n+1\right)!}}\right|=\left|{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\left(1\right)^n\varphi ^{2n}}{\left(2n+3\right)!}}\right|{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\varphi ^{2n}}{\left(2n+3\right)!}}\mathrm{cosh}\left(\varphi \right)`$
$``$ $`\mathrm{cosh}\left(\pi \left(1c\right)\right)`$
Next notice that $`S\left(s\right)=S\left(\left|s\right|\right),C\left(s\right)=C\left(\left|s\right|\right)`$ so that we have reduced our estimate for $`\left|S\left(z\right)\right|,\left|C\left(z\right)\right|`$ to that for real non-negative arguments. Finally, using the Taylor series expression for the hyperbolic functions we find
$`\left|{\displaystyle \frac{\mathrm{cosh}\left(s\right)1}{s}}\right|=\left|{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{s^{2n1}}{\left(2n\right)!}}\right|=\left|{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{s^{2n+1}}{\left(2n+2\right)!}}\right|\mathrm{sinh}\left(\left|s\right|\right)`$
$`\left|{\displaystyle \frac{\mathrm{sinh}\left(s\right)s}{s^2}}\right|=\left|{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\varphi ^{2n1}}{\left(2n+1\right)!}}\right|=\left|{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{s^{2n+1}}{\left(2n+3\right)!}}\right|\mathrm{sinh}\left(\left|s\right|\right)`$ (3.59)
so that in fact $`\left|S\left(s\right)\right|,\left|C\left(s\right)\right|1`$. Together with the trivial inequality $`1k_c`$ we thus obtain indeed $`\left|C\left(z\right)\right|^2,\left|S\left(z\right)\right|^22\left[1+k_c\mathrm{cosh}\left(\pi \left(1c\right)\right)\right]^2=2\left(k_c^{}\right)^2`$.
$`\mathrm{}`$
Now we can give a bound on $`I`$.
###### Lemma 3.3
For any $`z_0=s+i\varphi \text{ }\mathrm{C}`$ such that $`0\varphi (1c)\pi `$ for some $`0<c<1`$ we have
$$\left|I\right|2k_c^{}$$
(3.60)
Proof of Lemma 3.3 :
This follows immediately from the identity
$$I=S\left(z\right)C\left(z\right)$$
(3.61)
and lemma 3.2.
$`\mathrm{}`$
Next consider the following expression that appears in $`I_3`$
$$J:=\frac{\mathrm{cosh}\left(z_0\right)}{z_0\mathrm{sinh}\left(z_0\right)}\frac{\mathrm{sin}\left(\frac{4\pi nz_0}{t}\right)}{z_0}\frac{4\pi n}{tz_0^2}\mathrm{cos}\left(\frac{4\pi nz_0}{t}\right)$$
(3.62)
which we must estimate in such a way that finally $`s`$ and $`n`$ do not appear in the combination $`ns`$ inside a hyperbolic function as otherwise we must worry about convergence of the series $`I_3`$ as $`\left|s\right|`$ becomes arbitrarily large in the integrals we are considering. At the same time we must bound the superficial singularity at $`z_0=0`$. To that end we introduce $`z_0^{}=4\pi nz_0/t`$ and notice the identity, recalling the definition (3.40)
$$J=\frac{4\pi n}{t}\left\{\frac{\mathrm{sin}\left(z_0^{}\right)}{z_0^{}}I+\left(\frac{4\pi n}{t}\right)^2\left[\frac{\mathrm{sin}\left(z_0^{}\right)z_0^{}}{\left(z_0^{}\right)^3}\frac{\mathrm{cos}\left(z_0^{}\right)1}{\left(z_0^{}\right)^2}\right]\right\}$$
(3.63)
and the task is to estimate the three terms of the form $`\left(\mathrm{cos}\left(z\right)1\right)/z^2,\mathrm{sin}\left(z\right)/z,\left(\mathrm{sin}\left(z\right)z\right)/z^3`$ for arbitrary $`z=x+iy\text{ }\mathrm{C}`$. The inequality
$$\left|\frac{\mathrm{sin}\left(z\right)}{z}\right|2\mathrm{cosh}\left(y\right)$$
(3.64)
is the content of lemma 4.1 of . The remaining two estimates are the hardest ones and we have therefore devoted the subsequent lemma to them.
###### Lemma 3.4
Let for any complex number $`z`$
$$c\left(z\right):=\frac{\mathrm{cos}\left(z\right)1}{z^2}\text{ and }s\left(z\right):=\frac{\mathrm{sin}\left(z\right)z}{z^3}$$
(3.65)
Then
$$\left|c\left(z\right)\right|4\mathrm{cosh}\left(\mathrm{}\left(z\right)\right)\text{ and }\left|s\left(z\right)\right|4\mathrm{cosh}\left(\mathrm{}\left(z\right)\right)$$
(3.66)
Proof of Lemma 3.4 :
i)
Splitting $`z=x+iy`$ we have
$`\left|{\displaystyle \frac{\mathrm{cos}\left(z\right)1}{z^2}}\right|=\left|{\displaystyle \frac{\left[\mathrm{cos}\left(x\right)\mathrm{cosh}\left(y\right)1\right]i\left[\mathrm{sin}\left(x\right)\mathrm{sinh}\left(y\right)\right]}{z^2}}\right|`$ (3.67)
$``$ $`\left|{\displaystyle \frac{xy}{z^2}}\right|\left|{\displaystyle \frac{\mathrm{sin}\left(x\right)}{x}}\right|\left|{\displaystyle \frac{\mathrm{sinh}\left(y\right)}{y}}\right|+\left|{\displaystyle \frac{\left[\mathrm{cos}\left(x\right)1\right]\left[\mathrm{cosh}\left(y\right)1\right]+\left[\mathrm{cos}\left(x\right)1\right]+\left[\mathrm{cosh}\left(y\right)1\right]}{z^2}}\right|`$
$``$ $`\mathrm{sinh}\left(\left|y\right|\right)+\left|{\displaystyle \frac{xy}{z^2}}\right|\left|{\displaystyle \frac{\mathrm{cosh}\left(y\right)1}{y}}\right|\left|{\displaystyle \frac{\mathrm{cos}\left(x\right)1}{x}}\right|+\left|{\displaystyle \frac{x^2}{z^2}}\right|\left|{\displaystyle \frac{\mathrm{cos}\left(x\right)1}{x^2}}\right|+\left|{\displaystyle \frac{y^2}{z^2}}\right|\left|{\displaystyle \frac{\mathrm{cosh}\left(y\right)1}{y^2}}\right|`$
$``$ $`\mathrm{sinh}\left(\left|y\right|\right)+\left|{\displaystyle \frac{\mathrm{cosh}\left(y\right)1}{y}}\right|\left|{\displaystyle \frac{\mathrm{cos}\left(x\right)1}{x}}\right|+\left|{\displaystyle \frac{\mathrm{cos}\left(x\right)1}{x^2}}\right|+\left|{\displaystyle \frac{\mathrm{cosh}\left(y\right)1}{y^2}}\right|`$
using $`\left|x/z\right|,\left|y/z\right|1`$. It is easy to see that $`\left|\mathrm{cosh}\left(y\right)1\right|/y^2\mathrm{cosh}\left(y\right)`$ by using the Taylor series of $`\mathrm{cosh}\left(y\right)`$, see e.g. . By similar methods it is easy to establish that $`\left|\frac{\mathrm{cosh}\left(y\right)1}{y}\right|\mathrm{sinh}\left(\left|y\right|\right)`$. Furthermore, we claim that $`\left|\frac{\mathrm{cos}\left(x\right)1}{x}\right|,\left|\frac{\mathrm{cos}\left(x\right)1}{x^2}\right|1`$.
To see the former, notice that $`\left|\frac{\mathrm{cos}\left(x\right)1}{x}\right|=\left|\frac{\mathrm{cos}\left(\left|x\right|\right)1}{\left|x\right|}\right|`$ so it will be sufficient to demonstrate this for $`x0`$. Now for $`x0`$ the inequality $`\left|\frac{\mathrm{cos}\left(x\right)1}{x}\right|1`$ is equivalent with $`1x\mathrm{cos}\left(x\right)1+x`$, so we claim that $`f_\pm \left(x\right)=x\pm \left(1\mathrm{cos}\left(x\right)\right)`$ are not negative functions for $`x0`$. But $`f_\pm ^{}\left(x\right)=1\pm \mathrm{sin}\left(x\right)0`$ so that $`f_\pm `$ is never decreasing and takes its minimum at $`x=0`$ where $`f_\pm \left(0\right)=0`$.
To see the latter, notice that $`\left|\frac{\mathrm{cos}\left(x\right)1}{x^2}\right|=\left|\frac{\mathrm{cos}\left(\left|x\right|\right)1}{\left|x\right|^2}\right|`$ so it will be sufficient to demonstrate this for $`x0`$ as well. Now for $`x0`$ the inequality $`\left|\frac{\mathrm{cos}\left(x\right)1}{x^2}\right|1`$ is equivalent with $`1x^2\mathrm{cos}\left(x\right)1+x^2`$, so we claim that $`f_\pm \left(x\right)=x^2\pm \left(1\mathrm{cos}\left(x\right)\right)`$ are not negative functions for $`x0`$. But $`g_\pm =f_\pm ^{}\left(x\right)=2x\pm \mathrm{sin}\left(x\right)`$ and $`g_\pm ^{}\left(x\right)=2\pm \mathrm{cos}\left(x\right)>0`$. Thus, $`g_\pm `$ is strictly increasing, its minimum at $`x=0`$ being $`g_\pm \left(0\right)=0`$. Thus $`f_\pm `$ is never decreasing and takes its minimum at $`x=0`$ where $`f_\pm \left(0\right)=0`$. It follows that
$$\left|\frac{\mathrm{cos}\left(z\right)1}{z^2}\right|\mathrm{sinh}\left(\left|y\right|\right)+\mathrm{sinh}\left(\left|y\right|\right)+1+\mathrm{cosh}\left(y\right)4\mathrm{cosh}\left(y\right)$$
(3.68)
since $`1,\mathrm{sinh}\left(\left|y\right|\right)\mathrm{cosh}\left(y\right)`$.
ii)
Proceeding similarly as in i) we have
$`\left|{\displaystyle \frac{\mathrm{sin}\left(z\right)z}{z^3}}\right|=\left|{\displaystyle \frac{1}{z^3}}\left(\left[\mathrm{sin}\left(x\right)\mathrm{cosh}\left(y\right)x\right]+i\left[\mathrm{cos}\left(x\right)\mathrm{sinh}\left(y\right)y\right]\right)\right|`$ (3.69)
$``$ $`{\displaystyle \frac{1}{\left|z^3\right|}}\left|\left[xy^2{\displaystyle \frac{\mathrm{sin}\left(x\right)}{x}}{\displaystyle \frac{\mathrm{cosh}\left(y\right)1}{y^2}}\right]+\left[x^3{\displaystyle \frac{\mathrm{sin}\left(x\right)x}{x^3}}\right]+\left[x^2y{\displaystyle \frac{\mathrm{cos}\left(x\right)1}{x^2}}{\displaystyle \frac{\mathrm{sinh}\left(y\right)}{y}}\right]+\left[y^3{\displaystyle \frac{\mathrm{sinh}\left(y\right)y}{y^3}}\right]\right|`$
$``$ $`\left|{\displaystyle \frac{\mathrm{sin}\left(x\right)}{x}}\right|\left|{\displaystyle \frac{\mathrm{cosh}\left(y\right)1}{y^2}}\right|+\left|{\displaystyle \frac{\mathrm{sin}\left(x\right)x}{x^3}}\right|+\left|{\displaystyle \frac{\mathrm{cos}\left(x\right)1}{x^2}}\right|\left|{\displaystyle \frac{\mathrm{sinh}\left(y\right)}{y}}\right|+\left|{\displaystyle \frac{\mathrm{sinh}\left(y\right)y}{y^3}}\right|`$
$``$ $`\mathrm{cosh}\left(y\right)+\left|{\displaystyle \frac{\mathrm{sin}\left(x\right)x}{x^3}}\right|+\mathrm{cosh}\left(y\right)+\mathrm{cosh}\left(y\right)`$
where we have made use of properties already demonstrated in i) and the Taylor series of $`\mathrm{sinh}\left(y\right)`$. We claim that $`\left|\frac{\mathrm{sin}\left(x\right)x}{x^3}\right|=\left|\frac{\mathrm{sin}\left(\left|x\right|\right)\left|x\right|}{\left|x\right|^3}\right|1`$ and it is sufficient to prove this for $`x0`$. That statement is for $`x0`$ equivalent to $`f\left(x\right)=x^3+\mathrm{sin}\left(x\right)x0`$ because $`\mathrm{sin}\left(x\right)x`$ for $`x0`$. We have $`g\left(x\right)=f^{}\left(x\right)=3x^2+\mathrm{cos}\left(x\right)1,h\left(x\right)=g^{}\left(x\right)=6xsin\left(x\right),h^{}\left(x\right)=6\mathrm{cos}\left(x\right)>0`$. Since $`f\left(0\right)=g\left(0\right)=h\left(0\right)=0`$ we see that $`h`$ is an increasing function whence $`h\left(x\right)0`$ from which follows that $`g`$ is an increasing function whence $`g\left(x\right)0`$ from which follows that $`f`$ is an increasing function whence $`f\left(x\right)0`$ as claimed. It follows that
$$\left|\frac{\frac{\mathrm{sin}\left(z\right)}{z}1}{z^2}\right|1+3\mathrm{cosh}\left(y\right)4\mathrm{cosh}\left(y\right)$$
(3.70)
as claimed.
$`\mathrm{}`$
Collecting all the estimates we can now write the estimate for the modulus of (3.62) in the desired form, defining $`y^{}=4\pi n\varphi /t`$,
$$\left|J\right|\frac{4\pi n}{t}\left[2\mathrm{cosh}\left(y^{}\right)\left|I\right|+8\left(\frac{4\pi n}{t}\right)^2\mathrm{cosh}\left(y^{}\right)\right]\frac{32\pi n}{t}\left[k_c^{}+\left(\frac{4\pi n}{t}\right)^2\right]\mathrm{cosh}\left(y^{}\right)$$
(3.71)
which now enables us to estimate the various series $`I_1,I_2,I_3`$. We will do this one by one.
$`I_1`$ :
The elementary estimate $`\left|\mathrm{cos}\left(z\right)\right|\left|\mathrm{cosh}\left(y\right)\right|+\left|\mathrm{sinh}\left(y\right)\right|=e^{\left|y\right|}`$ applied to $`I_1`$ reveals
$`\left|I_1\right|`$ $``$ $`1+2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{4\pi ^2n^2}{t}}e^{\frac{4\pi ^2n\left(1c\right)}{t}}`$ (3.72)
$``$ $`1+2e^{\frac{4\pi ^2c}{t}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{4\pi ^2\left(n^2n\right)}{t}}`$
$``$ $`1+2e^{\frac{4\pi ^2c}{t}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}e^{\frac{4\pi ^2n^2}{t}}=:1+k_t`$
where in the last step we have made use of the inequality $`\left(n1\right)^2n^2n`$ valid for $`n1`$. The constant $`k_t`$ is independent of $`g,g^{}`$ and vanishes exponentially fast with $`t0`$ for any $`c>0`$.
$`I_2`$ :
Using again $`\left|\mathrm{cos}\left(z\right)\right|e^{\left|y\right|},\left|\mathrm{sin}\left(z\right)/z\right|2\mathrm{cosh}\left(y\right)2e^{\left|y\right|}`$ we easily find with $`\left|\varphi \right|\left(1c\right)\pi `$
$`\left|I_2\right|`$ $``$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{4\pi ^2n^2}{t}}e^{\frac{4\pi ^2n\left(1c\right)}{t}}\left(1+{\displaystyle \frac{32\pi ^2n^2}{t}}\right)`$ (3.73)
$``$ $`e^{\frac{4\pi ^2c}{t}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{4\pi ^2\left(n^2n\right)}{t}}\left(1+{\displaystyle \frac{32\pi ^2n^2}{t}}\right)`$
$``$ $`e^{\frac{4\pi ^2c}{t}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}e^{\frac{4\pi ^2n^2}{t}}(1+{\displaystyle \frac{32\pi ^2\left(n+1\right)^2}{t}})=:k_t^{}`$
where again $`k_t^{}`$ is a constant, approaching zero exponentially fast with $`t0`$ for any $`c>0`$.
$`I_3`$ :
Invoking our estimate (3.71) for the quantity $`J`$ (3.62) that appears in $`I_3`$ we obtain
$`\left|I_3\right|`$ $``$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{4\pi ^2n^2}{t}}e^{\frac{4\pi ^2n\left(1c\right)}{t}}{\displaystyle \frac{32\pi n^2}{t}}\left[k_c^{}+\left({\displaystyle \frac{4\pi n}{t}}\right)^2\right]`$ (3.74)
$``$ $`32\pi e^{\frac{4\pi ^2c}{t}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{4\pi ^2\left(n^2n\right)}{t}}{\displaystyle \frac{n^2}{t}}\left[k_c^{}+\left({\displaystyle \frac{4\pi n}{t}}\right)^2\right]`$
$``$ $`32\pi e^{\frac{4\pi ^2c}{t}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}e^{\frac{4\pi ^2n^2}{t}}{\displaystyle \frac{\left(n+1\right)^2}{t}}[k_c^{}+\left({\displaystyle \frac{4\pi \left(n+1\right)}{t}}\right)^2]=:\stackrel{~}{k}_t`$
where again $`\stackrel{~}{k}_t`$ is a constant independent of both $`g,g^{}`$ exponentially vanishing with $`t0`$, for any $`c>0`$.
Let us now define
$`\mathrm{\Delta }<\widehat{p}_j>_{gg^{}}^t`$ $`:=`$ $`<\widehat{p}_j>_{gg^{}}^t{\displaystyle \frac{\left[\frac{i}{2}\frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{\mathrm{sinh}\left(z_0\right)}z_0\right]e^{\frac{z_0^2p^2/2\left(p^{}\right)^2/2}{t}}\frac{z_0}{\mathrm{sinh}\left(z_0\right)}}{\sqrt{D^t\left(p\right)D^t\left(p^{}\right)}}}`$
$`=`$ $`{\displaystyle \frac{\left[\frac{i}{2}\frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{\mathrm{sinh}\left(z_0\right)}z_0\right]e^{\frac{z_0^2p^2/2\left(p^{}\right)^2/2}{t}}\frac{z_0}{\mathrm{sinh}\left(z_0\right)}\left\{\frac{t}{2}\frac{\frac{\mathrm{sinh}\left(z_0\right)}{z_0}\mathrm{cosh}\left(z_0\right)}{z_0\mathrm{sinh}\left(z_0\right)}I_1+2I_2+2\pi tI_3\right\}}{\sqrt{D^t\left(p\right)D^t\left(p^{}\right)}}}`$
Recall now the relation between the various objects $`s,\varphi ,\stackrel{~}{p},\stackrel{~}{\theta },\stackrel{~}{\alpha }`$ from section 4.2 of . Basically, one writes $`g=Hh,g^{}=H^{}h^{}`$ in polar decomposed form and defines the polar decomposition $`HH^{}=\stackrel{~}{H}u`$ as well as $`\stackrel{~}{h}=h^1h^{}u^1`$. Then $`\stackrel{~}{H}=\mathrm{exp}\left(i\stackrel{~}{p}_j\tau _j/2\right),\stackrel{~}{h}=\mathrm{exp}\left(\stackrel{~}{\theta }_j\tau _j\right)`$ and $`\mathrm{cosh}\left(s\right)\mathrm{cos}\left(\varphi \right)=\mathrm{cosh}\left(\stackrel{~}{p}/2\right)\mathrm{cos}\left(\stackrel{~}{\theta }\right)`$ and $`\mathrm{sinh}\left(s\right)\mathrm{sin}\left(\varphi \right)=\mathrm{sinh}\left(\stackrel{~}{p}/2\right)\mathrm{sin}\left(\stackrel{~}{\theta }\right)\mathrm{cos}\left(\stackrel{~}{\alpha }\right)`$ with $`\mathrm{cos}\left(\stackrel{~}{\alpha }\right)=\stackrel{~}{p}_j\stackrel{~}{\theta }_j/\left(\stackrel{~}{p}\stackrel{~}{\theta }\right)`$.
Writing $`g^{}\overline{g}^T=e^{i\tau _jz_0^j}`$ we have
$$\left[\frac{i}{2}\frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{\mathrm{sinh}\left(z_0\right)}z_0\right]=z_0^j$$
(3.76)
and are interested in the relation between $`z_0^j`$ and $`z_0=s+i\varphi `$. By definition we have $`\mathrm{cosh}\left(z_0\right)=\text{tr}\left(g^{}\overline{g}^T\right)/2`$ which reveals that $`z_0^2=_j\left(z_0^j\right)^2`$, however, it is not true that $`\left|z_0\right|^2_j\left|z_0^j\right|^2`$. In order to estimate the integral over $`\mathrm{\Delta }<\widehat{p}_j>_{g,g^{}}`$ we thus have to prove one more relation.
###### Lemma 3.5
For $`z_0=s+i\varphi ,\mathrm{\hspace{0.33em}0}\varphi (1c)\pi `$ and $`\stackrel{~}{g}:=g^{}\overline{g}^T=e^{i\tau _jz_0^j}`$ with $`z_0^2=z_0^jz_0^j`$ we have
$$\left|z_0^j\right|^2\left[k_c\frac{s}{\mathrm{sinh}\left(s\right)}\sqrt{\mathrm{cosh}\left(\stackrel{~}{p}\right)}\right]^2\left[2k_c\frac{s}{\mathrm{sinh}\left(s\right)}\mathrm{cosh}\left(\frac{p+p^{}}{2}\right)\right]^2$$
(3.77)
for any $`j=1,2,3`$.
Proof of Lemma 3.5 :
Using the $`SL(2,\text{ }\mathrm{C})`$ “Fierz identity” $`\text{tr}\left(M\tau _j\right)\tau _j=\text{tr}\left(M\right)2M`$ valid for any $`2\times 2`$ matrix $`M`$ we find from (3.76) that
$`{\displaystyle \underset{j}{}}\left|z_0^j\right|^2`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left|{\displaystyle \frac{z_0}{\mathrm{sinh}\left(z_0\right)}}\right|^2\text{tr}\left(\tau _j\stackrel{~}{g}\right)\text{tr}\left(\tau _j\overline{\stackrel{~}{g}}^T\right)`$ (3.78)
$`=`$ $`{\displaystyle \frac{1}{4}}\left|{\displaystyle \frac{z_0}{\mathrm{sinh}\left(z_0\right)}}\right|^2\left(\left|\text{tr}\left(\stackrel{~}{g}\right)\right|^22\text{tr}\left(\stackrel{~}{g}\overline{\stackrel{~}{g}}^T\right)\right)`$
$``$ $`{\displaystyle \frac{1}{2}}\left|{\displaystyle \frac{z_0}{\mathrm{sinh}\left(z_0\right)}}\right|^2\text{tr}\left(\stackrel{~}{H}^2\right)=\left|{\displaystyle \frac{z_0}{\mathrm{sinh}\left(z_0\right)}}\right|^2\mathrm{cosh}\left(\stackrel{~}{p}\right)`$
Now recall from that
$$\mathrm{cosh}\left(\stackrel{~}{p}\right)=\left(1+r\right)\mathrm{cosh}^2\left(\frac{p+p^{}}{2}\right)+\left(1r\right)\mathrm{cosh}^2\left(\frac{pp^{}}{2}\right)1$$
(3.79)
where $`r=p_jp_j^{}/\left(pp^{}\right)[1,1]`$. Combining (3.79) and (3.78) yields (notice that $`p+p^{}>\left|pp^{}\right|,\left|r\right|1`$)
$$\underset{j}{}\left|z_0^j\right|^24\left|\frac{z_0}{\mathrm{sinh}\left(z_0\right)}\right|^2\mathrm{cosh}^2\left(\frac{p+p^{}}{2}\right)\left[2k_c\frac{s}{\mathrm{sinh}\left(s\right)}\mathrm{cosh}\left(\frac{p+p^{}}{2}\right)\right]^2$$
(3.80)
$`\mathrm{}`$
Next, by estimates established in we have
$$\mathrm{\Delta }^2(\stackrel{}{p},\stackrel{}{p}^{}):=p^2+\left(p^{}\right)^2\frac{\stackrel{~}{p}^2}{2}0\text{ and }\delta ^2(g,g^{}):=\stackrel{~}{p}^2/4s^2+\varphi ^2\stackrel{~}{\theta }^20$$
(3.81)
where $`\mathrm{\Delta }=0`$ if and only if $`\stackrel{}{p}=\stackrel{}{p}^{}`$ and $`\delta =0`$ if either a) $`\stackrel{~}{\alpha }=0,\pi `$ and $`\stackrel{~}{p},\stackrel{~}{\theta }`$ are arbitrary or b) $`\stackrel{~}{\alpha }`$ is arbitrary and one or both of $`\stackrel{~}{p}=0;\stackrel{~}{\theta }=0,\pi `$ hold. In both of the cases a),b) we have $`\left|s\right|=\stackrel{~}{p}/2,\varphi =\stackrel{~}{\theta }`$. It follows that
$$\mathrm{}\left(z_0^2p^2/2\left(p^{}\right)^2/2\right)=\left\{\left[\frac{p^2}{2}+\frac{\left(p^{}\right)^2}{2}\frac{\stackrel{~}{p}^2}{4}\right]+\stackrel{~}{\theta }^2+\left[s^2+\varphi ^2\stackrel{~}{\theta }^2+\frac{\stackrel{~}{p}^2}{4}\right]\right\}=\left[\mathrm{\Delta }^2/2+\delta ^2+\stackrel{~}{\theta }^2\right]$$
(3.82)
Combining the estimates (3.35) for $`D^t\left(p\right)`$, (3.60) for $`I`$, (3.72), (3.73) and (3.74) respectively for $`I_1,I_2,I_3`$ respectively, (3.77) for $`\left|z_0^j\right|`$, (3.41) for $`f_c\left(s\right)`$ and (3.82) for the Gaussian prefactor of $`N^t(g,g^{})`$ we conclude
$`\left|\mathrm{\Delta }<\widehat{p}_j>_{gg^{}}^t\right|\left|z_0^j\left|e^{\frac{\mathrm{}\left(z_0^2p^2/2\left(p^{}\right)^2/2\right)}{t}}\right|{\displaystyle \frac{z_0}{\mathrm{sinh}\left(z_0\right)}}\right|{\displaystyle \frac{\left\{\frac{t}{2}\left|I\right|\left|I_1\right|+2\left|I_2\right|+2\pi t\left|I_3\right|\right\}}{\sqrt{D^t\left(p\right)D^t\left(p^{}\right)}}}`$
$``$ $`\left[2k_c{\displaystyle \frac{s}{\mathrm{sinh}\left(s\right)}}\mathrm{cosh}\left({\displaystyle \frac{p+p^{}}{2}}\right)\right]\left[e^{\frac{\mathrm{\Delta }^2/2+\delta ^2+\stackrel{~}{\theta }^2}{t}}\right]\left[{\displaystyle \frac{k_c\frac{s}{\mathrm{sinh}\left(s\right)}}{\left(1K_t\right)\sqrt{\frac{p}{\mathrm{sinh}\left(p\right)}\frac{p^{}}{\mathrm{sinh}\left(p^{}\right)}}}}\right]\times `$
$`\times `$ $`\left\{\left[tk_c^{}\left(1+k_t\right)\right]+\left[2k_t^{}\right]+\left[2\pi t\stackrel{~}{k}_t\right]\right\}`$
$`=`$ $`2k_c^2\left[e^{\frac{\mathrm{\Delta }^2/2+\stackrel{~}{\theta }^2}{t}}\right]\left[e^{\delta ^2/t}{\displaystyle \frac{s^2}{\sqrt{pp^{}}}}{\displaystyle \frac{\mathrm{cosh}\left(\frac{p+p^{}}{2}\right)\sqrt{\mathrm{sinh}\left(p\right)\mathrm{sinh}\left(p^{}\right)}}{\mathrm{sinh}^2\left(s\right)}}\right]\left[{\displaystyle \frac{tk_c^{}\left(1+k_t\right)+2k_t^{}+2\pi t\stackrel{~}{k}_t}{1K_t}}\right]`$
Let us discuss this result. The last line of (3.1.2) consists of three factors corresponding to the three square brackets. The first bracket contains a Gaussian with peak of width of order $`\sqrt{t}`$ at $`g=g^{}`$. The third bracket is of the form $`t\left(1+K_t^{}\left(c\right)\right)`$ where $`K_t^{}\left(c\right)`$ is exponentially vanishing with $`t`$ for any $`0<c<1`$. These two brackets are expected from the harmonic oscillator. The remaining second bracket
$$B_2:=e^{\delta ^2/t}\frac{s^2}{\sqrt{pp^{}}}\frac{\mathrm{cosh}\left(\frac{p+p^{}}{2}\right)\sqrt{\mathrm{sinh}\left(p\right)\mathrm{sinh}\left(p^{}\right)}}{\mathrm{sinh}^2\left(s\right)}$$
(3.84)
is unexpected, it is not manifestly bounded from above and thus the integral over $`g^{}`$ is not obviously converging. The subsequent paragraph will be devoted to the behaviour of that term.
Since (3.84) is manifestly regular at $`p^{}=0`$, convergence problems can arise only for $`p^{}\mathrm{}`$. Formula (3.79) implies that then also $`\stackrel{~}{p}\mathrm{}`$, no matter which values $`p,r`$ take, actually $`\stackrel{~}{p}p^{}`$ as $`p^{}\mathrm{}`$ for fixed $`p,r`$. By estimate (3.81) we then see that either a) $`\delta \mathrm{}`$ or b) $`\delta `$ stays bounded from above. In the first case $`s`$ must stay bounded or grows slowlier than $`\stackrel{~}{p}/2`$ while in the latter case $`s\stackrel{~}{p}/2p/2^{}`$. Consider first case a). Then for large $`p^{}`$ the Gaussian $`e^{\delta ^2/t}`$ in (3.84) certainly wins over the remaining factor and $`B_2`$ decays exponentially fast. Next consider case b). In that case $`B_2`$ grows as $`\sqrt{p^{}}^3`$ at large $`p^{}`$. Thus, altogether we have shown that $`B_2`$ grows no worse than polynomially as $`p^{}\mathrm{}`$. Notice that $`s`$ stays bounded if and only if $`\stackrel{~}{\theta }=\stackrel{~}{\alpha }=\pi /2`$ which defines a set of $`\mathrm{\Omega }`$ measure zero. In that case in fact $`s=0,\varphi =\pi /2`$, therefore $`z_0=i\pi /2`$ whence $`\stackrel{~}{g}=\tau _j\stackrel{~}{\theta }_j\mathrm{cosh}\left(\stackrel{~}{p}/2\right)/\stackrel{~}{\theta }`$ and $`_j\left|z_0^j/z_0\right|^2=\mathrm{cosh}^2\left(\stackrel{~}{p}/2\right)`$.
We conclude that in the range $`0\varphi \left(1c\right)\pi `$ $`\mathrm{\Delta }<\widehat{p}_j>_{gg^{}}^t`$ can be bounded by a function of $`g,g^{}`$ which is of the form of a Gaussian with peak at $`g=g^{}`$ times a function of $`g,g^{}`$ bounded by a polynomial in $`\stackrel{}{p},\stackrel{}{p}^{}`$ times $`t`$ times a constant that approaches unity exponentially fast. Thus the integral of that function with respect to $`g^{}`$ exists (since $`\mathrm{\Delta }^2`$ approaches $`p^2/2`$ at large $`p^{}`$) and will be of the form of a function of $`g`$ bounded by a polynomial in $`p`$ times $`t`$ times a constant that approaches unity exponentially fast. It should be noted that the six Gaussians in (3.1.2) of width $`\sqrt{t}`$ each cancel the $`1/t^3`$ in the measure $`d\mathrm{\Omega }/t^3`$ (recall from (3.18), (3.19) that $`\psi _g^t^2\nu _t\left(g\right)`$ approaches $`2/\left(\pi t^3\right)`$ exponentially fast).
It follows that as far as the leading order (in $`t`$) behaviour of the expectation value of any monomial is concerned, we can drop the term $`\mathrm{\Delta }<\widehat{p}_j>_{g,g^{}}^t`$ from $`<\widehat{p}_j>_{g,g^{}}^t`$, at least in the range $`0\varphi \left(1c\right)\pi `$.
Case B)
Let us now discuss the range $`\left(1c\right)\pi \varphi \pi `$. We will not be as explicit as in case A), the steps to be performed are essentially identical to the case A) and can be found in more detail in the analogous discussion of .
The essential point is now to write everything in terms of
$$z_0^{}:=z_0i\pi =si(\pi \varphi )=:si\varphi ^{}\text{ where }0\varphi ^{}c\pi $$
(3.85)
Starting with the expression (3.1.2) we have
$`<\widehat{p}_j>_{gg^{}}^t=\left[{\displaystyle \frac{i}{2}}{\displaystyle \frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{\mathrm{sinh}\left(z_0\right)}}\right]\times `$
$`\times `$ $`{\displaystyle \frac{\left\{\frac{t}{2}_ne^{\frac{\left(z_0^{}i\pi \left(2n1\right)\right)^2p^2/2\left(p^{}\right)^2/2}{t}}\left[2\frac{\left(z_0^{}i\pi \left(2n1\right)\right)^2}{t\mathrm{sinh}\left(z_0\right)}\left(z_0^{}i\pi \left(2n1\right)\right)\frac{\mathrm{cosh}\left(z_0\right)}{\mathrm{sinh}^2\left(z_0\right)}+\frac{1}{\mathrm{sinh}\left(z_0\right)}\right]\right\}}{\left[_n\frac{p2\pi in}{\mathrm{sinh}\left(p\right)}e^{\frac{4\pi ^2n^2}{t}}e^{i\frac{4\pi np}{t}}\right]^{1/2}\left[_n\frac{p^{}2\pi in}{\mathrm{sinh}\left(p^{}\right)}e^{\frac{4\pi ^2n^2}{t}}e^{i\frac{4\pi np^{}}{t}}\right]^{1/2}}}`$
$`=`$ $`\left[{\displaystyle \frac{i}{2}}{\displaystyle \frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{\mathrm{sinh}\left(z_0\right)}}\right]\times `$
$`\times `$ $`e^{\frac{\left(z_0^{}\right)^2p^2/2\left(p^{}\right)^2/2}{t}}{\displaystyle \frac{\left\{\frac{t}{2}_{n=\text{odd}}e^{\frac{\pi ^2n^2}{t}}e^{\frac{2i\pi z_0^{}n}{t}}\left[2\frac{\left(z_0^{}i\pi n\right)^2}{t\mathrm{sinh}\left(z_0^{}\right)}\left(z_0^{}i\pi n\right)\frac{\mathrm{cosh}\left(z_0\right)}{\mathrm{sinh}^2\left(z_0\right)}+\frac{1}{\mathrm{sinh}\left(z_0\right)}\right]\right\}}{\sqrt{\frac{p}{\mathrm{sinh}\left(p\right)}\frac{p^{}}{\mathrm{sinh}\left(p^{}\right)}\left(1K_t\left(p\right)\right)\left(1K_t\left(p^{}\right)\right)}}}`$
where $`\left|K_t\left(p\right)\right|,\left|K_t\left(p^{}\right)\right|K_t`$ are the functions implicit in the estimate (3.35) bounded from above by exponentially vanishing constants. We now pull out a factor of $`\left(z_0^{}\right)^2/\mathrm{sinh}\left(z_0\right)`$ out of the series in (3.1.2), collect terms as to sum over positive, odd integers $`n`$ only and observe that $`\mathrm{sinh}\left(z_0\right)=\mathrm{sinh}\left(z_0^{}\right)`$ and $`\mathrm{cosh}\left(z_0\right)=\mathrm{cosh}\left(z_0^{}\right)`$. Then (3.1.2) becomes
$`<\widehat{p}_j>_{gg^{}}^t={\displaystyle \frac{\left[\frac{iz_0^{}}{2}\frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{\mathrm{sinh}\left(z_0^{}\right)}\right]\left[\frac{z_0^{}}{\mathrm{sinh}\left(z_0^{}\right)}\right]e^{\frac{\left(z_0^{}\right)^2p^2/2\left(p^{}\right)^2/2}{t}}}{\sqrt{\frac{p}{\mathrm{sinh}\left(p\right)}\frac{p^{}}{\mathrm{sinh}\left(p^{}\right)}\left(1K_t\left(p\right)\right)\left(1K_t\left(p^{}\right)\right)}}}\times `$
$`\times `$ $`{\displaystyle \underset{n=1,\text{odd}}{\overset{\mathrm{}}{}}}e^{\frac{\pi ^2n^2}{t}}\{2[\mathrm{cos}\left({\displaystyle \frac{2\pi nz_0^{}}{t}}\right){\displaystyle \frac{\left(2\pi n\right)^2}{t}}{\displaystyle \frac{\mathrm{sin}\left(\frac{2\pi nz_0^{}}{t}\right)}{\frac{2\pi nz_0^{}}{t}}}]+t[{\displaystyle \frac{1}{\left(z_0^{}\right)^2}}{\displaystyle \frac{\mathrm{cosh}\left(z_0^{}\right)}{z_0^{}\mathrm{sinh}\left(z_0^{}\right)}}]\mathrm{cos}\left({\displaystyle \frac{2\pi nz_0^{}}{t}}\right)`$
$`+\pi nt[{\displaystyle \frac{\mathrm{cosh}\left(z_0^{}\right)}{\left(z_0^{}\right)^2\mathrm{sinh}\left(z_0^{}\right)}}\mathrm{sin}\left({\displaystyle \frac{2\pi nz_0^{}}{t}}\right){\displaystyle \frac{2\pi n}{\left(z_0^{}\right)^2t}}\mathrm{cos}\left({\displaystyle \frac{2\pi nz_0^{}}{t}}\right)]\}`$
which can now be estimated essentially as in case A). The most important differences are the following : First, lemma 3.1 has now to be replaced by $`\left|z_0^{}/\mathrm{sinh}\left(z_0^{}\right)\right|k_{1c}s/\mathrm{sinh}\left(s\right)`$ which can be seen by following the proof given there step by step. Secondly, by the methods given in the Gaussian is now estimated from above by
$$e^{\frac{\mathrm{\Delta }^2+2\delta ^2+\stackrel{~}{\theta }^2}{2t}}$$
(3.88)
where necessarily $`\stackrel{~}{\theta }\pi /2`$ and is therefore exponentially small for all values of $`g^{}`$ in the range $`\left(1c\right)\pi \varphi \pi `$ provided that we choose $`c<1/2`$ as we do. Finally, choosing as in $`c=1/32`$ (see specifically formula (4.44) there) and using the same estimates of case A) we can display (3.1.2) as a function of $`g,g^{}`$ bounded by a polynomial in $`p,p^{}`$ times a Gaussian in $`\mathrm{\Delta }^2`$ multiplied by an overall constant which decays exponentially fast to zero as $`t0`$. Thus, (3.1.2) is an $`\mathrm{\Omega }/t^3`$ integrable function with respect to $`g^{}`$ and the result of the integration is a function of $`g`$ bounded by a polynomial in $`p`$ times a constant which decays exponentially fast as $`t0`$.
We conclude that the range of integration $`\left(1c\right)\pi \varphi \pi `$ is irrelevant for the expectation value and all its corrections in powers of $`t`$.
We can now finish the estimate of the matrix element. Our discussion has demonstrated that to leading order in $`t`$ we can replace $`<\widehat{p}_j>_{gg^{}}^t`$ by
$$\frac{\left[\frac{i}{2}\frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{\mathrm{sinh}\left(z_0\right)}z_0\right]e^{\frac{z_0^2p^2/2\left(p^{}\right)^2/2}{t}}\frac{z_0}{\mathrm{sinh}\left(z_0\right)}\mathrm{\Theta }\left(\left(1c\right)\pi \varphi \right)}{\sqrt{D^t\left(p\right)D^t\left(p^{}\right)}}$$
(3.89)
where $`\mathrm{\Theta }`$ is the step function. The expression (3.89) is Gaussian peaked at $`g=g^{}`$ with decay width of order $`\sqrt{t}`$ at which it equals $`p_j`$. In order to perform the integral over $`g^{}`$ we will therefore expand the square bracket in the numerator of (3.89) as
$$p_j+\left[\frac{i}{2}\frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{\mathrm{sinh}\left(z_0\right)}z_0p_j\right]$$
(3.90)
and the integral over the $`g^{}`$ independent term with respect to $`\mathrm{\Omega }\left(g^{}\right)/t^3`$ converges exponentially fast to $`p_j`$ since the absolute value squared of (3.89) modulo the square bracket equals precisely the overlap function of modulo a multiplicative term whose absolute value can be estimated from above by a constant that approaches unity exponentially fast. The integral over the remaining term can be expanded as a function of $`\stackrel{}{p}\stackrel{}{p}^{},h\left(h^{}\right)^1`$ and vanishes at least linearly in $`t`$ by standard properties of Gaussian integrals.
Collecting all the results we have arrived at the first main theorem of this paper.
###### Theorem 3.2
The matrix elements of the momentum operators with respect to coherent states can be estimated by
$$\left|\frac{<\psi _g^t,\widehat{p}_j\psi _g^{}^t>}{\psi _g^t\psi _g^{}^t}p_j\left(g\right)\frac{<\psi _g^t,\psi _g^{}^t>}{\psi _g^t\psi _g^{}^t}\right|tf(\stackrel{}{p},\stackrel{}{p}^{})\frac{|<\psi _g^t,\psi _g^{}^t>|}{\psi _g^t\psi _g^{}^t}$$
(3.91)
where $`f`$ is a polynomial of $`p,p^{}`$.
As a corollary to theorem (3.2) we obtain that the expectation value $`<\widehat{p}_j>_{gg}^t`$ equals $`p_j\left(g\right)`$ up to bounded corrections in $`p_j\left(g\right)`$ and that are proportional to $`t`$. We will actually calculate the exact correction in a later section by a different method.
#### 3.1.3 Matrix Elements for the Holonomy Operator
The computation of the matrix element of the holonomy operator
$$<\widehat{h}_{AB}>_{gg^{}}^t:=\frac{<\psi _g^t,\widehat{h}_{AB}\psi _g^{}^t>}{\psi _g^t\psi _g^{}^t}$$
(3.92)
turns out to be rather messy. Let us determine first the following matrix element, using the Peter&Weyl theorem and the $`SL(2,\text{ }\mathrm{C})`$ identity $`\chi _j\left(g\right)=\chi _j\left(g^1\right)`$
$`<\psi _g^t,\widehat{h}_{A_0B_0}\psi _g^{}^t>`$
$`=`$ $`{\displaystyle \underset{j,j^{}}{}}d_jd_j^{}e^{t\left[\lambda _j+\lambda _j^{}\right]/2}\overline{\pi _j\left(g\right)_{A_1..A_{2j},B_1..B_{2j}}}\pi _j^{}\left(g^{}\right)_{A_1^{}..A_{2j^{}}^{},B_1^{}..B_{2j^{}}^{}}\times `$
$`\times `$ $`<\pi _j^{}(.)_{A_1^{}..A_{2j}^{},B_1^{}..B_{2j}^{}},\widehat{h}_{A_0B_0}\pi _j(.)_{A_1..A_{2j},B_1..B_{2j}}>`$
$`=`$ $`{\displaystyle \underset{j,j^{}}{}}d_jd_j^{}e^{t\left[\lambda _j+\lambda _j^{}\right]/2}\overline{\pi _j\left(g\right)_{A_1..A_{2j},B_1..B_{2j}}}\pi _j^{}\left(g^{}\right)_{A_1^{}..A_{2j^{}}^{},B_1^{}..B_{2j^{}}^{}}\times `$
$`\times `$ $`[{\displaystyle \frac{\delta _{j^{},j+\frac{1}{2}}}{d_j^{}}}\pi _j^{}\left(1\right)_{A_1^{}..A_{2j^{}}^{},A_0A_1..A_{2j}}\pi _j^{}\left(1\right)_{B_1^{}..B_{2j^{}}^{},B_0B_1..B_{2j}}+{\displaystyle \frac{\delta _{j^{},j\frac{1}{2}}d_{j\frac{1}{2}}}{d_jd_j^{}}}\times `$
$`\times `$ $`\pi _j^{}\left(1\right)_{A_1^{}..A_{2j^{}}^{},(A_2..A_{2j}}ϵ_{A_1)A_0}\pi _j^{}\left(1\right)_{B_1^{}..B_{2j^{}}^{},(B_2..B_{2j}}ϵ_{B_1)B_0}]`$
$`=`$ $`{\displaystyle \underset{j}{}}d_je^{t\lambda _j/2}\pi _j\left(\overline{g}^T\right)_{B_1..B_{2j},A_1..A_{2j}}\times `$
$`\times `$ $`\left[e^{t\lambda _{j+\frac{1}{2}}/2}\pi _{j+\frac{1}{2}}\left(g^{}\right)_{A_0A_1..A_{2j},B_0B_1..B_{2j}}{\displaystyle \frac{d_{j\frac{1}{2}}}{d_j}}e^{t\lambda _{j\frac{1}{2}}/2}ϵ_{A_0(A_1}\pi _{j\frac{1}{2}}\left(g^{}\right)_{A_2..A_{2j}),(B_2^{}..B_{2j}^{}}ϵ_{B_1)B_0}\right]`$
In the second step we have recalled the following (Clebsch-Gordan) identity, valid for arbitrary $`gG^{\text{ }\mathrm{C}}=SL(2,\text{ }\mathrm{C})`$ and proved in
$`g_{A_0B_0}\pi _j\left(g\right)_{A_1..A_{2j},B_1..B_{2j}}`$ (3.94)
$`=`$ $`\pi _{j+\frac{1}{2}}\left(g\right)_{A_0..A_{2j},B_0..B_{2j}}{\displaystyle \frac{d_{j\frac{1}{2}}}{d_j}}ϵ_{A_0(A_1}\pi _{j\frac{1}{2}}\left(g\right)_{A_2..A_{2j}),(B_2..B_{2j}}ϵ_{B_1)B_0}`$
with $`A,B,..=\pm 1/2`$, round brackets around groups of indices denote total symmetrization taken as an idempotent operation, $`ϵ_{AB}`$ is the skew symmetric spinor of rank two, $`d_j=dim\left(\pi _j\right)=2j+1`$ and the relation with the usual matrix elements of the irreducible representation $`\pi _j`$ is given by $`\pi _j\left(g\right)_{A_1+..+A_{2j},B_1+..+B_{2j}}=\pi _j\left(g\right)_{A_1..A_{2j},B_1..B_{2j}}`$.
The huge amount of summation indices that appear in (3.1.3) and which we cannot nicely contract as irreducible representations of different dimension are multiplied with each other make (3.1.3) impossible to work with because then we cannot apply the Weyl character – and Poisson summation formula, our main tools in all the estimates. Fortunately, we have the following trick at our disposal (it obviously extends to groups of higher rank) :
Let $`\mathrm{\Delta }_h`$ be the Laplacian on $`G=SU\left(2\right)`$ acting on $`hG`$. Since $`\pi _j\left(hg\right)_{mn}=\pi _j\left(h\right)_{mm^{}}\pi _j\left(g\right)_{m^{}n}`$ is an eigenstate of $`\mathrm{\Delta }_h`$ with eigenvalue $`j\left(j+1\right)`$ we obtain the following formulas that isolate the irreducible pieces on the right hand side of (3.94)
$`\left\{\left(j+{\displaystyle \frac{1}{4}}\right)\left(hg\right)_{A_0B_0}+[\mathrm{\Delta }_h,\left(hg\right)_{A_0B_0}]\right\}\pi _j\left(hg\right)_{A_1..B_{2j}}`$
$`=`$ $`d_j\pi _{j+\frac{1}{2}}\left(hg\right)_{A_0..A_{2j},B_0..B_{2j}}`$
$`\left\{\left(j+{\displaystyle \frac{3}{4}}\right)\left(hg\right)_{A_0B_0}[\mathrm{\Delta }_h,\left(hg\right)_{A_0B_0}]\right\}\pi _j\left(hg\right)_{A_1..B_{2j}}`$
$`=`$ $`d_{j\frac{1}{2}}ϵ_{A_0(A_1}\pi _{j\frac{1}{2}}\left(hg\right)_{A_2..A_{2j}),(B_2..B_{2j}}ϵ_{B_1)B_0}`$ (3.96)
Taking the limit $`h1`$ in (3.1.3), (3.1.3) we can cast (3.1.3) into the following simpler form which allows us to contract indices
$`<\psi _g^t,\widehat{h}_{A_0B_0}\psi _g^{}^t>`$ (3.97)
$`=`$ $`{\displaystyle \underset{j}{}}e^{t\lambda _j/2}\pi _j\left(\overline{g}^T\right)_{B_1..B_{2j},A_1..A_{2j}}\times `$
$`\times `$ $`\{[e^{t\lambda _{j+\frac{1}{2}}/2}((j+{\displaystyle \frac{1}{4}})\left(hg^{}\right)_{A_0B_0}+[\mathrm{\Delta }_h,\left(hg^{}\right)_{A_0B_0}])`$
$`+e^{t\lambda _{j\frac{1}{2}}/2}((j+{\displaystyle \frac{3}{4}})\left(hg^{}\right)_{A_0B_0}[\mathrm{\Delta }_h,\left(hg^{}\right)_{A_0B_0}])]\pi _j\left(hg^{}\right)_{A_1..A_{2j},B_1..B_{2j}}\}_{|h=1}`$
$`=`$ $`{\displaystyle \underset{j}{}}e^{t\lambda _j/2}\times `$
$`\times `$ $`\{[e^{t\lambda _{j+\frac{1}{2}}/2}((j+{\displaystyle \frac{1}{4}})\left(hg^{}\right)_{A_0B_0}+[\mathrm{\Delta }_h,\left(hg^{}\right)_{A_0B_0}])`$
$`+e^{t\lambda _{j\frac{1}{2}}/2}((j+{\displaystyle \frac{3}{4}})\left(hg^{}\right)_{A_0B_0}[\mathrm{\Delta }_h,\left(hg^{}\right)_{A_0B_0}])]\chi _j\left(hg^{}\overline{g}^T\right)\}_{|h=1}`$
$`=`$ $`\left(g^{}\right)_{A_0B_0}{\displaystyle \underset{j}{}}e^{t\lambda _j/2}\left[e^{t\lambda _{j+\frac{1}{2}}/2}\left(j+{\displaystyle \frac{1}{4}}\right)+e^{t\lambda _{j\frac{1}{2}}/2}\left(j+{\displaystyle \frac{3}{4}}\right)\right]\chi _j\left(g^{}\overline{g}^T\right)`$
$`+`$ $`\left\{[\mathrm{\Delta }_h,\left(hg^{}\right)_{A_0B_0}]{\displaystyle \underset{j}{}}e^{t\lambda _j/2}\left[e^{t\lambda _{j+\frac{1}{2}}/2}e^{t\lambda _{j\frac{1}{2}}/2}\right]\chi _j\left(hg^{}\overline{g}^T\right)\right\}_{|h=1}`$
Formula (3.97) can be further simplified by making use of the commutator identity (use $`\mathrm{\Delta }_h=X_h^jX_h^j`$)
$$\left\{[\mathrm{\Delta }_h,\left(hg^{}\right)_{A_0B_0}]f\left(h\right)\right\}_{|h=1}=\frac{3}{4}g_{A_0B_0}^{}f\left(1\right)\frac{1}{2}\left(\tau _jg^{}\right)_{A_0B_0}\left(\frac{d}{ds}\right)_{s=0}f\left(e^{s\tau _j}\right)$$
(3.98)
valid for any differentiable function of $`h`$. Inserting this into (3.97) results in the final formula
$`<\psi _g^t,\widehat{h}_{AB}\psi _g^{}^t>`$ (3.99)
$`=`$ $`g_{AB}^{}{\displaystyle \underset{j}{}}e^{t\lambda _j/2}\left[\left(j+1\right)e^{t\lambda _{j+\frac{1}{2}}/2}+je^{t\lambda _{j\frac{1}{2}}/2}\right]\chi _j\left(g^{}\overline{g}^T\right)`$
$``$ $`{\displaystyle \frac{1}{2}}\left(\tau _jg^{}\right)_{AB}\left({\displaystyle \frac{d}{ds}}\right)_{s=0}{\displaystyle \underset{j}{}}e^{t\lambda _j/2}\left[e^{t\lambda _{j+\frac{1}{2}}/2}e^{t\lambda _{j\frac{1}{2}}/2}\right]\chi _j\left(e^{s\tau _j}g^{}\overline{g}^T\right)`$
to which we can now apply the Weyl character formula.
Let again $`\mathrm{cosh}\left(z\right)=\text{tr}\left(e^{s\tau _j}g^{}\overline{g}^T\right)/2`$ and $`\mathrm{cosh}\left(z_0\right)=\text{tr}\left(g^{}\overline{g}^T\right)/2`$, then
$`<\psi _g^t,\widehat{h}_{AB}\psi _g^{}^t>`$ (3.100)
$`=`$ $`{\displaystyle \frac{g_{AB}^{}}{2\mathrm{sinh}\left(z_0\right)}}{\displaystyle \underset{j}{}}\left[\left(d_j+1\right)e^{\frac{t}{4}\left(d_j^2+d_j1/2\right)}+\left(d_j1\right)e^{\frac{t}{4}\left(d_j^2d_j1/2\right)}\right]\mathrm{sinh}\left(d_jz_0\right)`$
$``$ $`{\displaystyle \frac{1}{2}}\left(\tau _jg^{}\right)_{AB}\left({\displaystyle \frac{d}{ds}}\right)_{s=0}{\displaystyle \underset{j}{}}\left[e^{\frac{t}{4}\left(d_j^2+d_j1/2\right)}e^{\frac{t}{4}\left(d_j^2d_j1/2\right)}\right]{\displaystyle \frac{\mathrm{sinh}\left(d_jz\right)}{\mathrm{sinh}\left(z\right)}}`$
$`=`$ $`{\displaystyle \frac{g_{AB}^{}}{4\mathrm{sinh}\left(z_0\right)}}e^{t/8}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{tn^2/4}\left[\left(n+1\right)e^{nt/4}+\left(n1\right)e^{nt/4}\right]\left[e^{nz_0}e^{nz_0}\right]`$
$``$ $`{\displaystyle \frac{1}{2}}\left(\tau _jg^{}\right)_{AB}\left({\displaystyle \frac{d}{ds}}\right)_{s=0}{\displaystyle \frac{e^{t/8}}{2\mathrm{sinh}\left(z\right)}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{tn^2/4}\left[e^{nt/4}e^{nt/4}\right]\left[e^{nz}e^{nz}\right]`$
$`=`$ $`{\displaystyle \frac{g_{AB}^{}}{4\mathrm{sinh}\left(z_0\right)}}e^{t/8}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{tn^2/4}\{[ne^{n\left(z_0t/4\right)}+(n)e^{\left(n\right)\left(z_0t/4\right)}]+[ne^{n\left(z_0+t/4\right)}+(n)e^{\left(n\right)\left(z_0+t/4\right)}]`$
$`+[e^{n\left(z_0t/4\right)}+e^{\left(n\right)\left(z_0t/4\right)}][e^{n\left(z_0+t/4\right)}+e^{\left(n\right)\left(z_0+t/4\right)}]\}`$
$``$ $`{\displaystyle \frac{1}{2}}\left(\tau _jg^{}\right)_{AB}\left({\displaystyle \frac{d}{ds}}\right)_{s=0}{\displaystyle \frac{e^{t/8}}{2\mathrm{sinh}\left(z\right)}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{tn^2/4}\left\{\left[e^{n\left(zt/4\right)}+e^{\left(n\right)\left(zt/4\right)}\right]\left[e^{n\left(z+t/4\right)}+e^{\left(n\right)\left(z+t/4\right)}\right]\right\}`$
$`=`$ $`{\displaystyle \frac{g_{AB}^{}}{4\mathrm{sinh}\left(z_0\right)}}e^{t/8}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}e^{tn^2/4}\left[\left(n+1\right)e^{n\left(z_0t/4\right)}+\left(n1\right)e^{n\left(z_0+t/4\right)}\right]`$
$``$ $`{\displaystyle \frac{1}{2}}\left(\tau _jg^{}\right)_{AB}\left({\displaystyle \frac{d}{ds}}\right)_{s=0}{\displaystyle \frac{e^{t/8}}{2\mathrm{sinh}\left(z\right)}}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}e^{tn^2/4}\left[e^{n\left(zt/4\right)}e^{n\left(z+t/4\right)}\right]`$
where in the last step we have recognized that the terms in the curly brackets add up to zero at $`n=0`$ and that the terms with $`\left(n\right)`$ as argument can be taken care of by extending the series to negative values of $`n`$. Introducing
$$T:=\sqrt{t}/2,z_0^\pm =z_0/T\pm T,z^\pm =z/T\pm T$$
(3.101)
and remembering (3.31) we may write (3.100) in the form (notice that $`s`$ is to be replaced by $`s`$ as compared to (3.24))
$`<\psi _g^t,\widehat{h}_{AB}\psi _g^{}^t>`$ (3.102)
$`=`$ $`{\displaystyle \frac{g_{AB}^{}}{4\mathrm{sinh}\left(z_0\right)T}}e^{t/8}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}e^{\left(nT\right)^2}\left[\left(\left(nT\right)+T\right)e^{\left(nT\right)z_0^{}}+\left(\left(nT\right)T\right)e^{\left(nT\right)z_0^+}\right]`$
$``$ $`{\displaystyle \frac{1}{2}}\left(\tau _jg^{}\right)_{AB}{\displaystyle \frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{2\mathrm{sinh}\left(z_0\right)}}\left({\displaystyle \frac{d}{dz}}\right)_{z=z_0}{\displaystyle \frac{e^{t/8}}{2\mathrm{sinh}\left(z\right)}}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}e^{\left(nT\right)^2}\left[e^{\left(nT\right)z^{}}e^{\left(nT\right)z^+}\right]`$
An appeal to the Poisson summation formula now reveals that
$`<\psi _g^t,\widehat{h}_{AB}\psi _g^{}^t>`$
$`=`$ $`g_{AB}^{}{\displaystyle \frac{e^{t/8}}{4\mathrm{sinh}\left(z_0\right)T}}{\displaystyle \frac{2\pi }{T}}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\left[\left({\displaystyle \frac{z_0^{}\frac{2\pi in}{T}}{4\sqrt{\pi }}}+{\displaystyle \frac{T}{2\sqrt{\pi }}}\right)e^{\left(z_0^{}\frac{2\pi in}{T}\right)^2/4}+\left({\displaystyle \frac{z_0^+\frac{2\pi in}{T}}{4\sqrt{\pi }}}{\displaystyle \frac{T}{2\sqrt{\pi }}}\right)e^{\left(z_0^+\frac{2\pi in}{T}\right)^2/4}\right]`$
$``$ $`{\displaystyle \frac{e^{t/8}}{8\mathrm{sinh}\left(z_0\right)}}\left(\tau _jg^{}\right)_{AB}\text{tr}\left(\tau _jg^{}\overline{g}^T\right)\left({\displaystyle \frac{d}{dz}}\right)_{z=z_0}{\displaystyle \frac{1}{\mathrm{sinh}\left(z\right)}}{\displaystyle \frac{2\pi }{T}}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\left[{\displaystyle \frac{e^{\left(z^{}\frac{2\pi in}{T}\right)^2/4}}{2\sqrt{\pi }}}{\displaystyle \frac{e^{\left(z^+\frac{2\pi in}{T}\right)^2/4}}{2\sqrt{\pi }}}\right]`$
$`=`$ $`{\displaystyle \frac{\sqrt{\pi }e^{t/8}}{8\mathrm{sinh}\left(z_0\right)T^3}}\left[g_{AB}^{}\right]{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\left[\left(z_0+T^22\pi in\right)e^{\frac{\left(z_0T^22\pi in\right)^2}{t}}+\left(z_0T^22\pi in\right)e^{\frac{\left(z_0+T^22\pi in\right)^2}{t}}\right]`$
$``$ $`{\displaystyle \frac{\sqrt{\pi }e^{t/8}}{8\mathrm{sinh}\left(z_0\right)T^3}}\left[\left(\tau _jg^{}\right)_{AB}\text{tr}\left(\tau _jg^{}\overline{g}^T\right)\right]\left({\displaystyle \frac{d}{dz}}\right)_{z=z_0}{\displaystyle \frac{T^2}{\mathrm{sinh}\left(z\right)}}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\left[e^{\frac{\left(zT^22\pi in\right)^2}{t}}e^{\frac{\left(z+T^22\pi in\right)^2}{t}}\right]`$
$`=`$ $`{\displaystyle \frac{\sqrt{\pi }e^{t/8}}{8\mathrm{sinh}\left(z_0\right)T^3}}\left[g_{AB}^{}\right]{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\left[\left(z_0+T^22\pi in\right)e^{\frac{\left(z_0T^22\pi in\right)^2}{t}}+\left(z_0T^22\pi in\right)e^{\frac{\left(z_0+T^22\pi in\right)^2}{t}}\right]`$
$``$ $`{\displaystyle \frac{\sqrt{\pi }e^{t/8}}{8\mathrm{sinh}\left(z_0\right)T^3}}\left[\left(\tau _jg^{}\right)_{AB}\text{tr}\left(\tau _jg^{}\overline{g}^T\right)\right]T^2\{{\displaystyle \frac{\mathrm{cosh}\left(z_0\right)}{\mathrm{sinh}^2\left(z_0\right)}}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}[e^{\frac{\left(z_0T^22\pi in\right)^2}{t}}e^{\frac{\left(z_0+T^22\pi in\right)^2}{t}}]`$
$`+{\displaystyle \frac{1}{2T^2\mathrm{sinh}\left(z_0\right)}}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}[(z_0T^22\pi in)e^{\frac{\left(z_0T^22\pi in\right)^2}{t}}(z_0+T^22\pi in)e^{\frac{\left(z_0+T^22\pi in\right)^2}{t}}]\}`$
Recalling the norm of the coherent states from (3.1.2),
$$\left|\right|\psi _g^t||^2=\frac{\sqrt{\pi }e^{t/4}}{4\mathrm{sinh}\left(p\right)T^3}\underset{n}{}(p2\pi in)e^{\frac{\left(p2\pi in\right)^2}{t}}=:\frac{\sqrt{\pi }e^{t/4}}{4T^3}\frac{p}{\mathrm{sinh}\left(p\right)}e^{p^2/t}(1+K_t\left(p\right))$$
(3.104)
we arrive at the final exact formula for the matrix element of the holonomy operator (3.92) (all sums run over $`n\text{ }\mathrm{Z}`$)
$`<\widehat{h}_{AB}>_{gg^{}}^t`$ (3.105)
$`=`$ $`{\displaystyle \frac{\frac{e^{t/8}e^{\frac{p^2+\left(p^{}\right)^2}{2t}}}{2\mathrm{sinh}\left(z_0\right)}}{\sqrt{\frac{p}{\mathrm{sinh}\left(p\right)}\left(1+K_t\left(p\right)\right)\frac{p^{}}{\mathrm{sinh}\left(p^{}\right)}\left(1+K_t\left(p^{}\right)\right)}}}\times `$
$`\times `$ $`\{g_{AB}^{}{\displaystyle \underset{n}{}}[(z_0+T^22\pi in)e^{\frac{\left(z_0T^22\pi in\right)^2}{t}}+(z_0T^22\pi in)e^{\frac{\left(z_0+T^22\pi in\right)^2}{t}}]`$
$``$ $`\left(\tau _jg^{}\right)_{AB}{\displaystyle \frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{2\mathrm{sinh}\left(z_0\right)}}{\displaystyle \underset{n}{}}[(z_0T^22\pi in2T^2{\displaystyle \frac{\mathrm{cosh}\left(z_0\right)}{\mathrm{sinh}\left(z_0\right)}})e^{\frac{\left(z_0T^22\pi in\right)^2}{t}}`$
$`(z_0+T^22\pi in2T^2{\displaystyle \frac{\mathrm{cosh}\left(z_0\right)}{\mathrm{sinh}\left(z_0\right)}})e^{\frac{\left(z_0+T^22\pi in\right)^2}{t}}]\}`$
At this point we must again distinguish between the cases A) $`0\varphi \left(1c\right)\pi `$ and B) $`\left(1c\right)\pi \varphi \pi `$ for some $`c<1/2`$ where $`z_0=s+i\varphi `$.
Case A)
We have
$$\left(z_0\pm T^22\pi in\right)^2/t=\frac{z_0^2}{t}+\frac{t}{16}\frac{4\pi ^2n^2}{t}\frac{4\pi inz_0}{t}\pm \left(z_0/2i\pi n\right)$$
(3.106)
and thus can write (3.105) more explicitely as
$`<\widehat{h}_{AB}>_{gg^{}}^t`$ (3.107)
$`=`$ $`{\displaystyle \frac{\frac{e^{t/16}e^{\frac{p^2+\left(p^{}\right)^22z_0^2}{2t}}}{2\mathrm{sinh}\left(z_0\right)}}{\sqrt{\frac{p}{\mathrm{sinh}\left(p\right)}\left(1+K_t\left(p\right)\right)\frac{p^{}}{\mathrm{sinh}\left(p^{}\right)}\left(1+K_t\left(p^{}\right)\right)}}}\times `$
$`\times `$ $`\{g_{AB}^{}{\displaystyle \underset{n}{}}(1)^ne^{\frac{4\pi ^2n^2}{t}}e^{\frac{4\pi inz_0}{t}}[(z_0+T^22\pi in)e^{z_0/2}+(z_0T^22\pi in)e^{z_0/2}]`$
$``$ $`\left(\tau _jg^{}\right)_{AB}{\displaystyle \frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{2\mathrm{sinh}\left(z_0\right)}}{\displaystyle \underset{n}{}}(1)^ne^{\frac{4\pi ^2n^2}{t}}e^{\frac{4\pi inz_0}{t}}[(z_0T^22\pi in2T^2{\displaystyle \frac{\mathrm{cosh}\left(z_0\right)}{\mathrm{sinh}\left(z_0\right)}})e^{z_0/2}`$
$`(z_0+T^22\pi in2T^2{\displaystyle \frac{\mathrm{cosh}\left(z_0\right)}{\mathrm{sinh}\left(z_0\right)}})e^{z_0/2}]\}`$
Let us focus on the curly bracket in (3.107) for which we find, after some considerable amount of algebra,
$`\{.\}=2z_0\times `$ (3.108)
$`\times `$ $`\{g_{AB}^{}\{\mathrm{cosh}(z_0/2)+[{\displaystyle \frac{T^2}{2}}{\displaystyle \frac{\mathrm{sinh}\left(z_0/2\right)}{z_0/2}}]+2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1)^ne^{\frac{4\pi ^2n^2}{t}}\times `$
$`\times `$ $`(\left[(\mathrm{cosh}(z_0/2){\displaystyle \frac{T^2}{2}}{\displaystyle \frac{\mathrm{sinh}\left(z_0/2\right)}{z_0/2}})\mathrm{cos}\left({\displaystyle \frac{4\pi nz_0}{t}}\right)\right]+[{\displaystyle \frac{8\pi ^2n^2}{t}}\mathrm{cosh}(z_0/2){\displaystyle \frac{\mathrm{sin}\left(\frac{4\pi nz_0}{t}\right)}{\frac{4\pi nz_0}{t}}}])\}`$
$``$ $`\left(\tau _jg^{}\right)_{AB}z_0{\displaystyle \frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{2\mathrm{sinh}\left(z_0\right)}}\{\mathrm{sinh}(z_0/2)/z_0+\left[2T^2{\displaystyle \frac{\text{coth}\left(z_0\right)\mathrm{sinh}\left(z_0/2\right)\frac{1}{2}\mathrm{cosh}\left(z_0/2\right)}{z_0^2}}\right]`$
$`+2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1)^ne^{\frac{4\pi ^2n^2}{t}}\times `$
$`\times (\left[(\mathrm{sinh}(z_0/2)/z_0+2T^2{\displaystyle \frac{\text{coth}\left(z_0\right)\mathrm{sinh}\left(z_0/2\right)\frac{1}{2}\mathrm{cosh}\left(z_0/2\right)}{z_0^2}})\mathrm{cos}\left({\displaystyle \frac{4\pi nz_0}{t}}\right)\right]`$
$`+\left[{\displaystyle \frac{\left(2\pi n\right)^2}{t}}{\displaystyle \frac{\mathrm{sinh}\left(z_0/2\right)}{z_0/2}}{\displaystyle \frac{\mathrm{sin}\left(\frac{4\pi nz_0}{t}\right)}{\frac{4\pi nz_0}{t}}}\right])\}\}`$
$`=:`$ $`2z_0\{g_{AB}^{}\{\mathrm{cosh}(z_0/2)+\left[I_1\right]+2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1)^ne^{\frac{4\pi ^2n^2}{t}}(\left[I_2\right]+\left[I_3\right])\}`$
$``$ $`\left(\tau _jg^{}\right)_{AB}z_0{\displaystyle \frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{2\mathrm{sinh}\left(z_0\right)}}\{{\displaystyle \frac{\mathrm{sinh}\left(z_0/2\right)}{z_0}}+2T^2\left[J_1\right]+2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1)^ne^{\frac{4\pi ^2n^2}{t}}(\left[J_2\right]+\left[J_3\right])\}\}`$
where we have abbreviated the terms in the square brackets in the first equality by $`I_1,I_2,I_3,J_1,J_2,J_3`$ in this order since we wish to estimate them separately.
Combining (3.107) with (3.108) yields
$`<\widehat{h}_{AB}>_{gg^{}}^t`$
$`=`$ $`{\displaystyle \frac{\frac{e^{t/16}e^{\frac{p^2+\left(p^{}\right)^22z_0^2}{2t}}z_0}{\mathrm{sinh}\left(z_0\right)}}{\sqrt{\frac{p}{\mathrm{sinh}\left(p\right)}\left(1+K_t\left(p\right)\right)\frac{p^{}}{\mathrm{sinh}\left(p^{}\right)}\left(1+K_t\left(p^{}\right)\right)}}}\times `$
$`\times `$ $`\{g_{AB}^{}\{\mathrm{cosh}(z_0/2)+\left[I_1\right]+2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1)^ne^{\frac{4\pi ^2n^2}{t}}(\left[I_2\right]+\left[I_3\right])\}`$
$``$ $`\left(\tau _jg^{}\right)_{AB}z_0{\displaystyle \frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{2\mathrm{sinh}\left(z_0\right)}}\{\mathrm{sinh}(z_0/2)/z_0+\left[J_1\right]+2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1)^ne^{\frac{4\pi ^2n^2}{t}}(\left[J_2\right]+\left[J_3\right])\}\}`$
and proceeding as in section 3.1.2 we define
$`\mathrm{\Delta }<\widehat{h}_{AB}>_{gg^{}}^t`$ (3.110)
$`:=`$ $`<\widehat{h}_{AB}>_{gg^{}}^t{\displaystyle \frac{\frac{e^{t/16}e^{\frac{p^2+\left(p^{}\right)^22z_0^2}{2t}}z_0}{\mathrm{sinh}\left(z_0\right)}}{\sqrt{\frac{p}{\mathrm{sinh}\left(p\right)}\left(1+K_t\left(p\right)\right)\frac{p^{}}{\mathrm{sinh}\left(p^{}\right)}\left(1+K_t\left(p^{}\right)\right)}}}\times `$
$`\times `$ $`[g_{AB}^{}\{\mathrm{cosh}(z_0/2)+\left(\tau _jg^{}\right)_{AB}{\displaystyle \frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{2\mathrm{sinh}\left(z_0\right)}}\mathrm{sinh}(z_0/2)]`$
$`=`$ $`{\displaystyle \frac{\frac{e^{t/16}e^{\frac{p^2+\left(p^{}\right)^22z_0^2}{2t}}z_0}{\mathrm{sinh}\left(z_0\right)}}{\sqrt{\frac{p}{\mathrm{sinh}\left(p\right)}\left(1+K_t\left(p\right)\right)\frac{p^{}}{\mathrm{sinh}\left(p^{}\right)}\left(1+K_t\left(p^{}\right)\right)}}}\times `$
$`\times `$ $`\{g_{AB}^{}\{\left[I_1\right]+2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1)^ne^{\frac{4\pi ^2n^2}{t}}(\left[I_2\right]+\left[I_3\right])\}`$
$``$ $`\left(\tau _jg^{}\right)_{AB}z_0{\displaystyle \frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{2\mathrm{sinh}\left(z_0\right)}}\{2T^2\left[J_1\right]+2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1)^ne^{\frac{4\pi ^2n^2}{t}}(\left[J_2\right]+\left[J_3\right])\}\}`$
The tools to estimate these terms have already been laid in section 3.1.2 so that we can be brief here. We just need the following result.
###### Lemma 3.6
For any $`z_0=s+i\varphi ,\mathrm{\hspace{0.33em}0}\varphi \pi (1c)`$ we have
$$\left|J_1\right|\mathrm{cosh}\left(\right|z_0|/2)[\sqrt{2}k_c^{}+\frac{k_ck_{\frac{1+c}{2}}^{}}{4}]=:\stackrel{~}{k}_c\mathrm{cosh}\left(\right|z_0|/2)$$
(3.111)
Proof of Lemma 3.6 :
We easily establish the following identity
$$J_1=\frac{1}{2}C\left(z_0\right)\frac{\mathrm{sinh}\left(z_0/2\right)}{z_0/2}\frac{S\left(z_0\right)}{2}\mathrm{cosh}\left(z_0/2\right)+\frac{1}{8}\frac{z_0}{\mathrm{sinh}\left(z_0\right)}\frac{\mathrm{sinh}\left(z_0/2\right)}{z_0/2}\left[S\left(z_0/2\right)C\left(z_0/2\right)\right]$$
(3.112)
Noticing that $`\mathrm{}\left(z_0/2\right)\pi \left(1c\right)/2=\pi \left(1\frac{1+c}{2}\right)`$ the assertion follows from the lemmas of the previous subsection.
$`\mathrm{}`$
With the help of this lemma one finds
$`\left|I_1\right|`$ $`=`$ $`\left|{\displaystyle \frac{T^2}{2}}{\displaystyle \frac{\mathrm{sinh}\left(z_0/2\right)}{z_0/2}}\right|{\displaystyle \frac{T^2}{2}}\mathrm{cosh}\left(\left|z_0\right|/2\right)`$
$`\left|I_2\right|`$ $`=`$ $`\left|\left(\mathrm{cosh}\left(z_0/2\right){\displaystyle \frac{T^2}{2}}{\displaystyle \frac{\mathrm{sinh}\left(z_0/2\right)}{z_0/2}}\right)\mathrm{cos}\left({\displaystyle \frac{4\pi nz_0}{t}}\right)\right|`$
$``$ $`\mathrm{cosh}\left(\right|z_0|/2)(1+T^2/2)\left|\mathrm{cos}\left({\displaystyle \frac{4\pi nz_0}{t}}\right)\right|\mathrm{cosh}\left(\right|z_0|/2)(1+T^2/)e^{\frac{4\pi ^2n\left(1c\right)}{t}}`$
$`\left|I_3\right|`$ $`=`$ $`\left|{\displaystyle \frac{8\pi ^2n^2}{t}}\mathrm{cosh}\left(z_0/2\right){\displaystyle \frac{\mathrm{sin}\left(\frac{4\pi nz_0}{t}\right)}{\frac{4\pi nz_0}{t}}}\right|`$
$``$ $`{\displaystyle \frac{8\pi ^2n^2}{t}}\mathrm{cosh}\left(\left|z_0\right|/2\right)\left|{\displaystyle \frac{\mathrm{sin}\left(\frac{4\pi nz_0}{t}\right)}{\frac{4\pi nz_0}{t}}}\right|{\displaystyle \frac{16\pi ^2n^2}{t}}\mathrm{cosh}\left(\left|z_0\right|/2\right)e^{\frac{4\pi ^2n\left(1c\right)}{t}}`$
$`\left|J_2\right|`$ $`=`$ $`\left|\left(\mathrm{sinh}\left(z_0/2\right)/z_0+2T^2J_1\right)\mathrm{cos}\left({\displaystyle \frac{4\pi nz_0}{t}}\right)\right|`$
$``$ $`\left|(2\mathrm{cosh}\left(\right|z_0|/2)+2T^2|J_1\left|\right)\right|\mathrm{cos}\left({\displaystyle \frac{4\pi nz_0}{t}}\right)|`$
$``$ $`2\mathrm{cosh}\left(\left|z_0\right|/2\right)\left(1+T^2\stackrel{~}{k}_c\right)e^{\frac{4\pi ^2n\left(1c\right)}{t}}`$
$`\left|J_3\right|`$ $`=`$ $`\left|{\displaystyle \frac{\left(2\pi n\right)^2}{t}}{\displaystyle \frac{\mathrm{sinh}\left(z_0/2\right)}{z_0/2}}{\displaystyle \frac{\mathrm{sin}\left(\frac{4\pi nz_0}{t}\right)}{\frac{4\pi nz_0}{t}}}\right|`$ (3.113)
$``$ $`{\displaystyle \frac{\left(2\pi n\right)^2}{t}}\mathrm{cosh}\left(\left|z_0\right|/2\right)e^{\frac{4\pi ^2n\left(1c\right)}{t}}`$
Observing the basic estimates (exploit $`\left(\tau _j\right)^2=1`$)
$`\left|g_{AB}^{}\right|^2{\displaystyle \underset{A,B}{}}\left|g_{AB}^{}\right|^2=\text{tr}\left(g^{}\left(\overline{g}^{}\right)^T\right)=2\mathrm{cosh}\left(p^{}\right)4\mathrm{cosh}^2\left(p^{}/2\right)\text{ and}`$
$`\left|\left(\tau _jg^{}\right)_{AB}\right|^2\text{tr}\left(\tau _jg^{}\left(\overline{\tau _jg^{}}\right)^T\right)=\text{tr}\left(\tau _j^2g^{}\left(\overline{g}^{}\right)^T\right)4\mathrm{cosh}^2\left(p^{}/2\right)`$ (3.114)
and employing the estimate for the denominator of (3.110) from the previous section together with (3.77) we obtain as an estimate (notation the same as in section 3.1.2)
$`\left|\mathrm{\Delta }<\widehat{h}_{AB}>_{gg^{}}^t\right|`$
$``$ $`{\displaystyle \frac{e^{t/16}e^{\frac{p^2+\left(p^{}\right)^22\mathrm{}\left(z_0^2\right)}{2t}}k_c\left|\frac{s}{\mathrm{sinh}\left(s\right)}\right|\mathrm{cosh}\left(\left|z_0\right|/2\right)}{\left(1K_t\right)\sqrt{\frac{p}{\mathrm{sinh}\left(p\right)}\frac{p^{}}{\mathrm{sinh}\left(p^{}\right)}}}}\times `$
$`\times `$ $`\left\{\right|g_{AB}^{}|\{{\displaystyle \frac{T^2}{2}}+2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{4\pi ^2n^2}{t}}e^{\frac{4\pi ^2n\left(1c\right)}{t}}[1+{\displaystyle \frac{T^2}{2}}+{\displaystyle \frac{16\pi ^2n^2}{t}}]\}`$
$`+`$ $`\left|\left(\tau _jg^{}\right)_{AB}\right|\left|z_0{\displaystyle \frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{2\mathrm{sinh}\left(z_0\right)}}\right|\{2T^2tildek_c+2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{4\pi ^2n^2}{t}}e^{\frac{4\pi ^2n\left(1c\right)}{t}}[2(1+T^2\stackrel{~}{k}_c)+{\displaystyle \frac{\left(2\pi n\right)^2}{t}}]\}\}`$
$``$ $`{\displaystyle \frac{2e^{t/16}e^{\frac{\mathrm{\Delta }^2+2\delta ^2+2\stackrel{~}{\theta }^2}{2t}}k_c\left|\frac{s}{\mathrm{sinh}\left(s\right)}\right|\mathrm{cosh}\left(\left|z_0\right|/2\right)\mathrm{cosh}\left(p^{}/2\right)}{\left(1K_t\right)\sqrt{\frac{p}{\mathrm{sinh}\left(p\right)}\frac{p^{}}{\mathrm{sinh}\left(p^{}\right)}}}}\times `$
$`\times `$ $`\{\{{\displaystyle \frac{T^2}{2}}+2e^{\frac{4\pi ^2c}{t}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}e^{\frac{4\pi ^2n^2}{t}}[1+{\displaystyle \frac{T^2}{2}}+{\displaystyle \frac{16\pi ^2\left(n+1\right)^2}{t}}]\}`$
$`+`$ $`4k_c\left|{\displaystyle \frac{s}{\mathrm{sinh}\left(s\right)}}\right|\mathrm{cosh}\left({\displaystyle \frac{p+p^{}}{2}}\right)\{2T^2tildek_c+2e^{\frac{4\pi ^2c}{t}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}e^{\frac{4\pi ^2n^2}{t}}[2(1+T^2\stackrel{~}{k}_c)+{\displaystyle \frac{\left(4\pi \left(n+1\right)\right)^2}{t}}]\}\}`$
In discussing the behaviour of this function as $`p^{}`$ becomes large (integrability) we need to separate again the regions described by a) bounded $`\delta `$ (i.e. $`s\stackrel{~}{p}/2p^{}/2`$) and b) unbounded $`\delta `$ (i.e. $`s/p^{}`$ vanishes as $`p^{}\mathrm{}`$ or $`s`$ is bounded) similar as in the previous section. In case a) (3.1.3) grows as $`\mathrm{cosh}\left(\left|z_0\right|/2\right)\sqrt{p^{}}^3\mathrm{cosh}\left(p^{}/4\right)\sqrt{p^{}}^3`$ times the Gaussian in $`\mathrm{\Delta }^2+2\stackrel{~}{\theta }^2`$. In case b) it is damped by the Gaussian in $`\delta ^2`$ and is exponentially small times the Gaussian in $`\mathrm{\Delta }^2+2\stackrel{~}{\theta }^2`$.
Finally looking at the terms inside the two inner curly brackets of (3.1.3) we see that they are up to a numerical factor given by constants of the form $`t+K_t\left(c\right)`$ where $`K_t\left(c\right)`$ vanishes exponentially fast as $`t0`$. We conclude that the integral over $`g^{}`$ of (3.1.3) in the range $`0\varphi \left(1c\right)\pi `$ results in a function of $`g`$ which is at most exponentially growing with $`p`$ times a constant that approaches $`t`$ exponentially fast.
Case B)
Following the by now already familiar trick we will now write (3.105) in terms of $`z_0^{}=z_0i\pi =si\left(\pi \varphi \right)=si\varphi ^{}`$ with $`0\varphi ^{}c\pi `$. This gives (observe that $`\mathrm{sinh}\left(z_0\right)=\mathrm{sinh}\left(z_0^{}\right),\mathrm{cosh}\left(z_0\right)=\mathrm{cosh}\left(z_0^{}\right)`$)
$`<\widehat{h}_{AB}>_{gg^{}}^t`$ (3.116)
$`=`$ $`{\displaystyle \frac{\frac{e^{t/8}e^{\frac{p^2+\left(p^{}\right)^2}{2t}}}{2\mathrm{sinh}\left(z_0^{}\right)}}{\sqrt{\frac{p}{\mathrm{sinh}\left(p\right)}\left(1+K_t\left(p\right)\right)\frac{p^{}}{\mathrm{sinh}\left(p^{}\right)}\left(1+K_t\left(p^{}\right)\right)}}}\times `$
$`\times `$ $`\{g_{AB}^{}{\displaystyle \underset{n}{}}[(z_0^{}+T^2\pi i(2n1))e^{\frac{\left(z_0^{}T^2\pi i\left(2n1\right)\right)^2}{t}}+(z_0^{}T^2\pi i(2n1))e^{\frac{\left(z_0^{}+T^2\pi i\left(2n1\right)\right)^2}{t}}]`$
$`+`$ $`\left(\tau _jg^{}\right)_{AB}{\displaystyle \frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{2\mathrm{sinh}\left(z_0^{}\right)}}{\displaystyle \underset{n}{}}[(z_0^{}T^2\pi i(2n1)2T^2{\displaystyle \frac{\mathrm{cosh}\left(z_0^{}\right)}{\mathrm{sinh}\left(z_0^{}\right)}})e^{\frac{\left(z_0^{}T^2\pi i\left(2n1\right)\right)^2}{t}}`$
$`(z_0^{}+T^2\pi i(2n1)2T^2{\displaystyle \frac{\mathrm{cosh}\left(z_0^{}\right)}{\mathrm{sinh}\left(z_0^{}\right)}})e^{\frac{\left(z_0^{}+T^2\pi i\left(2n1\right)\right)^2}{t}}]\}`$
$`=`$ $`{\displaystyle \frac{\frac{e^{t/8}e^{\frac{p^2+\left(p^{}\right)^2}{2t}}}{2\mathrm{sinh}\left(z_0^{}\right)}}{\sqrt{\frac{p}{\mathrm{sinh}\left(p\right)}\left(1+K_t\left(p\right)\right)\frac{p^{}}{\mathrm{sinh}\left(p^{}\right)}\left(1+K_t\left(p^{}\right)\right)}}}\times `$
$`\times `$ $`\{g_{AB}^{}{\displaystyle \underset{n=\text{odd}}{}}[(z_0^{}+T^2\pi in)e^{\frac{\left(z_0^{}T^2\pi in\right)^2}{t}}+(z_0^{}T^2\pi in)e^{\frac{\left(z_0^{}+T^2\pi in\right)^2}{t}}]`$
$`+`$ $`\left(\tau _jg^{}\right)_{AB}{\displaystyle \frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{2\mathrm{sinh}\left(z_0^{}\right)}}{\displaystyle \underset{n=\text{odd}}{}}[(z_0^{}T^2\pi in2T^2{\displaystyle \frac{\mathrm{cosh}\left(z_0^{}\right)}{\mathrm{sinh}\left(z_0^{}\right)}})e^{\frac{\left(z_0^{}T^2\pi in\right)^2}{t}}`$
$`(z_0^{}+T^2\pi in2T^2{\displaystyle \frac{\mathrm{cosh}\left(z_0^{}\right)}{\mathrm{sinh}\left(z_0^{}\right)}})e^{\frac{\left(z_0^{}+T^2\pi in\right)^2}{t}}]\}`$
With
$$\frac{\left(z_0^{}\pm T^2\pi in\right)^2}{t}=\frac{\left(z_0^{}\right)^2}{t}+\frac{t}{16}\frac{\pi ^2n^2}{t}\frac{2\pi inz_0^{}}{t}\pm \left(z_0^{}/2i\pi n/2\right)$$
(3.117)
and since $`i^n=i^n`$ for $`n`$ odd we find after some pages of algebra
$`<\widehat{h}_{AB}>_{gg^{}}^t`$ (3.118)
$`=`$ $`{\displaystyle \frac{\frac{e^{t/16}e^{\frac{p^2+\left(p^{}\right)^22\left(z_0^{}\right)^2}{2t}}z_0^{}}{2\mathrm{sinh}\left(z_0^{}\right)}}{\sqrt{\frac{p}{\mathrm{sinh}\left(p\right)}\left(1+K_t\left(p\right)\right)\frac{p^{}}{\mathrm{sinh}\left(p^{}\right)}\left(1+K_t\left(p^{}\right)\right)}}}\times `$
$`\times `$ $`\{\{2g_{AB}^{}{\displaystyle \underset{n=1,\text{odd}}{\overset{\mathrm{}}{}}}i^ne^{\pi ^2n^2/t}(\left[2i\mathrm{sinh}(z_0^{}/2)\mathrm{sin}(2\pi nz_0^{}/t)\right]`$
$`+[\pi n\mathrm{cosh}(z_0^{}/2){\displaystyle \frac{\mathrm{sin}\left(2\pi nz_0^{}/t\right)}{2\pi nz_0^{}/t}}]+\left[4\pi in{\displaystyle \frac{\mathrm{sinh}\left(z_0^{}/2\right)}{z_0^{}/2}}\mathrm{cos}(2\pi nz_0^{}/t)\right])\}`$
$`+`$ $`\{\left(\tau _jg^{}\right)_{AB}z_0^{}{\displaystyle \frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{2\mathrm{sinh}\left(z_0^{}\right)}}{\displaystyle \underset{n=1,\text{odd}}{\overset{\mathrm{}}{}}}i^ne^{\pi ^2n^2/t}([{\displaystyle \frac{4\pi in}{t}}(2\mathrm{cosh}(z_0^{}/2)`$
$`+T^2{\displaystyle \frac{\mathrm{sinh}\left(z_0^{}/2\right)}{z_0^{}/2}}){\displaystyle \frac{\mathrm{sin}\left(2\pi nz_0^{}/t\right)}{2\pi nz_0^{}/t}}]`$
$`+\left[2it{\displaystyle \frac{\text{coth}\left(z_0^{}\right)\mathrm{sin}\left(2\pi nz_0^{}/t\right)\frac{2\pi n}{t}\mathrm{cos}\left(2\pi nz_0^{}/t\right)}{\left(z_0^{}\right)^2}}\mathrm{cosh}(z_0^{}/2)\right])\}\}`$
In estimating (3.118) the only term that is superficially non-regular at $`z_0^{}=0`$ is the last fraction in the second inner curly bracket. However, using the identity
$`K`$ $`:=`$ $`{\displaystyle \frac{\text{coth}\left(z_0^{}\right)\mathrm{sin}\left(2\pi nz_0^{}/t\right)\frac{2\pi n}{t}\mathrm{cos}\left(2\pi nz_0^{}/t\right)}{\left(z_0^{}\right)^2}}={\displaystyle \frac{2\pi n}{t}}\times `$ (3.119)
$`\times `$ $`\left\{\left\{\left[C\left(z_0^{}\right){\displaystyle \frac{\mathrm{sin}\left(\stackrel{~}{z}_0^{}\right)}{\stackrel{~}{z}_0^{}}}S\left(z_0^{}\right)\right]+\left({\displaystyle \frac{2\pi n}{t}}\right)^2\left[s\left(\stackrel{~}{z}_0^{}\right){\displaystyle \frac{z_0^{}}{\mathrm{sinh}\left(z_0^{}\right)}}c\left(\stackrel{~}{z}_0^{}\right)\right]\right\}\right\}`$
where $`\stackrel{~}{z}_0^{}=2\pi nz_0^{}/t`$ and employing the estimates (3.66), (3.65) as well as lemmata 3.2, 3.4 it is not difficult to show that in the range $`0\varphi ^{}c\pi `$
$$\left|K\right|\frac{16\pi n}{t}\left[k_{1c}^{}+\left(\frac{2\pi n}{t}\right)^2k_{1c}\right]e^{2\pi ^2nc/t}$$
(3.120)
With these preparations and using previous results we can finish the estimate of (3.118)
$`\left|<\widehat{h}_{AB}>_{gg^{}}^t\right|`$
$``$ $`{\displaystyle \frac{\frac{e^{t/16}e^{\frac{p^2+\left(p^{}\right)^22\mathrm{}\left(\left(z_0^{}\right)^2\right)}{2t}}\left|z_0^{}\right|}{2\left|\mathrm{sinh}\left(z_0^{}\right)\right|}}{\left(1K_t\right)\sqrt{\frac{p}{\mathrm{sinh}\left(p\right)}\frac{p^{}}{\mathrm{sinh}\left(p^{}\right)}}}}\times `$
$`\times `$ $`\{\left\{2\right|g_{AB}^{}|{\displaystyle \underset{n=1,\text{odd}}{\overset{\mathrm{}}{}}}e^{\pi ^2n^2/t}(\left[2\mathrm{sinh}\left(\right|z_0^{}|/2)\right|\mathrm{sin}(2\pi nz_0^{}/t)\left|\right]`$
$`+\left[\pi n\mathrm{cosh}\left(\right|z_0^{}|/2)\right|{\displaystyle \frac{\mathrm{sin}\left(2\pi nz_0^{}/t\right)}{2\pi nz_0^{}/t}}\left|\right]+\left[4\pi n{\displaystyle \frac{\mathrm{sinh}\left(\left|z_0^{}\right|/2\right)}{\left|z_0^{}\right|/2}}\right|\mathrm{cos}(2\pi nz_0^{}/t)\left|\right])\}`$
$`+`$ $`\left\{\right|\left(\tau _jg^{}\right)_{AB}\left|\right|z_0^{}{\displaystyle \frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{2\mathrm{sinh}\left(z_0^{}\right)}}|{\displaystyle \underset{n=1,\text{odd}}{\overset{\mathrm{}}{}}}e^{\pi ^2n^2/t}\times `$
$`\times (\left[{\displaystyle \frac{4\pi n}{t}}(2\mathrm{cosh}\left(\right|z_0^{}|/2)+T^2{\displaystyle \frac{\mathrm{sinh}\left(\left|z_0^{}\right|/2\right)}{\left|z_0^{}\right|/2}})\right|{\displaystyle \frac{\mathrm{sin}\left(2\pi nz_0^{}/t\right)}{2\pi nz_0^{}/t}}\left|\right]+\left[2t\right|K\left|\mathrm{cosh}\left(\right|z_0^{}|/2)\right])\}\}`$
$``$ $`\mathrm{cosh}\left(\right|z_0^{}|/2){\displaystyle \frac{\frac{e^{t/16}e^{\frac{\mathrm{\Delta }^2+2\left(\stackrel{~}{p}^2/4s^2+\left(\varphi ^{}\right)^2\right)}{2t}}\left|z_0^{}\right|}{2\left|\mathrm{sinh}\left(z_0^{}\right)\right|}}{\left(1K_t\right)\sqrt{\frac{p}{\mathrm{sinh}\left(p\right)}\frac{p^{}}{\mathrm{sinh}\left(p^{}\right)}}}}\times `$
$`\times `$ $`\{\left\{2\right|g_{AB}^{}\left|{\displaystyle \underset{n=1,\text{odd}}{\overset{\mathrm{}}{}}}e^{\pi ^2n^2/t}e^{2\pi ^2nc/t}(2+\pi n+4\pi n)\right\}`$
$`+`$ $`\left\{\right|\left(\tau _jg^{}\right)_{AB}\left|\right|z_0^{}{\displaystyle \frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{2\mathrm{sinh}\left(z_0^{}\right)}}\left|{\displaystyle \underset{n=1,\text{odd}}{\overset{\mathrm{}}{}}}e^{\pi ^2n^2/t}e^{2\pi ^2nc/t}({\displaystyle \frac{4\pi n}{t}}[2+T^2]+32\pi n[k_{1c}^{}+\left({\displaystyle \frac{2\pi n}{t}}\right)^2k_{1c}])\right\}\}`$
$``$ $`\mathrm{cosh}\left(\right|z_0^{}|/2)k_{1c}{\displaystyle \frac{s}{\mathrm{sinh}\left(s\right)}}\mathrm{cosh}(p^{}/2){\displaystyle \frac{e^{t/16}e^{\frac{\mathrm{\Delta }^2+2\left(\stackrel{~}{p}^2/4s^2+\left(\varphi ^{}\right)^2\right)}{2t}}}{\left(1K_t\right)\sqrt{\frac{p}{\mathrm{sinh}\left(p\right)}\frac{p^{}}{\mathrm{sinh}\left(p^{}\right)}}}}\times `$
$`\times `$ $`\{\left\{2{\displaystyle \underset{n=1,\text{odd}}{\overset{\mathrm{}}{}}}e^{\pi ^2n^2/t}e^{2\pi ^2nc/t}(2+5\pi n)\right\}`$
$`+`$ $`\left\{4k_{1c}{\displaystyle \frac{s}{\mathrm{sinh}\left(s\right)}}\mathrm{cosh}\left({\displaystyle \frac{p+p^{}}{2}}\right){\displaystyle \underset{n=1,\text{odd}}{\overset{\mathrm{}}{}}}e^{\pi ^2n^2/t}e^{2\pi ^2nc/t}({\displaystyle \frac{4\pi n}{t}}[2+T^2]+32\pi n[k_{1c}^{}+\left({\displaystyle \frac{2\pi n}{t}}\right)^2k_{1c}])\right\}\}`$
$``$ $`\mathrm{cosh}\left(\right|z_0^{}|/2)k_{1c}{\displaystyle \frac{s}{\mathrm{sinh}\left(s\right)}}\mathrm{cosh}(p^{}/2){\displaystyle \frac{e^{t/16}e^{\frac{\mathrm{\Delta }^2+2\left(\stackrel{~}{p}^2/4s^2+\left(\varphi ^{}\right)^2+\left(12c\right)\pi ^2\right)}{2t}}}{\left(1K_t\right)\sqrt{\frac{p}{\mathrm{sinh}\left(p\right)}\frac{p^{}}{\mathrm{sinh}\left(p^{}\right)}}}}\times `$
$`\times `$ $`\{\left\{2{\displaystyle \underset{n=0,\text{even}}{\overset{\mathrm{}}{}}}e^{\pi ^2n^2/t}(2+5\pi (n+1))\right\}`$
$`+`$ $`\{4k_{1c}{\displaystyle \frac{s}{\mathrm{sinh}\left(s\right)}}\mathrm{cosh}\left({\displaystyle \frac{p+p^{}}{2}}\right){\displaystyle \underset{n=0,\text{even}}{\overset{\mathrm{}}{}}}e^{\pi ^2n^2/t}({\displaystyle \frac{4\pi \left(n+1\right)}{t}}[2+T^2]`$
$`+32\pi (n+1)[k_{1c}^{}+\left({\displaystyle \frac{2\pi \left(n+1\right)}{t}}\right)^2k_{1c}])\}\}`$
where we have used in the last step that for any $`n1`$
$$\pi ^2n^2/t+2\pi ^2nc/t=\pi ^2n\left(n1\right)/t\pi ^2n\left(12c\right)/t\pi ^2\left(n1\right)^2/t\pi ^2\left(12c\right)/t$$
(3.122)
using that $`2c<1`$ and $`\left(n1\right)^2n\left(n1\right)`$ valid for $`n1`$. Consider now the piece of the argument of the Gaussian given by
$`\left[\stackrel{~}{p}^2/4s^2+\left(\varphi ^{}\right)^2+\left(12c\right)\pi ^2\right]=\left[\stackrel{~}{p}^2/4s^2+\varphi ^2\stackrel{~}{\theta }^2\right]\stackrel{~}{\theta }^2+2\pi \varphi 2\left(1c\right)\pi ^2`$ (3.123)
$``$ $`\delta ^2\stackrel{~}{\theta }^2+2\pi \varphi 2\left(1c\right)\pi ^2\delta ^2\stackrel{~}{\theta }^2+2c\pi ^2=\delta ^2\left(1d\right)\stackrel{~}{\theta }^2d\stackrel{~}{\theta }^2+2c\pi ^2`$
$``$ $`\delta ^2\left(1d\right)\stackrel{~}{\theta }^2\left(d/42c\right)\pi ^2`$
where in the first inequality we used $`\varphi ^{}=\pi \varphi `$ and (3.81), in the second $`\varphi \pi `$ and in the third that $`\varphi </=/>\pi /2`$ iff $`\stackrel{~}{\theta }</=/>\pi /2`$ so that for $`2c<1`$ we have $`\stackrel{~}{\theta }>\pi /2`$. Here $`0<d<1`$ is an arbitrary real number. Choosing $`c<d/8`$, say $`c=d/16`$ and $`d=1/2`$ for definiteness we can complete the estimate (3.127) by writing
$`\left|<\widehat{h}_{AB}>_{gg^{}}^t\right|`$ (3.124)
$``$ $`\mathrm{cosh}\left(\right|z_0^{}|/2)k_{1c}{\displaystyle \frac{s}{\mathrm{sinh}\left(s\right)}}\mathrm{cosh}(p^{}/2){\displaystyle \frac{e^{t/16}e^{\frac{\mathrm{\Delta }^2+2\delta ^2+\stackrel{~}{\theta }^2}{2t}}}{\left(1K_t\right)\sqrt{\frac{p}{\mathrm{sinh}\left(p\right)}\frac{p^{}}{\mathrm{sinh}\left(p^{}\right)}}}}\times `$
$`\times `$ $`\{\left\{2{\displaystyle \underset{n=0,\text{even}}{\overset{\mathrm{}}{}}}e^{\pi ^2n^2/t}(2+5\pi (n+1))\right\}`$
$`+`$ $`\{4k_{1c}{\displaystyle \frac{s}{\mathrm{sinh}\left(s\right)}}\mathrm{cosh}\left({\displaystyle \frac{p+p^{}}{2}}\right){\displaystyle \underset{n=0,\text{even}}{\overset{\mathrm{}}{}}}e^{\pi ^2n^2/t}({\displaystyle \frac{4\pi \left(n+1\right)}{t}}[2+T^2]`$
$`+32\pi (n+1)[k_{1c}^{}+\left({\displaystyle \frac{2\pi \left(n+1\right)}{t}}\right)^2k_{1c}])\}\}`$
Comparing (3.1.3) with (3.124) we see that the overall structure is completely identical, the only essential difference being that $`2\stackrel{~}{\theta }^2`$ in the exponent of the Gaussian is replaced by $`\stackrel{~}{\theta }^2`$. Since now $`\stackrel{~}{\theta }\pi /2`$ we conclude that the integral of (3.1.3) over $`g^{}`$ with respect to $`\mathrm{\Omega }/t^3`$ exists, resulting in a function of $`g`$ growing no stronger than exponentially with $`p`$ times a constant that vanishes exponentially fast with $`t0`$.
Summarizing, as far as the leading order behaviour (in $`t`$) of the matrix element of the holonomy operator is concerned, we can replace it by
$`<\widehat{h}_{AB}>_{gg^{}}^t`$ $``$ $`\mathrm{\Theta }(\pi (1c)\varphi ){\displaystyle \frac{\frac{e^{t/16}e^{\frac{p^2+\left(p^{}\right)^22z_0^2}{2t}}z_0}{\mathrm{sinh}\left(z_0\right)}}{\sqrt{\frac{p}{\mathrm{sinh}\left(p\right)}\left(1+K_t\left(p\right)\right)\frac{p^{}}{\mathrm{sinh}\left(p^{}\right)}\left(1+K_t\left(p^{}\right)\right)}}}\times `$ (3.125)
$`\times `$ $`\left[g_{AB}^{}\mathrm{cosh}\left(z_0/2\right)+\left(\tau _jg^{}\right)_{AB}{\displaystyle \frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{2\mathrm{sinh}\left(z_0\right)}}\mathrm{sinh}\left(z_0/2\right)\right]`$
The Gaussian displays a peak at $`g=g^{}`$ where $`\stackrel{~}{p}/2=p=p^{},\stackrel{~}{\alpha }=\stackrel{~}{\theta }=0`$ implying $`z_0=p`$ whence the prefactor in front of the square bracket in (3.125) becomes $`e^{t/16}`$ and the square bracket itself becomes
$`[.]`$ $`=`$ $`g_{AB}\mathrm{cosh}\left(p/2\right)+\left(\tau _jg\right)_{AB}{\displaystyle \frac{\text{tr}\left(\tau _jg\overline{g}^T\right)}{2\mathrm{sinh}\left(p\right)}}\mathrm{sinh}\left(p/2\right)`$ (3.126)
$`=`$ $`\left(\left[\mathrm{cosh}\left(p/2\right)+i{\displaystyle \frac{p_j}{p}}\mathrm{sinh}\left(p/2\right)\tau _j\right]g\right)_{AB}=\left(H\left(g\right)^1g\right)_{AB}=h\left(g\right)_{AB}`$
as expected. We now write the square bracket as
$$h\left(g\right)_{AB}+\left[g_{AB}^{}\mathrm{cosh}\left(z_0/2\right)+\left(\tau _jg^{}\right)_{AB}\frac{\text{tr}\left(\tau _jg^{}\overline{g}^T\right)}{2\mathrm{sinh}\left(z_0\right)}\mathrm{sinh}\left(z_0/2\right)h\left(g\right)_{AB}\right]$$
and do the Gaussian in $`g^{}`$. Notice that the $`g^{}`$ independent term is given by
$`h_{AB}\left(g\right)<1>_{gg^{}}^t`$ while the integral over the additional term will be at least of order $`t`$ since the first contribution to the Gaussian of width of order $`\sqrt{t}`$ comes from quadratic terms in $`\stackrel{}{p}\stackrel{}{p}^{},\stackrel{}{\theta }\stackrel{}{\theta }^{}`$. This gives rise to the second main theorem.
###### Theorem 3.3
The matrix elements of the holonomy operators with respect to coherent states can be estimated by
$$\left|\frac{<\psi _g^t,\widehat{h}_{AB}\psi _g^{}^t>}{\psi _g^t\psi _g^{}^t}h_{AB}\left(g\right)\frac{<\psi _g^t,\psi _g^{}^t>}{\psi _g^t\psi _g^{}^t}\right|tf^{}(\stackrel{}{p},\stackrel{}{p}^{})\frac{|<\psi _g^t,\psi _g^{}^t>|}{\psi _g^t\psi _g^{}^t}$$
(3.127)
where $`f^{}`$ is a function of $`\stackrel{}{p},\stackrel{}{p}^{}`$ growing no faster than exponentially in either or $`p,p^{}`$.
As a corollary to theorem (3.2) we obtain that the expectation value $`<\widehat{h}_{AB}>_{gg}^t`$ equals $`h_{AB}\left(g\right)`$ up to bounded corrections to $`h_{AB}\left(g\right)`$ that are proportional to $`t`$. We will actually calculate the exact correction in the next section by a different method. Notice that due to the Gaussian behaviour of the overlap function the exponential growth of $`f^{}`$ is irrelevant in computing the expectation value of operator monomials, the corrections of order at least $`t`$ are always integrable.
#### 3.1.4 Computation of Operator Monomial Expectation Values by a Different Method
One can compute the expectation value of operator monomials also by a different method which does not rely on the overcompleteness of the coherent states. To see how this works, notice first of all that due to $`[\widehat{p}_j,\widehat{h}_{AB}]=it\left(\tau _j\widehat{h}\right)_{AB}/2`$ every operator monomial can be reduced to finite linear combinations of operator monomials in the following “standard ordered form”
$$<(\widehat{O}_1..\widehat{O}_{m+n})_0>_g^t=<\widehat{p}_{j_1}..\widehat{p}_{j_m}\widehat{h}_{A_1B_1}..\widehat{h}_{A_nB_n}>_g^t$$
(3.128)
for some $`m,n0`$ and some ordering of the $`j_1,..,j_m`$. The idea is now to use the following identity established in for general semisimple compact $`G`$ (we display the case $`G=SU\left(2\right)`$), relating the annihilation and holonomy operators,
$$\widehat{g}_{AB}=e^{3t/8}\left(e^{i\widehat{p}_j\tau _j/2}\widehat{h}\right)_{AB}$$
(3.129)
and to use the eigenvalue property of the coherent states $`\widehat{g}_{AB}\psi _g^t=g_{AB}\psi _g^t`$. In order to do this we first have to invert (3.129) for $`\widehat{h}_{AB}`$. The naive guess turns out to be the correct one.
###### Theorem 3.4
The inversion of (3.129) reads
$$\widehat{h}_{AB}=e^{3t/8}\left(e^{i\widehat{p}_j\tau _j/2}\widehat{g}\right)_{AB}$$
(3.130)
The proof of that theorem rests on the following lemma.
###### Lemma 3.7
Define the strictly positive and self-adjoint operator $`\widehat{p}`$ by
$$\widehat{p}:=\sqrt{\widehat{p}_j\widehat{p}_j+\frac{t^2}{4}}$$
(3.131)
which commutes with all the $`\widehat{p}_j`$. Then the following operator identities hold for $`G=SU(2)`$ (obvious generalizations hold for groups of higher rank) :
$`e^{\pm i\widehat{p}_j\tau _j/2}`$ $`=`$ $`e^{\pm \frac{t}{4}}\left\{\left[\mathrm{cosh}\left(\widehat{p}/2\right){\displaystyle \frac{t}{2}}{\displaystyle \frac{\mathrm{sinh}\left(\widehat{p}/2\right)}{\widehat{p}}}\right]1_2\pm i\tau _j\widehat{p}_j{\displaystyle \frac{\mathrm{sinh}\left(\widehat{p}/2\right)}{\widehat{p}}}\right\}`$ (3.132)
$`e^{i\widehat{p}_j\tau _j/2}e^{i\widehat{p}_j\tau _j/2}=e^{i\widehat{p}_j\tau _j/2}e^{i\widehat{p}_j\tau _j/2}=1_2\widehat{1}_{}`$
Notice that no operator ordering ambiguities occur in (3.132).
Proof of Lemma 3.7 :
By definition (the operator $`\widehat{p}_j`$ is unbounded but the exponential can be defined by Nelson’s analytic vector theorem)
$$e^{\pm i\widehat{p}_j\tau _j/2}=\underset{n=0}{\overset{\mathrm{}}{}}\frac{\left(\pm i/2\right)^n}{n!}\left(\widehat{p}_j\tau _j\right)^n$$
(3.133)
Due to the Pauli matrix relation $`\tau _j\tau _k=\delta _{jk}1_2+ϵ_{jkl}\tau _l`$ and the commutator relation
$$ϵ_{jkl}\widehat{p}_j\widehat{p}_k\tau _l=\frac{1}{2}ϵ_{jkl}[\widehat{p}_j,\widehat{p}_k]\tau _l=it\frac{1}{2}ϵ_{jkl}ϵ_{jkm}\widehat{p}_m\tau _l=it\widehat{p}_j\tau _j$$
(3.134)
every power of the matrix valued operator $`\underset{¯}{\overset{^}{p}}=\tau _j\widehat{p}_j`$ can be written as a linear combination
$$\left(\underset{¯}{\overset{^}{p}}\right)^n=q_n\left(\widehat{x}\right)\underset{¯}{\overset{^}{p}}+r_n\left(\widehat{x}\right)1_2\widehat{1}_{}\text{ where }\widehat{x}=\widehat{p}_j\widehat{p}_j$$
(3.135)
and $`q_n,r_n`$ are polynomials which are inductively defined by
$$\left(\underset{¯}{\overset{^}{p}}\right)^{n+1}=q_{n+1}\left(\widehat{x}\right)\underset{¯}{\overset{^}{p}}+r_{n+1}\left(\widehat{x}\right)1_2\widehat{1}_{}=\left[q_n\left(\widehat{x}\right)\underset{¯}{\overset{^}{p}}+r_n\left(\widehat{x}\right)1_2\widehat{1}_{}\right]\underset{¯}{\overset{^}{p}}$$
(3.136)
from which we find
$$q_{n+1}=q_nit+r_n,r_{n+1}=\widehat{x}q_n,q_1=1,r_1=0$$
(3.137)
Notice that no operator ordering problems arise.
As one can check, the two-dimensional recursion defined in (3.137) is solved by
$$q_n=\frac{\lambda _+^n\lambda _{}^n}{\lambda _+\lambda _{}},r_n=\lambda _+\lambda _{}\frac{\lambda _+^{n1}\lambda _{}^{n1}}{\lambda _+\lambda _{}}$$
(3.138)
where
$$\lambda _\pm =i\left(\frac{t}{2}\pm \widehat{p}\right)$$
(3.139)
Inserted back into (3.133) gives
$$e^{\pm i\underset{¯}{\overset{^}{p}}/2}=\frac{1}{\lambda _+\lambda _{}}\left[\underset{¯}{\overset{^}{p}}\left(e^{\pm i\lambda _+/2}e^{\pm i\lambda _{}/2}\right)1_2\widehat{1}_{}\left(\lambda _{}e^{\pm i\lambda _+/2}\lambda _+e^{\pm i\lambda _{}/2}\right)\right]$$
(3.140)
and writing out $`\lambda _\pm `$ results in the first line of (3.132).
Using this result and that $`[\widehat{p}_j,\widehat{p}]=0`$ we easily verify the second line in (3.132).
$`\mathrm{}`$
Remark :
In the form (3.132) the exponential of $`\underset{¯}{\overset{^}{p}}`$ can directly be defined through the spectral theorem without recourse to Nelson’s theorem.
Proof of Theorem 3.4 :
The proof follows trivially from the second line of (3.132).
$`\mathrm{}`$
Inserting formula (3.130) into (3.128) does not directly help us as not all the $`\widehat{g}_{A_kB_k}`$ stand to the right. However, making use again of a finite number of commutation relations and the eigenvalue property we see that we have to leading order in $`t`$
$$<\widehat{p}_{j_1}..\widehat{p}_{j_m}\widehat{h}_{A_1B_1}..\widehat{h}_{A_nB_n}>_g^t=g_{C_1B_1}..g_{C_nB_n}<\widehat{p}_{j_1}..\widehat{p}_{j_m}\left(e^{i\underset{¯}{\overset{^}{p}}/2}\right)_{A_1C_1}..\left(e^{i\underset{¯}{\overset{^}{p}}/2}\right)_{A_nC_n}>_g^t[1+O\left(t\right)]$$
(3.141)
Using the explicit expression (3.132) for $`e^{i\underset{¯}{\overset{^}{p}}/2}`$ and $`[\widehat{p}_j,\widehat{p}]=0`$ we see that we can compute the leading order of (3.141) once we know all expectation values of the form
$$<\widehat{p}_{j_1}..\widehat{p}_{j_m}f\left(\widehat{p}^2\right)>_g^t=<f\left(\widehat{p}^2\right)\widehat{p}_{j_1}..\widehat{p}_{j_m}>_g^t$$
(3.142)
where $`f`$ is an arbitrary analytical function of $`\widehat{p}^2`$. If $`f`$ would be at most a polynomial in $`\widehat{p}^2`$ then it would suffice to know all the matrix elements of the form
$$<\widehat{p}_{j_1}..\widehat{p}_{j_m}>_g^t$$
(3.143)
however, the functions of $`\widehat{p}^2`$ that appear in (3.132) are not simply polynomials. In what follows we will show that (3.142) can be computed once we know $`<f\left(\widehat{p}^2\right)>_g^t`$. The latter can then be computed by an appeal to the solution of the moment problem by Hamburger.
Recall that
$$\widehat{p}_j\psi _g^t\left(h\right)=it/2\left(d/ds\right)_{s=0}\psi _g^t\left(e^{s\tau _j}h\right)=it/2\left(d/ds\right)_{s=0}\psi _{e^{s\tau _j}g}^t\left(h\right)$$
(3.144)
and since $`\widehat{p}_j`$ is self-adjoint
$$<\widehat{p}_{j_1}..\widehat{p}_{j_m}f\left(\widehat{p}^2\right)>_g^t=(\frac{it}{2})^m[\left(\frac{d}{ds_m}\right)_{s_m=0}..\left[\left(\frac{d}{ds_1}\right)_{s_1=0}\frac{<\psi _{e^{s_1\tau _{j_1}}..e^{s_m\tau _{j_m}}g}^t,f\left(\widehat{p}^2\right)\psi _g^t>}{\psi _g^t^2}\right]..]$$
(3.145)
Let $`g^{}=e^{s_1\tau _{j_1}}..e^{s_m\tau _{j_m}}g`$, then since $`\widehat{p}\pi _{jmn}\left(h\right)=t\left(j+1/2\right)\pi _{jmn}\left(h\right)`$ we have with $`f(.)=F\sqrt{(.)}`$
$`<\psi _g^{}^t,F\left(\widehat{p}\right)\psi _g^t>`$ $`=`$ $`{\displaystyle \underset{j}{}}d_je^{t\left(d_j^21\right)/4}F\left(td_j/2\right)\chi _j\left(g^{}\overline{g}^T\right)`$ (3.146)
$`=`$ $`{\displaystyle \frac{1}{2\mathrm{sinh}\left(z\right)}}{\displaystyle \underset{j}{}}d_je^{t\left(d_j^21\right)/4}F\left(td_j/2\right)\left[e^{zd_j}e^{zd_j}\right]`$
$`=`$ $`{\displaystyle \frac{e^{t/4}}{2\mathrm{sinh}\left(z\right)T}}{\displaystyle \underset{n\text{ }\mathrm{Z}}{}}e^{\left(nT\right)^2)/4}F\left(\left(nT\right)T/2\right)e^{\left(nT\right)\frac{z}{T}}`$
where $`\mathrm{cosh}\left(z\right)=\text{tr}\left(g^{}\overline{g}^T\right)/2`$ and $`T=\sqrt{t}`$. Applying the Poisson summation formula to (3.146) we find
$$<\psi _g^{}^t,F\left(\widehat{p}\right)\psi _g^t>=\frac{e^{t/4}2\mathrm{sinh}(zT^2}{}_{n\text{ }\mathrm{Z}}g_n\left(z\right)$$
(3.147)
where
$$g_n\left(z\right)=\left[_{\text{ }\mathrm{R}}𝑑xe^{ikx}g\left(x\right)\right]_{k=2\pi n/T}=_{\text{ }\mathrm{R}}𝑑xe^{x^2/4}F\left(xT/2\right)xe^{x\frac{z2\pi in}{T}}$$
(3.148)
Let $`\mathrm{cosh}\left(z_0\right)=\text{tr}\left(g\overline{g}^T\right)/2`$ and define $`z=z(s_1,..,s_m),z_k=z(s_1,.,s_k,0,..,0),k=0,..,m`$. Let $`G\left(z\right)`$ be any function of $`z`$, then
$`[\left({\displaystyle \frac{d}{ds_m}}\right)_{s_m=0}..\left[\left({\displaystyle \frac{d}{ds_1}}\right)_{s_1=0}G\left(z\right)\right]..]=[\left({\displaystyle \frac{d}{ds_m}}\right)_{s_m=0}..\left[\left({\displaystyle \frac{d}{ds_2}}\right)_{s_2=0}\left({\displaystyle \frac{dz}{ds_1}}\right)_{s_1=0}G^{\left(1\right)}\left(z_{m1}\right)\right]..]`$ (3.149)
$`=`$ $`[\left({\displaystyle \frac{d}{ds_m}}\right)_{s_m=0}..[\left({\displaystyle \frac{d}{ds_3}}\right)_{s_3=0}\{\left({\displaystyle \frac{d^2z}{ds_1ds_2}}\right)_{s_1=s_2=0}G^{\left(1\right)}\left(z_{m2}\right)`$
$`+\left({\displaystyle \frac{dz}{ds_1}}\right)_{s_1=s_2=0}\left({\displaystyle \frac{dz}{ds_2}}\right)_{s_1=s_2=0}G^{\left(2\right)}\left(z_{m2}\right)\}]..]`$
$`=`$ $`\mathrm{}`$
$`=`$ $`\left({\displaystyle \frac{d^mz}{ds_1..ds_m}}\right)_{s_1=..=s_m=0}G^{\left(1\right)}\left(z_0\right)+\mathrm{}+\left({\displaystyle \frac{dz}{ds_1}}\right)_{s_1=..=s_m=0}..\left({\displaystyle \frac{dz}{ds_m}}\right)_{s_1=..=s_m=0}G^{\left(m\right)}\left(z_0\right)`$
Applied to $`G\left(z\right)=<\psi _g^{}^t,F\left(\widehat{p}\right)\psi _g^t>`$ we infer that the derivatives of $`g_n\left(z\right)/\mathrm{sinh}\left(z\right)`$ at $`z_0`$ of all orders $`k`$ between $`1`$ and $`m`$ appear in (3.145) with coefficients that involve sums, over all partitions of $`m=l_1+..+l_k,l_j1`$, of products of $`k`$ factors of the form
$$\left(\frac{d^lz}{ds_{I_1}..ds_{I_l}}\right)_{s_1=..=s_m=0},\mathrm{\hspace{0.33em}1}I_1,..,I_lm\text{ mutually disjoint}$$
(3.150)
Performing the $`k`$th derivative of $`g_n\left(z\right)`$ we obtain
$$g_n^{\left(k\right)}\left(z_0\right)=\frac{1}{T^k}_{\text{ }\mathrm{R}}𝑑xe^{x^2/4}F\left(xT/2\right)x^{k+1}e^{x\frac{z_02\pi in}{T}}$$
(3.151)
The idea is now to do integrations by parts until only $`x`$ appears instead of $`x^{k+1}`$ using $`xe^{x^2/4}=2\left(e^{x^2/4}\right)^{}`$. There are no boundary terms due to the Gaussian. Each time the derivative hits $`e^{x\frac{z2\pi in}{T}}`$ it brings down a factor of $`\frac{z2\pi in}{T}`$ while hitting $`x^lF\left(xT/2\right)`$ produces a non-negative power of $`T`$. We conclude that
$`g_n^{\left(k\right)}\left(z_0\right)`$ $`=`$ $`{\displaystyle \frac{\left(2\left[z_02\pi in\right]\right)^k}{T^{2k}}}\left[{\displaystyle _{\text{ }\mathrm{R}}}𝑑xe^{x^2/4}F\left(xT/2\right)xe^{x\frac{z_02\pi in}{T}}\right]\left(1+O\left(T\right)\right)`$ (3.152)
$`=`$ $`{\displaystyle \frac{\left(2\left[z_02\pi in\right]\right)^k}{T^{2k}}}g_n\left(z_0\right)\left(1+O\left(T\right)\right)`$
Since in (3.145) we multiply with $`T^{2m}`$ and since the derivatives of $`z`$ at $`s_1=..=s_m=0`$ are independent of $`t`$ we see that to leading order in $`t`$ we only need to keep the term with $`k=m`$. The same argument reveals that we do not need to take into account the derivatives of $`1/\mathrm{sinh}\left(z\right)`$.
Next, by a substitution and a contour argument we obtain (at least for functions $`F`$ integrable against the Gaussian)
$$g_n\left(z_0\right)=e^{4\frac{\left(z_02\pi in\right)^2}{t}}_{\text{ }\mathrm{R}}𝑑xe^{x^2/4}F\left(T\left(x+\frac{z_02\pi in}{T}\right)/2\right)\left(x+\frac{z_02\pi in}{T}\right)$$
(3.153)
By our assumption on $`F`$ the integral exists and by the already familiar argument, the $`e^{z_0^2}`$ in (3.153) is controlled by the $`e^{p^2/t}`$ coming from the denominator $`\psi _g^t^2`$ in (3.142) so that (3.153) is exponentially suppressed with $`t0`$ for any $`n0`$.
Putting everything together we therefore have to leading order in $`T`$
$`<\widehat{p}_{j_1}..\widehat{p}_{j_m}f\left(\widehat{p}^2\right)>_g^t`$ $`=`$ $`\left(it/2\right)^m\left(2z_0/t\right)^m\left[{\displaystyle \underset{l=1}{\overset{m}{}}}\left({\displaystyle \frac{dz}{ds_l}}\right)_{s_1=..=s_m=0}\right]{\displaystyle \frac{\frac{e^{t/4}}{2\mathrm{sinh}\left(z_0\right)T^2}_{n\text{ }\mathrm{Z}}g_n\left(z_0\right)\left]\right(1+O\left(T\right))}{\psi _g^t^2}}`$ (3.154)
$`=`$ $`\left[{\displaystyle \underset{l=1}{\overset{m}{}}}\left(iz_0{\displaystyle \frac{dz}{ds_l}}\right)_{s_1=..=s_m=0}\right]<f\left(\widehat{p}^2\right)>_g^t\left(1+O\left(T\right)\right)`$
Now by the method of section 3.1.2 we find
$$z(s_1,..,s_m)=p+i\frac{1}{p}\underset{l=1}{\overset{m}{}}s_lp_{j_l}+O\left(s^2\right)$$
(3.155)
and we obtain the desired result
$$<\widehat{p}_{j_1}..\widehat{p}_{j_m}f\left(\widehat{p}^2\right)>_g^t=p_{j_1}\left(g\right)..p_{j_m}\left(g\right)<f\left(\widehat{p}^2\right)>_g^t(1+O\left(T\right))$$
(3.156)
It therefore remains to compute the expectation value $`<f\left(\widehat{p}^2\right)>_g^t`$ where $`f`$ for the purpose of computing (3.128) can be chosen analytical in $`\widehat{O}:=\widehat{p}^2`$ and at most exponentially growing with $`\widehat{p}`$. Consider first the expectation values of the powers $`\widehat{O}^n`$. Since they are of the form (3.155) with the choice $`f=1`$ we immediately find
$$\underset{t0}{lim}<\widehat{O}^n>_g^t=O\left(g\right)^n=\left(p_j\left(g\right)p_j\left(g\right)\right)^n$$
(3.157)
The assertion $`lim_{t0}<f\left(\widehat{p}^2\right)>_g^t=f\left(p_j\left(g\right)p_j\left(g\right)\right)`$ follows therefore immediately from the solution of the moment problem due to Hamburger which is the subject of the next subsection.
Remark :
Obviously, since we can compute commutators of polynomial operators and express it in terms of elementary operators again, the Ehrenfest theorem to first order in $`t`$ for such operators is trivially satisfied because the operator algebra of elementary operators precisely mirrors the classical Poisson algebra.
### 3.2 Expectation Values of Non-Polynomial Operators and the Moment Problem due to Hamburger
Recall the following theorem (see, e.g. )
###### Theorem 3.5 (Hamburger)
Let be given a sequence of real numbers $`a_n\text{ }\mathrm{R},n=0,1,2,..`$. A necessary and sufficient criterion for the existence of a positive, finite measure $`d\rho (x)`$ on $`\mathrm{R}`$ such that the $`a_n`$ are its moments, that is,
$$a_n=_{\text{ }\mathrm{R}}𝑑\rho \left(x\right)x^n$$
(3.158)
is that for any natural number $`0N<\mathrm{}`$ and arbitrary complex numbers $`z_k,k=0,..,N`$ it holds that
$$\underset{k,l=0}{\overset{N}{}}\overline{z}_kz_la_{n+m}0$$
(3.159)
The measure is faithful if equality in (3.159) occurs only for $`z_k=0`$. Moreover, if there exist constants $`\alpha ,\beta >0`$ such that $`|a_n|\alpha \beta ^n(n!)`$ for all $`n`$, then the measure $`\rho `$ is unique.
Necessity is easy to see by considering the $`L_2`$ norm of the functions $`_{k=0}^Nz_kx^k`$. Sufficiency follows from the spectral theorem and uniqueness can be established by an appeal to Nelson’s analytic vector theorem.
In this section we assume that all operators under consideration are densely defined on a common domain which they together with arbitrary powers leave invariant. We are then able to prove the following theorem.
###### Theorem 3.6
Let $`\widehat{O}`$ be a self-adjoint operator $`\widehat{O}`$ on $`_\gamma `$ for some $`\gamma \mathrm{\Gamma }_0^\omega `$ built from $`\widehat{p}_j^e,\widehat{h}_e,eE(\gamma )`$, that is, $`\widehat{O}=O(\{\widehat{p}_e,\widehat{h}_e\}_{eE(\gamma )}`$). Let $`O(\stackrel{}{g})=O(\{p_e(g_e),h_e(g_e)\}_{eE(\gamma )})`$ be its real valued classical counterpart and suppose that for every $`n\text{ }\mathrm{N}`$
$$\underset{t0}{lim}<\widehat{O}^n>_{\gamma ,\stackrel{}{g}}^t=O\left(\stackrel{}{g}\right)^n$$
(3.160)
Then for any Borel measurable function $`f`$ on $`\mathrm{R}`$ such that $`<f(\widehat{O})^{}f(\widehat{O})>_g^t<\mathrm{}`$ we have
$$\underset{t0}{lim}<f\left(\widehat{O}\right)>_{\gamma ,\stackrel{}{g}}^t=f\left(O\left(\stackrel{}{g}\right)\right)$$
(3.161)
Proof of Theorem 3.6 :
Let $`E\left(x\right),x\text{ }\mathrm{R}`$ be the spectral projections of $`\widehat{O}`$. Then, by assumption and the spectral theorem
$$\underset{t0}{lim}_{\text{ }\mathrm{R}}d<\xi _{\gamma ,\stackrel{}{g}}^t,E\left(x\right)\xi _{\gamma ,\stackrel{}{g}}^t>x^n=O\left(\stackrel{}{g}\right)^n$$
(3.162)
where $`\xi _\stackrel{}{g}^t=\psi _\stackrel{}{g}^t/\psi _\stackrel{}{g}^t`$. Now $`a_n:=O\left(\stackrel{}{g}\right)^n`$ obviously satisfies all the criteria of theorem 3.5 and we conclude that there exists a measure $`d\rho _\stackrel{}{g}\left(x\right)`$ on $`\mathrm{R}`$ such that
$$_{\text{ }\mathrm{R}}𝑑\rho _\stackrel{}{g}\left(x\right)x^n=O\left(\stackrel{}{g}\right)^n$$
(3.163)
The Dirac measure $`d\rho _\stackrel{}{g}\left(x\right)=\delta _{\text{ }\mathrm{R}}(x,O\left(\stackrel{}{g}\right))dx`$ obviously satisfies (3.163) and choosing $`\alpha =1,\beta =\left|O\left(\stackrel{}{g}\right)\right|`$ in theorem 3.5 obviously satisfies the uniqueness part of the criterion. Thus, the Dirac measure is in fact the unique solution to our moment problem. Thus the spectral measure $`d\rho _\stackrel{}{g}^t\left(x\right):=d<\xi _{\gamma ,\stackrel{}{g}}^t,E\left(x\right)\xi _{\gamma ,\stackrel{}{g}}^t>`$ approaches the Dirac $`\delta `$ distribution when evaluated on monomials $`x^n`$. It follows that the support of $`\rho _\stackrel{}{g}^t`$ gets confined to $`\left\{x_0\right\}`$ as $`t0`$ by definiton of the Lebesgue integral.
Now for any function $`f`$ satisfying the assumptions of the theorem the spectral theorem applies and we have
$$<f\left(\widehat{O}\right)>_\stackrel{}{g}^t=_{\text{ }\mathrm{R}}𝑑\rho _\stackrel{}{g}^t\left(x\right)f\left(x\right)$$
(3.164)
Thus, since $`f`$ is in particular measurable, the limit of both sides of (3.164) turns into (3.161).
$`\mathrm{}`$
###### Corollary 3.1
Let $`\widehat{O}_1,..,\widehat{O}_m`$ be self-adjoint, not necessarily commuting, operators such that
$$\underset{t0}{lim}<\underset{k=1}{\overset{m}{}}\widehat{O}_k^{n_k}>_{\gamma ,\stackrel{}{g}}^t=\underset{k=1}{\overset{m}{}}O_k\left(\stackrel{}{g}\right)^{n_k}$$
(3.165)
Then for any Borel measurable function $`f`$ on $`\text{ }\mathrm{R}^m`$
$$\underset{t0}{lim}<f\left(\left\{\widehat{O}_k\right\}_{k=1}^m\right)>_{\gamma ,\stackrel{}{g}}^t=f\left(\left\{O_k\left(\stackrel{}{g}\right)\right\}_{k=1}^m\right\})$$
(3.166)
Proof of Corollary 3.1 :
By the spectral theorem
$$\underset{t0}{lim}_{\text{ }\mathrm{R}^m}d^m<E_1\left(x_1\right)..E\left(x_m\right)>_\stackrel{}{g}^t\underset{k=1}{\overset{m}{}}x_k^{n_k}=\underset{k=1}{\overset{m}{}}O_k\left(\stackrel{}{g}\right)^{n_k}$$
(3.167)
where $`E_k\left(x\right)`$ is the family of spectral projections of $`\widehat{O}_k`$. Thus, by the unique solution to the moment problem the measure in (3.167) approaches the product Dirac measure $`d^mx_k\delta _{\text{ }\mathrm{R}}(x_k,O_k\left(\stackrel{}{g}\right))`$ similar as in theorem 3.6.
$`\mathrm{}`$
Next we turn to commutators.
###### Theorem 3.7
Suppose that $`\widehat{O}_1,\widehat{O}_2`$ are self-adjoint operators satisfying the assumptions of corollary 3.1. Suppose, moreover, that $`\widehat{O}_1`$ is positive semi-definite and that
$$\underset{t0}{lim}\frac{<[\widehat{O}_1,\widehat{O}_2]>_g^t}{it}=\{O_1,O_2\}\left(\stackrel{}{g}\right)$$
(3.168)
Then for any real number $`r`$
$$\underset{t0}{lim}\frac{<[\left(\widehat{O}_1\right)^r,\widehat{O}_2]>_g^t}{it}=\{\left(O_1\right)^r,O_2\}\left(\stackrel{}{g}\right)$$
(3.169)
Proof of Theorem 3.7 :
It suffices to prove the theorem for rational $`r=m/n`$ with $`m,n`$ integers and $`n>0`$. We have the identity
$`{\displaystyle \frac{[\widehat{O}_1^m,\widehat{O}_2]}{it}}`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{m}{}}}\widehat{O}_1^{k1}{\displaystyle \frac{[\widehat{O}_1,\widehat{O}_2]}{it}}\widehat{O}_1^{mk}`$ (3.170)
$`=`$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}\widehat{O}_1^{m\left(k1\right)/n}{\displaystyle \frac{[\widehat{O}_1^r,\widehat{O}_2]}{it}}\widehat{O}_1^{m\left(nk\right)/n}`$
Now for any measurable function $`f`$ we have by assumption and completeness relation
$`\underset{t0}{lim}<f\left(\widehat{O}_1\right)f\left(\widehat{O}_1\right)>_g^t`$ $`=`$ $`\underset{t0}{lim}\left({\displaystyle \frac{2}{\pi t^3}}\right)^N{\displaystyle _{\left(G^{\text{ }\mathrm{C}}\right)^N}}d^N\mathrm{\Omega }\left(\stackrel{}{g}^{}\right)<f\left(\widehat{O}_1\right)>_{\stackrel{}{g}\stackrel{}{g}^{}}^t<f\left(\widehat{O}_1\right)>_{\stackrel{}{g}^{}\stackrel{}{g}}^t`$ (3.171)
$`=`$ $`\underset{t0}{lim}<f\left(\widehat{O}_1\right)>_{\stackrel{}{g}\stackrel{}{g}}^t<f\left(\widehat{O}_1\right)>_{\stackrel{}{g}\stackrel{}{g}}^t`$
meaning that $`\left(\frac{2}{\pi t^3}\right)^N\left|<f\left(\widehat{O}_1\right)>_{\stackrel{}{g}\stackrel{}{g}^{}}^t\right|^2`$ approaches a delta distribution times $`\left(<f\left(\widehat{O}_1\right)>_\stackrel{}{g}^t\right)^2`$, for any $`f`$, with respect to $`\mathrm{\Omega }^N`$ as $`t0`$ where $`N=\left|E\left(\gamma \right)\right|`$. It follows that $`<f\left(\widehat{O}_1\right)>_{gg^{}}^t`$ is concentrated at $`g=g^{}`$ as explicitly displayed in sections 3.1.2, 3.1.3 and we therefore find for the expectation value of (3.170) by using again the completeness relation in a similar fashion
$$m\underset{t0}{lim}<\widehat{O}_1^{m1}>_\stackrel{}{g}^t<\frac{[\widehat{O}_1,\widehat{O}_2]}{it}>_\stackrel{}{g}^t=n\underset{t0}{lim}<\widehat{O}_1^{m\left(n1\right)/n}>_\stackrel{}{g}^t<\frac{[\widehat{O}_1^r,\widehat{O}_2]}{it}>_\stackrel{}{g}^t$$
(3.172)
Using the assumptions of the theorem we thus find
$$\underset{t0}{lim}<\frac{[\widehat{O}_1^r,\widehat{O}_2]}{it}>_g^t=\frac{m}{n}O_1\left(\stackrel{}{g}\right)^{\frac{m}{n}1}\{O_1,O_2\}\left(\stackrel{}{g}\right)=\{O_1^r,O_2\}\left(\stackrel{}{g}\right)$$
(3.173)
as claimed.
$`\mathrm{}`$
The application of these theorems concerns operators which are not polynomials of the elementary ones. Such operators occur in quantum general relativity where diffeomorphism invariance requires that Hamiltonian constraint operators are density one valued and therefore free of UV singularities. This enforces that non-analytic functions, specifically roots of the volume operator , appear. The spectral measure of this operator is not explicitly known and therefore a direct computation of its expectation values and its commutators with holonomy operators that appear in is a hopeless task. Theorems 3.6, 3.7 and corollary 3.1 circumvent this problem at least as far as the leading order behaviour of expectation values is concerned by using the following trick : The fourth power of the volume operator $`\widehat{O}:=\widehat{V}^4`$ is in fact a polynomial of the $`p_j^e`$ and thus the expectation values of $`\widehat{O}^n`$ and the commutators with holonomy operators can be straightforwardly computed, leading to the expected result. Defining then $`\widehat{V}:=\widehat{O}^{1/4}`$ and using the above results shows that the Hamiltonian constraint indeed has the correct classical limit. Details will appear elsewhere .
Acknowledgements
O. W. thanks the Studienstiftung des Deutschen Volkes for financial support.
## Appendix A The $`U(1)`$ case
In this appendix we will apply the results of this paper to the case of $`U\left(1\right)`$ as the gauge group. As will become clear, the much simpler structure of $`U\left(1\right)`$ leads to a considerable simplification of the derivation of all the results. The main reason for this is, of course, the fact that $`U\left(1\right)`$ is Abelian and as a consequence of this that all its irreducible representations are one-dimensional. This means that one has to deal with numbers only, instead of matrices.
### A.1 Expectation Values of the Momentum Operator
We will first show that the expectation value of the momentum operator with respect to the $`U\left(1\right)`$ coherent states has the proper (semi-)classical limit. Recall from the form of the coherent states for $`U\left(1\right)`$:
$$\psi _g^t\left(h\right)=\underset{n}{}e^{\frac{t}{2}n^2}\left(gh^1\right)^n$$
(A.1)
with $`g=e^pe^{i\theta _0}`$ and $`h=e^{i\theta }`$. By the same token as in section 3.1.4, the expectation value of $`\widehat{p}`$ with respect to these states is given by :
$$\frac{\psi _g^t,\widehat{p}\psi _g^t}{\psi _g^t^2}=it\left(\frac{d}{dr}\right)_{r=0}\frac{_ne^{tn^2}\left(e^{ir}e^{2p}\right)^n}{_ne^{tn^2}e^{2np}}=\frac{_ne^{tn^2}nte^{2np}}{_ne^{tn^2}e^{2np}}$$
(A.2)
To see the behaviour of this expression for $`t0`$, we have to perform a Poisson transformation. For the denominator this has already been done in the appendix of , the result being
$$\underset{n}{}e^{tn^2}e^{2np}=\sqrt{\pi }e^{\frac{p^2}{t}}\underset{n}{}e^{\frac{\pi ^2n^2+2i\pi np}{t}}.$$
(A.3)
The transformation for the numerator, however, has to be calculated anew. With $`s=\sqrt{t}`$, as usual, we get
$`\stackrel{~}{f}\left(k\right)`$ $`=`$ $`s{\displaystyle 𝑑xxe^{x^2}e^{\left(\frac{2p}{s}ik\right)x}}`$ (A.4)
$`=`$ $`se^{\frac{\left(\frac{2p}{s}ik\right)}{4}}{\displaystyle 𝑑x\left(x\frac{ik2p/s}{2}\right)e^{x^2}}`$
$`=`$ $`\sqrt{\pi }s{\displaystyle \frac{\left(2p/sik\right)}{2}}e^{\frac{\left(\frac{2p}{s}ik\right)}{4}}.`$
Inserting this into (A.2) yields
$`{\displaystyle \frac{\psi _g^t,\widehat{p}\psi _g^t}{\psi _g^t^2}}`$ $`=`$ $`{\displaystyle \frac{s_n\left(p\pi ni\right)/se^{\frac{\left(2p2\pi ni\right)^2}{4t}}}{e^{p^2/t}_ne^{\frac{\pi ^2n^2+2\pi inp}{t}}}}`$ (A.5)
$`=`$ $`{\displaystyle \frac{p_ne^{\frac{\pi ^2n^2+2\pi inp}{t}}2\pi i_nne^{\frac{\pi ^2n^2+2\pi inp}{t}}}{_ne^{\frac{\pi ^2n^2+2\pi inp}{t}}}}`$
$`=`$ $`p{\displaystyle \frac{2\pi i_nne^{\frac{\pi ^2n^2+2\pi inp}{t}}}{_ne^{\frac{\pi ^2n^2+2\pi inp}{t}}}}.`$
and this result makes it immediately obvious that
$$limt0\frac{\psi _g^t,\widehat{p}\psi _g^t}{\psi _g^t^2}=p$$
(A.6)
Our next task is to generalize this calculation to an arbitrary integer power of the momentum operator. We thus have
$$\frac{\psi _g^t,\widehat{p}^m\psi _g^t}{\psi _g^t^2}=\frac{_ne^{tn^2}\left(nt\right)^me^{2np}}{_ne^{tn^2}e^{2np}}$$
(A.7)
which now has to be Poisson transformed again. The transform for the denominator stays the same, so we can concentrate on the numerator. As the general steps are the same as above we will be more concise here. The Poisson transform is given by
$`\stackrel{~}{f}\left(k\right)`$ $`=`$ $`s^me^{\frac{\left(ik2p/s\right)^2}{4}}{\displaystyle 𝑑ye^{y^2}\left(y+\left(p/sik/2\right)\right)^m}`$ (A.8)
$`=`$ $`s^me^{\frac{\left(ik2p/s\right)^2}{4}}{\displaystyle 𝑑ye^{y^2}\underset{l=0}{\overset{m}{}}y^{ml}\left(p/sik/2\right)^l\left(\genfrac{}{}{0pt}{}{m}{l}\right)}`$
As we are only interested in the limit $`t0`$ and due to the prefactor $`s^m`$ the only surviving term will be the $`l=m`$ term. We therefore have
$`\underset{t0}{lim}\stackrel{~}{f}\left(k\right)`$ $`=`$ $`s^me^{\frac{\left(ik2p/s\right)^2}{4}}{\displaystyle 𝑑ye^{y^2}\left(p/sik/2\right)^m}`$ (A.9)
$`=`$ $`\sqrt{\pi }\left(pi\pi n\right)^me^{\frac{\left(2i\pi n2p\right)^2}{4t}}`$
where we already substituted the $`k`$ variable by $`2\pi n/\sqrt{t}`$. This expression now has to be put back into (A.7) which yields
$`\underset{t0}{lim}{\displaystyle \frac{\psi _g^t,\widehat{p}^m\psi _g^t}{\psi _g^t^2}}`$ $`=`$ $`\underset{t0}{lim}{\displaystyle \frac{\frac{2\pi }{\sqrt{t}}\sqrt{\pi }e^{p^2/t}_n\left(pi\pi n\right)^me^{\frac{\pi ^2n^2+2\pi inp}{t}}}{\frac{2\pi }{\sqrt{t}}\sqrt{\pi }e^{p^2/t}_ne^{\frac{\pi ^2n^2+2\pi inp}{t}}}}`$ (A.10)
$`=`$ $`p^m,`$
as the only term in the sum, which survives in the limit, is the one with $`n=0`$.
Although these results are quite satisfying, one often encounters other powers of the momentum operators, especially square and higher roots, so it would be reassuring to know that they, too, have the expected semiclassical behaviour. A direct calculation as performed above becomes quite difficult for roots of arbitrary polynomials of $`\widehat{p}_e`$ where $`e`$ labels the edges of a graph (for one edge and, say, $`\sqrt{\left|\widehat{p}\right|}`$ the computational effort is still low and is left to the reader as an exercise) so we have to resort to other methods. A clue comes from reformulating the expectation value for integer powers of $`\widehat{p}`$:
$`\underset{t0}{lim}{\displaystyle \frac{\psi _g^t,\widehat{p}^m\psi _g^t}{\psi _g^t^2}}`$ $`=`$ $`\underset{t0}{lim}{\displaystyle \frac{_n\psi _g^t,\widehat{p}^mnn,\psi _g^t}{\psi _g^t^2}}`$ (A.11)
$`=`$ $`\underset{t0}{lim}{\displaystyle \frac{_n\left(nt\right)^m\left|n,\psi _g^t\right|^2}{\psi _g^t^2}}`$
$`=`$ $`p^m`$
where the first two lines are an expansion in terms of $`|n`$, the basis consisting of eigenvectors of $`\widehat{p}`$ \- recall, that $`U\left(1\right)`$ momenta have discrete spectrum - , and the last line follows from our calculations above. This suggests that $`n,\psi _g^t|^2/\psi _g^t^2`$ approaches - in the sense of distributions - just $`\delta _{n,p/t}`$, an observation that receives additional support from the explicit form of $`lim_{t0}\left|n,\psi _g^t\right|^2/\psi _g^t^2`$ that was calculated in . That this also holds in a rigorous sense is guaranteed by the solution to the moment problem by Hamburger as quoted in the main text. In our case the $`a_n`$ are given by the $`p^n`$, therefore $`a_{n+m}=a_na_m`$ and thus the condition of the theorem is obviously satisfied. We can therefore conclude that our results for the integer powers of the momentum operator indeed determine $`\left|n,\psi _g^t\right|^2/\psi _g^t^2`$ to approach $`\delta _{n,p/t}`$. This important result will considerably simplify the calculations in the following subsections. We now come back to the problem of the roots of the momentum operator. Let $`m`$ be an odd integer. Then
$`\underset{t0}{lim}{\displaystyle \frac{\psi _g^t,\widehat{p}^{\frac{m}{2}}\psi _g^t}{\psi _g^t^2}}`$ $`=`$ $`\underset{t0}{lim}{\displaystyle \frac{_n\psi _g^t,\widehat{p}^{\frac{m}{2}}nn,\psi _g^t}{\psi _g^t^2}}`$ (A.12)
$`=`$ $`\underset{t0}{lim}{\displaystyle \frac{_n\left(nt\right)^{\frac{m}{2}}\left|n,\psi _g^t\right|^2}{\psi _g^t^2}}`$
$`=`$ $`p^{\frac{m}{2}}`$
where the last equality is now justified by the aforementioned theorem.
### A.2 Expectation Values of the Holonomy Operator
In this subsection we will compute the semiclassical limit of expectation values of (powers of) the configuration operator $`\widehat{h}`$. We can basically reduce this case to the one in the last subsection by the useful observation that $`\widehat{h}=e^{1/2t}e^{\widehat{p}}\widehat{g}`$, see . For $`m`$ integer or half-integer we get
$`\underset{t0}{lim}{\displaystyle \frac{\psi _g^t,\widehat{h}^m\psi _g^t}{\psi _g^t^2}}`$ $`=`$ $`\underset{t0}{lim}{\displaystyle \frac{e^{m/2t}\psi _g^t,e^{\widehat{p}}\widehat{g}\mathrm{}e^{\widehat{p}}\widehat{g}\psi _g^t}{\psi _g^t^2}}`$ (A.13)
$`=`$ $`\underset{t0}{lim}{\displaystyle \frac{e^{m/2t}\psi _g^t,e^{m\widehat{p}}\widehat{g}^m\psi _g^t}{\psi _g^t^2}}`$
$`=`$ $`\underset{t0}{lim}{\displaystyle \frac{e^{m/2t}\psi _g^t,e^{m\widehat{p}}\psi _g^tg^m}{\psi _g^t^2}}`$
$`=`$ $`\underset{t0}{lim}e^{m/2t}e^{mp}g^m`$
$`=`$ $`h^m`$
where we used in line two that all remaining commutator terms are at least of order $`t`$ and therefore vanish in the limit $`t0`$, and in line three that our coherent states are eigenstates of $`\widehat{g}`$. It should be obvious from this that arbitrary mixed polynomials in $`\widehat{p}`$ and $`\widehat{h}`$ can be treated equivalently, to leading order in $`t`$.
### A.3 Expectation Values of Commutators
In this subsection we intend to obtain the semiclassical limit of expectation values of commutator terms by direct computation. The main example we have in mind here is $`\widehat{h}^1[\sqrt{\widehat{V}},\widehat{h}]`$ which plays an important role in the Hamiltonian constraint operator constructed in . Here $`\widehat{V}`$ denotes the volume operator. As this requires rather tedious calculations due to its structure, requiring at least a graph with three-valent vertices (in the gauge-variant case), we will restrict ourselves to the following case : we would like to check that $`\widehat{h}^1[\sqrt{\widehat{p}},\widehat{h}]/\left(it\right)`$ has the right semiclassical limit, i.e. reproduces $`h^1\{\sqrt{\left|p\right|},h\}`$ which is $`i\text{sgn}\left(p\right)/\left(2\sqrt{\left|p\right|}\right)`$. We start with the observation that
$$\widehat{h}\psi _g^t\left(h\right)=\underset{n}{}e^{n^2t/2}g^nh^{n+1}$$
(A.14)
We then have
$`{\displaystyle \frac{\psi _g^t,\widehat{h}^1[\sqrt{\left|\widehat{p}\right|},\widehat{h}]/\left(it\right)\psi _g^t}{\psi _g^t^2}}`$ (A.15)
$`=`$ $`i{\displaystyle \frac{_ne^{n^2t}e^{2pn}\left(\sqrt{\left|\left(n1\right)t\right|}\sqrt{\left|nt\right|}\right)/t}{_ne^{n^2t}e^{2pn}}}`$
$`=`$ $`i{\displaystyle \frac{\frac{1}{\sqrt{\pi }}_n_{\mathrm{}}^{\mathrm{}}𝑑x\frac{\left(\sqrt{\left|xs\right|}\sqrt{\left|x\right|}\right)}{s^{3/2}}e^{x^2}e^{\left(2p/s2i\pi n/s\right)x}}{e^{\frac{p^2}{s^2}}_ne^{\frac{\pi ^2n^2+2i\pi np}{s^2}}}}`$
$`=`$ $`{\displaystyle \frac{\frac{1}{\sqrt{\pi }}_ne^{\left(2p/s2i\pi n/s\right)^2}_{\mathrm{}}^{\mathrm{}}𝑑x\frac{\left(\sqrt{\left|x+p/si\pi n/ss\right|}\sqrt{\left|x+p/si\pi n/s\right|}\right)}{s^{3/2}}e^{x^2}}{e^{\frac{p^2}{s^2}}_ne^{\frac{\pi ^2n^2+2i\pi np}{s^2}}}}`$
where the last integral can involves the choice of a branch cut for $`n0`$. Since the integral certainly converges for any $`n`$ and is multiplied by the exponentially fast vanishing function $`e^{4\pi ^2n^2/s}`$, by the argument already familiar from for the limit $`t0`$ it will be sufficient to keep the term $`n=0`$ for what follows. We thus obtain for the expectation value
$$\underset{t0}{lim}\frac{\psi _g^t,\widehat{h}^1[\sqrt{\left|\widehat{p}\right|},\widehat{h}]/t\psi _g^t}{\psi _g^t^2}=i\frac{1}{\sqrt{\pi }}_{\mathrm{}}^{\mathrm{}}𝑑x\frac{\left(\sqrt{\left|x+p/ss\right|}\sqrt{\left|x+p/s\right|}\right)}{s^{3/2}}e^{x^2}$$
(A.16)
As we are ultimately only interested in the limit $`t0`$, and therefore $`s0`$, we aim at putting the integrand into a form that allows taking the $`s0`$ limit inside:
$`{\displaystyle \frac{i}{\sqrt{\pi }}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑x{\displaystyle \frac{\left(\sqrt{\left|x+p/ss\right|}\sqrt{\left|x+p/s\right|}\right)}{s^{3/2}}}e^{x^2}`$ (A.17)
$`=`$ $`{\displaystyle \frac{i}{\sqrt{\pi }}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑x{\displaystyle \frac{\left(\sqrt{\left|xs+ps^2\right|}\sqrt{\left|xs+p\right|}\right)}{s^2}}e^{x^2}`$
$`=`$ $`{\displaystyle \frac{i}{\sqrt{\pi }}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑x{\displaystyle \frac{\left(xs+ps^2\right)^2\left(xs+p\right)^2}{s^2\left(\left|xs+ps^2\right|^{1/2}+\left|xs+p\right|^{1/2}\right)\left(\left|xs+ps^2\right|+\left|xs+p\right|\right)}}e^{x^2}`$
$`=`$ $`{\displaystyle \frac{i}{\sqrt{\pi }}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑x{\displaystyle \frac{s^22\left(xs+p\right)}{\left(\left|xs+ps^2\right|^{1/2}+\left|xs+p\right|^{1/2}\right)\left(\left|xs+ps^2\right|+\left|xs+p\right|\right)}}e^{x^2}`$
It is easy to see that the limit $`f_p\left(x\right)`$ as $`s0`$ of the integrand $`f_p^s\left(x\right)`$ exists pointwise. Furthermore it is clear that the modulus of the integrand is $`L^1`$-integrable. Next, we write the integrand of (A.17) as $`e^{x^2}f_p^s\left(x\right)=e^{x^2/2}g_p^s\left(x\right)`$ and we seek to give a bound on $`g_p^s\left(x\right)`$ independent of $`s,x`$ for $`s`$ smaller than some $`s_0`$. To that end we estimate $`e^{x^2/2}1`$ for $`\left|x\right|1`$ and $`e^{x^2/}e^{\left|x\right|/2}`$ for $`\left|x\right|1`$ when estimating $`g_p^s\left(x\right)`$. Consider first the region $`\left|x\right|1`$. The first derivative of the estimated $`\left|g_p^s\left(x\right)\right|`$ then leads to a quadratic equation whose roots depend on the signs of both $`p,x`$. The local maxima turn out to lie at $`x=\pm 1`$ and $`xp/s`$. Only the former one is an absolute maximum. The value of $`\left|g_p^s\left(x\right)\right|`$ can then be estimated roughly by $`1/\sqrt{\left|p\right|}`$ up to a multiplicative, numerical constant and the same is true for the region $`\left|x\right|1`$. Altogether we have found, up to a numerical factor the following $`L_1`$ function, independent of $`s`$ that dominates $`f_p^s`$
$$\left|f_p^s\left(x\right)\right|\stackrel{<}{}\frac{e^{x^2/2}}{\sqrt{\left|p\right|}}$$
(A.18)
so that all conditions of the dominated convergence theorem are satisfied, and the $`s0`$ limit can be taken inside the integral. We thus obtain
$`\underset{t0}{lim}{\displaystyle \frac{\psi _g^t,\widehat{h}^1[\sqrt{\left|\widehat{p}\right|},\widehat{h}]/t\psi _g^t}{\psi _g^t^2}}`$ (A.19)
$`=`$ $`{\displaystyle \frac{i}{\sqrt{\pi }}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑x\underset{s0}{lim}{\displaystyle \frac{s^22\left(xs+p\right)}{\left(\left|xs+ps^2\right|^{1/2}+\left|xs+p\right|^{1/2}\right)\left(\left|xs+p^2s^2\right|+\left|xs+p\right|\right)}}e^{x^2}`$
$`=`$ $`{\displaystyle \frac{i}{\sqrt{\pi }}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑x\left({\displaystyle \frac{1}{2\sqrt{\left|p\right|}}}\text{sgn}\left(p\right)\right)e^{x^2}`$
$`=`$ $`{\displaystyle \frac{1}{2\sqrt{\left|p\right|}}}\text{sgn}\left(p\right).`$
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# The influence of chiral surface states on the London penetration depth in Sr2RuO4
Several years of intense experimental research have established the unconventional nature of superconductivity in Sr<sub>2</sub>RuO<sub>4</sub> .$`^{\text{?}\text{}\text{?}\text{)}}`$ This compound has a layered perovskite structure representing a basically two-dimensional metal with three almost cylindrical Fermi surfaces. The symmetry of the superconducting state is very likely odd in parity, which implies the spin-triplet configuration analogous to superfluid <sup>3</sup>He .$`^{\text{?}\text{}\text{?}\text{)}}`$ Muon spin rotation experiments provide evidence for broken time reversal symmetry ,$`^{\text{?}\text{)}}`$ a fact that strongly suggests that the gap function has the basic form,
$$𝐝(𝐤)=\mathrm{\Delta }_0\widehat{𝐳}\frac{k_x\pm ik_y}{k_F}$$
(1)
which is a chiral $`p`$-wave state, here written in the vector representation, assuming cylindrical symmetry. The Cooper pairs possess an internal orbital angular momentum which is oriented along the $`z`$-axis. A consequence of this topological property of the superconducting phase is the presence of chiral surface states at the surface .$`^{\text{?}\text{}\text{?}\text{}\text{?}\text{)}}`$ While the chiral $`p`$-wave state has a basically gapful quasiparticle spectrum, the surface states correspond to subgap quasiparticle excitations with a continuous spectrum down to zero energy .$`^{\text{?, ?)}}`$ These quasiparticle states are Andreev bound states and extend only over a coherence length towards the bulk. In this letter we consider the contribution of these states to the temperature dependence of the London penetration depth. The London penetration depth $`\lambda _{}`$ for currents within the plane and the in-plane coherence length $`\xi _{}`$ are very similar giving a Ginzburg-Landau parameter $`\kappa =\lambda _{}/\xi _{}2.6`$. Therefore, the presence of the surface states can lead to a visible reduction of the screening effect and could even dominate the low-temperature behavior $`\lambda _{}`$, in particular, if the bulk quasiparticle spectrum is gapped. We show here that power-law temperature dependence can result from the surface states, which is usually taken as an evidence for nodes in the bulk quasiparticle gap.
The discussion of the London penetration depth requires a careful analysis of the current-current response to an external field which can be written in general as
$$j_\mu (𝐫,t)=\frac{c}{4\pi }\underset{\nu }{}𝑑t^{}d^3r^{}K_{\mu \nu }(𝐫,t;𝐫^{},t^{})A_\nu (𝐫^{},t^{})$$
(2)
where only the transverse component of $`A_\nu (𝐫^{},t^{})`$ enters. The kernel $`K_{\mu \nu }(𝐫,t;𝐫^{},t^{})`$ is obtained from the current-current correlation function. In our case this response consists of two contributions: the bulk part due to the continuum of quasiparticle states above the gap and the part due to the surface states. We consider from now on the specific case of a surface with normal vector along the $`x`$-axis and an external field parallel to the $`z`$-axis. Consequently we have to deal with the transverse vector potential and screening current along the $`y`$-axis. The relevant terms are then,
$$j_y(𝐫,t)=\frac{1}{c}d^4r^{}\left[\mathrm{\Pi }_{yy}(r;r^{})\frac{c^2\delta ^{(4)}(rr^{})}{4\pi \lambda _0^2}\right]A_y(r^{})$$
(3)
where the integral runs of the four coordinates $`r^{}=(𝐫^{},t^{})`$ and $`\lambda _0`$ corresponds to the bare “London penetration depth” of the bulk regime which can be considered as basically temperature-independent for very low temperatures. For this bulk part we take the local approximation, while for the first part connected with the surface states nonlocality is important, as we will see below.
The surface states can be easily described within the Bogolyubov-de Gennes formalism if we neglect self-consistency of the gap $`\mathrm{\Delta }_0`$ which we choose to be constant everywhere inside the superconductor. For the sake of simplicity we assume also that the surface provides specular reflection of quasiparticles and the gap has no anisotropy on the Fermi surface. The electron band in our model has cylindrical symmetry and is represented by the parabolic form, $`\epsilon _𝐤=\{(k_x^2+k_y^2)k_\mathrm{F}^2\}/2m`$ neglecting any dispersion along the $`z`$-axis. Using these simplifications the wave function of the subgap states localized at the surface is given by
$$\left(\begin{array}{c}u_𝐤\left(𝐫\right)\\ v_𝐤\left(𝐫\right)\end{array}\right)\sqrt{\frac{2}{\xi _0L_yd}}\mathrm{e}^{ik_yy\frac{x}{\xi _0}}\mathrm{sin}(k_xx)\left(\begin{array}{c}1\\ i\end{array}\right).$$
(4)
Since only states very close to Fermi surface are important for the low-temperature properties the wave vector can be represented essentially as $`(k_x,k_y)=k_F(\mathrm{cos}\theta ,\mathrm{sin}\theta )`$ for $`k_F\xi _01`$. Further, $`L_y`$ is the extension of the system along the $`y`$-direction with periodic boundary conditions and $`d`$ is the interlayer spacing (the wave function is renormalized per layer). The energy of the surface states is given by $`E_{k_y}=\eta \mathrm{\Delta }_0k_y/k_F`$ with $`\eta =\pm 1`$ denoting the sign of the chirality of $`𝐝(𝐤)=\widehat{𝐳}(k_x\pm ik_y)/k_F`$ .$`^{\text{?}\text{)}}`$ Defining the current operators as $`j_\mu =(\mathrm{}e/2mi)(\widehat{\mathrm{\Psi }}^{}_\mu \widehat{\mathrm{\Psi }}h.c.)`$, where
$$\widehat{\mathrm{\Psi }}(𝐫)=\left(\begin{array}{c}\psi _{}(𝐫)\\ \psi _{}^{}(𝐫)\end{array}\right)=\underset{𝐤}{}\left[\begin{array}{cc}u_𝐤(𝐫)& v_𝐤^{}(𝐫)\\ v_𝐤(𝐫)& u_𝐤^{}(𝐫)\end{array}\right]\left(\begin{array}{c}\gamma _𝐤\\ \gamma _𝐤^{}\end{array}\right),$$
(5)
is the Nambu field operator with $`\gamma _{𝐤\sigma }^{()}`$ the Bogolyubov quasi-particle operator, we can express the current-current correlation function $`\mathrm{\Pi }_{yy}(r;r^{})`$ as
$`\mathrm{\Pi }_{yy}(r;r^{})`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2e^2}{4m^2}}\underset{r_1,r_2r,r_1^{},r_2^{}r^{}}{lim}`$ (6)
$`\times (_{y_1}_{y_2})(_{y_2^{}}_{y_1^{}})\mathrm{Tr}\left[𝐆(r_1;r_1^{})𝐆(r_2^{};r_2)\right]`$
where $`𝐆(r;r^{})`$ is the Nambu-Gor’kov Green’s function in real space and $`_y`$ denotes the derivative with respect to the spatial $`y`$-coordinate. The Green’s function can be expressed as
$$𝐆(𝐫,𝐫^{};i\omega _n)\frac{\varphi (x)\varphi (x^{})}{L_yd}\underset{0<k<k_F}{}\underset{s=\pm }{}\frac{\widehat{\sigma }_0s\widehat{\sigma }_2}{i\omega _nsE_k}\mathrm{e}^{isk(yy^{})}$$
(7)
with $`\varphi (x)=\sqrt{2/\xi _0}\mathrm{exp}^{x/\xi _0}\mathrm{sin}k_Fx`$, $`\widehat{\sigma }_0`$ the unit matrix and $`\widehat{\sigma }_2`$ the second Pauli matrix, and $`\omega _n`$ being the fermionic Matsubara frequency. Using Eqs. (6) and (7), we calculate the current-current correlation function. The translational invariance along the $`y`$-direction allows us to transform the $`y`$-coordinate into momentum space,
$`\mathrm{\Pi }_{yy}(x,x^{};q,i\mathrm{\Omega }_n)`$ (8)
$`{\displaystyle \frac{32\mathrm{}^2e^2}{m^2L_yd}}g(x)g(x^{}){\displaystyle \underset{0<k<k_F}{}}k^2{\displaystyle \frac{f(E_{k+q})f(E_k)}{i\mathrm{\Omega }_nE_{k+q}+E_k}}`$
$`{\displaystyle \frac{8\pi \mathrm{}^2k_F^3}{3dm^2\mathrm{\Delta }_0}}\left({\displaystyle \frac{k_BT}{\mathrm{\Delta }_0}}\right)^2{\displaystyle \frac{g(x)g(x^{})}{1i\mathrm{\Omega }_nk_F/q\mathrm{\Delta }_0}},`$
in the limit $`TT_c`$, where $`g(x)\mathrm{exp}(\frac{2x}{\xi _0})\mathrm{sin}^2(k_Fx)/\xi _0`$ is the square of the amplitude of the surface state wave function and $`i\mathrm{\Omega }_n`$ is the bosonic Matsubara frequency. In deriving Eq. (8), we restrict ourselves to the leading contribution for $`qk_F`$ and $`k_BT\mathrm{\Delta }_0`$. The nonlocal nature of response enters via the product form $`g(x)g(x^{})`$ which accounts for the fact that each of the quasiparticle state is localized at the surface. The field at the particular point $`(x^{},y^{},z^{})`$ couples to the surface state in the same layer with a weight $`g(x^{})`$ and yields consequently a response at any other point $`(x,y,z^{})`$ with weight $`g(x)`$. This is a feature of the effectively one-dimensional character of the surface states within each layer. A local approach in this place would underestimate the role of the surface states in the low-temperature response.
Combining Eqs. (3) and (8), where we further use the analytic continuation $`i\mathrm{\Omega }_n\mathrm{}\omega +i\delta `$, with the Maxwell equation $`^2A_y(𝐫,\omega )=\frac{4\pi }{c}j_y(𝐫,\omega )`$ we obtain an integro-differential equation for $`A_y(𝐫,t)`$. The boundary condition is given by $`_xA_y(𝐫,\omega )|_{x=0}=B_z(q,\omega )`$, where $`B_z(q,\omega )`$ is the external magnetic field at the surface parallel to the $`z`$-axis. We solve this equation using an approximation $`g(x)\mathrm{exp}(2x/\xi _0)/2\xi _0`$ in the integrand (we ignore the fast oscillations), which is certainly valid for $`k_F\xi _01`$. This allows us to calculate the surface impedance, $`Z(q,\omega )=4\pi cE_y(x=0,q,\omega )/B(q,\omega )=4\pi i\omega A_y(x=0,q,\omega )/B(q,\omega )`$. We then obtain the penetration depth using the relation $`4\pi \omega \lambda (q,\omega )=\mathrm{Im}Z(q,\omega )`$. Taking a static limit $`\mathrm{}\omega k_F/q\mathrm{\Delta }_00`$ and $`q0`$, we find in the regime of $`k_BT\mathrm{\Delta }_0`$, $`\mathrm{\Delta }\lambda (T)=\lambda (T)\lambda (0)`$ has $`T^2`$-behavior:
$$\mathrm{\Delta }\lambda (T)/\lambda _0\frac{4\pi ^2}{3}\frac{\kappa }{(2\kappa +1)^2}\left(k_BT/\mathrm{\Delta }_0\right)^2.$$
(9)
Setting $`\kappa 2.6`$, which is a typical value of Sr<sub>2</sub>RuO<sub>4</sub>, we obtain $`\mathrm{\Delta }\lambda (T)/\lambda _00.14\times \left(T/T_c\right)^2`$, if we assume the weak-coupling relation $`\mathrm{\Delta }_0=1.76k_BT_c`$. We ignored the temperature dependence of $`\lambda _0`$ as it is exponential in the low-temperature regime in our model.
It is important to notice that this contribution is independent of the sign of the chirality and the charge of the carriers, i.e., whether the superconducting state is $`𝐝\widehat{𝐳}(k_x+ik_y)`$ or $`\widehat{𝐳}(k_xik_y)`$) and the Fermi surface is electron- or hole-like. Therefore, the formation of domains of the two superconducting states would not lead to a significant change of the result. Furthermore, for the case of several superconducting bands the contributions of the surface states of each band add up to enlarges the prefactor of the $`T^2`$-law. The coherence length as the extension of the surface states towards the interior is different for each band, since the Fermi velocities and the gap magnitudes are different. The coherence length, experimentally determined via the measurement of $`H_{c2}`$, giving $`\kappa 2.6`$ is the shortest among all. Therefore, the enhancement of the surface state contribution can be sizable in the multiband case. We consider here the case of three bands, as schematically shown in Fig.1, where each band yields its own surface state described by a Green’s function $`𝐆^{(j)}(r;r^{})`$, (the superscript $`j`$ labels the $`j^{th}`$band). We assume that the reflection of quasiparticles on the surface does not lead to transitions among the different bands. Each band is characterized by a Fermi vector $`k_F^{(j)}`$, the effective band mass $`m_j`$ and the superconducting gap $`\mathrm{\Delta }_j`$. We now use Eq.(6) to calculate the contribution of each band to the current-current correlation function and then analyze the resulting equation for the transverse vector potential as in the single band case. This leads to the low-temperature behavior of the London penetration depth,
$$\frac{\mathrm{\Delta }\lambda (T)}{\lambda _0}\frac{4\pi ^2}{3}\underset{j=1}{\overset{3}{}}\frac{\kappa ^{(j)}}{(2\kappa ^{(j)}+1)^2}\left(\frac{\lambda _0\mathrm{\Delta }_0}{\lambda ^{(j)}\mathrm{\Delta }_j}\right)^2\left(\frac{k_BT}{\mathrm{\Delta }_0}\right)^2$$
(10)
where
$$\frac{1}{\lambda _0^2}=\underset{j}{}\frac{1}{\lambda ^{(j)2}}=\underset{j}{}\frac{2e^2k_F^{(j)2}}{m_jc^2d}$$
(11)
and $`\kappa ^{(j)}=\lambda _0/\xi _0^{(j)}`$ with $`\xi _0^{(j)}=\mathrm{}^2k_F^{(j)}/m_j\mathrm{\Delta }_j`$. We choose $`\mathrm{\Delta }_0`$ to reproduce again the weak coupling relation, $`\mathrm{\Delta }_0=1.76k_BT_c`$. The contribution of all three bands can easily give a prefactor to the $`(T/T_c)^2`$-law of order one, consistent with recent measurements for fields along the $`z`$-axis.
In the discussion of the multi-band situation we neglected the interband effects. Oscillatory fields, for example appearing in microwave experiments, can yield interband transitions. The matrix elements for the transition depends on various details of the orbital and band structure. We do not go into these complex details here, but assume that the interband transition can be described by an ordinary current operator, $`j_\mu ^{(i,j)}=(\mathrm{}e/2m^{}i)(\widehat{\mathrm{\Psi }}^{(i)}_\mu \widehat{\mathrm{\Psi }}^{(j)}h.c.)`$ where $`m^{}`$ is a phenomenological parameter accounting for the matrix element. $`\widehat{\mathrm{\Psi }}^{(i)}`$ is the Nambu field operator of the $`i`$-th band as in Eq. (5). In the small-momentum $`(q)`$ and small-frequency $`(\omega )`$ limit only surface states are important which share the same zero-energy momentum. As shown in Fig.1 this is the case for the $`\beta `$\- and $`\gamma `$-band. The $`\alpha `$-band is unimportant because for interband transitions a large momentum transfer ($`q\pi `$) is necessary. Analogous to Eq.(8) we can derive the correlation function,
$$\begin{array}{c}\mathrm{\Pi }_{yy}^{(\beta ,\gamma )}(x_1,x_2,q;i\mathrm{\Omega }_n)\frac{16\mathrm{}^2e^2}{m^2dL_y}\stackrel{~}{g}(x_1)\stackrel{~}{g}(x_2)\underset{k}{}k^2\hfill \\ \\ \times \left[\frac{f(E_{k+q}^\beta )f(E_k^\gamma )}{i\mathrm{\Omega }_nE_{k+q}^\beta +E_k^\gamma }+\frac{f(E_{k+q}^\gamma )f(E_k^\beta )}{i\mathrm{\Omega }_nE_{k+q}^\gamma +E_k^\beta }\right]\hfill \end{array}$$
(12)
where $`\stackrel{~}{g}(x)\mathrm{exp}(2x/\stackrel{~}{\xi })/2\sqrt{\xi ^{(\beta )}\xi ^{(\gamma )}}`$ with $`\stackrel{~}{\xi }^1=\xi ^{(\beta )1}+\xi ^{(\gamma )1}`$. The surface state spectra in the two bands is approximated by $`E_k^{(i)}=v_ik`$, with $`v_i=\mathrm{\Delta }_i/k_F^{(i)}`$. We now take the analytic continuation $`i\mathrm{\Omega }_n\mathrm{}\omega +i\delta `$ and set $`q=0`$ in Eq.(12). For the limit of very small $`\mathrm{}\omega `$ ($`k_BT`$) the surface impedance gets the following contribution from the interband transitions,
$$\begin{array}{c}ImZ_{ib}(\omega )=\frac{4\pi ^3\stackrel{~}{\xi }\xi ^{(\gamma )}}{3\lambda ^{(\gamma )2}}\frac{\stackrel{~}{\kappa }^3}{(2\stackrel{~}{\kappa }+1)^2}\left(\frac{\mathrm{\Delta }_0m_\gamma }{\mathrm{\Delta }_\gamma m^{}}\right)^2\frac{\omega v_\gamma }{\stackrel{~}{v}}\hfill \\ \\ \times \left[\left(\frac{k_BT}{\mathrm{\Delta }_0}\right)^2\left(\frac{v_\gamma ^2}{v_\beta ^2}1\right)+\frac{2}{\pi }\left(\frac{\mathrm{}\omega }{\mathrm{\Delta }_0}\right)^2\mathrm{ln}\frac{v_\gamma }{v_\beta }\right]\hfill \end{array}$$
(13)
for the imaginary part which yields in the zero frequency-limit also a $`T^2`$-contribution to the London penetration depth ($`\stackrel{~}{\kappa }=\lambda /\stackrel{~}{\xi }`$ and $`\stackrel{~}{v}=v_\gamma v_\beta >0`$). For the real part we obtain,
$$\begin{array}{c}ReZ_{ib}(\omega )=\frac{2\pi ^3\stackrel{~}{\xi }\xi ^{(\gamma )}}{\lambda ^{(\gamma )2}}\frac{\stackrel{~}{\kappa }^3}{(2\stackrel{~}{\kappa }+1)^2}\left(\frac{\mathrm{\Delta }_0m_\gamma }{\mathrm{\Delta }_\gamma m^{}}\right)^2\hfill \\ \\ \times \left(\frac{v_\gamma }{\stackrel{~}{v}}\right)^3\omega \left(\frac{\mathrm{}\omega }{\mathrm{\Delta }_0}\right)^2\frac{\mathrm{}\omega }{k_BT}\hfill \end{array}$$
(14)
for given $`T`$ and small $`\mathrm{}\omega `$ ($`k_BT\mathrm{\Delta }_0`$). The imaginary part, the inductive resistance, shows a $`\omega `$-linear plus $`\omega ^3`$-behavior, while the real part, the surface resistance follows an $`\omega ^4`$-law. In the opposite limit where $`k_BT\mathrm{}\omega \mathrm{\Delta }_0`$ the surface impedance due to interband transitions has to vanish. The reason is that for $`q=0`$ the initial and final states are either both empty or occupied in the zero-temperature limit. A simple analysis shows the following behavior,
$$\begin{array}{c}ReZ_{ib}(\omega )\frac{\omega ^4}{T}\mathrm{e}^{\frac{v_\gamma +v_\beta }{2\stackrel{~}{v}}\frac{\mathrm{}\omega }{k_BT}},ImZ_{ib}(\omega )\frac{T^4}{\omega ^2}\hfill \end{array}$$
(15)
The surface resistance and the inductive resistance vanish exponentially and with a power-law, respectively, in the zero-temperature limit.
We now consider the possibility of a so-called nonlinear Meissner effect.The application of a magnetic field introduces a Doppler shift which changes the quasiparticle energy,
$$E_𝐤^{}=E_𝐤^{(sg)}+\frac{ev_{Fy}}{c}A_y$$
(16)
if we again consider the case of $`𝐧=(1,0,0)`$. We can expand the current-current correlation function for small $`A_y`$ and analyze the contribution to the London penetration depth in the same way as done above. Restricting to the single band model we obtain,
$$\begin{array}{c}\frac{\mathrm{\Delta }\lambda (T)}{\lambda _0}\frac{4\pi ^2}{3}\frac{\kappa }{(2\kappa +1)^2}\left(\frac{k_BT}{\mathrm{\Delta }_0}\right)^2\left(1\eta \frac{3\kappa }{2}\frac{H}{H_{c2}}\right)\hfill \end{array}$$
(17)
which leads to a non-linear correction in the external field. This effect is a consequence of the angular momentum of the Cooper pairs coupling to the field along the $`z`$-axis. The sign of this correction depends on the chirality (direction of angular momentum) and the character of the Fermi surface, i.e., a different sign appears for the $`\alpha `$-band than for the $`\beta `$\- and $`\gamma `$-band in Fig.1. Therefore, the presence of electron-like and hole-like Fermi surfaces as well as the formation of domains of the two chiral states lead to compensations which diminish the effect.
In recent experiments by Bonalde et al., the temperature dependence of the London penetration depth was determined using a self-inductive technique.$`^{\text{?}\text{)}}`$ It was found that the low-temperature behavior is indeed governed by a $`T^2`$-behavior. This is not compatible with a simple interpretation in terms of line nodes in the gap as originally proposed based on the $`T`$-power laws in specific heat and NQR ,$`^{\text{?}\text{)}}`$ since this would lead to a linear $`T`$-dependence .$`^{\text{?}\text{)}}`$ The more sophisticated approach based on a nonlocal response theory by Kostin and Leggett (KL) (for the reason that $`\kappa `$ is small), however, would yield a $`T^2`$-behavior .$`^{\text{?}\text{)}}`$ On the other hand, in this letter we propose an alternative mechanism for a $`T^2`$-behavior based on the contributions of the surface states. In both theories it is expected that this power-law behavior is absent for $`\lambda _{}`$, for screening currents flowing along the $`z`$-axis. The $`z`$-axis current is not proportional to the surface state energy as required to obtain the $`T^2`$-behavior in our theory. Furthermore, $`\kappa _{}`$ is about 20 times larger than the in-plane $`\kappa `$ so that contribution of the surface states as well as the nonlocal effect by KL are rather small. However, the measurements for fields in the plane, probing the $`z`$-axis current show a similar $`T^2`$-behavior. This is in apparent conflict with both interpretations. Since in this case, however, not only $`\lambda _{}`$ but also the contribution from in-plane currents from the surfaces normal to the $`z`$-axis are involved, the final answer will be given only when these geometrical aspects have been thoroughly investigated .$`^{\text{?}\text{)}}`$
In our model the gap size is isotropic on the Fermi surface. Anisotropy in turn modifies the surface state spectrum to $`E_{k_y}=vk_y+v^{}k_y^3+\mathrm{}`$ without destroying the particle-hole symmetry ($`vv^{}<0`$ in general). This yields an additional $`T^4`$-contribution which may not be so small. Together with other correction in this order this leads to
$$\frac{\mathrm{\Delta }\lambda (T)}{\lambda _0}=a(T/T_c)^2+b(T/T_c)^4+\mathrm{}$$
(18)
with $`ab>0`$. In an intermediate range the second term generates a T-dependence which over some temperature range appears to be close to a $`T^3`$-behavior and only at rather low temperatures the $`T^2`$-law would dominate. For one sample Bonalde et al. could indeed fit their data reasonably well with a $`T^3`$-curve .$`^{\text{?)}}`$ While this sample happened to be dirty it is not clear from the experiment what is the intrinsic origin for the apparently different power-law. Therefore, the additional $`T^4`$-contribution, which due to the surface orientation and disorder is larger than usual, may be one possible explanation.
A further aspect noteworthy here is that our mechanism is active at the surface only, while the KL scheme also applies also in the bulk of the superconductor. Therefore, the London penetration depth governing the magnetic interaction between vortices should have different temperature dependence in the two scenarios. Measurements of in the mixed phase by $`\mu `$SR suggest that London penetration depth saturates faster than $`T^2`$ at low temperature .$`^{\text{?}\text{)}}`$ In contrast to the surface-sensitive experiment mentioned above ,$`^{\text{?)}}`$ $`\mu `$SR is indeed a bulk probe. Unfortunately, it has less accuracy in determining the temperature dependence of the London penetration depth so that we cannot draw a strong conclusion to date. Nevertheless, the present experimental result is consistent with the interpretation based on surface states.
We would like to thank A. Furusaki, M. Matsumoto, T.M. Rice, C. Honerkamp, J. Goryo, Y. Maeno, D. Van Harlingen and I. Bonalde for many stimulating discussions. This work was supported by a Grant-in-Aid of the Japanese Ministry of Education, Science, Culture and Sports.
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# 1 Introduction
## 1 Introduction
In this paper we continue the investigation started in of the $`SL(2,R)`$ WZW model describing string theory on $`AdS_3\times `$. For other work on this model, see . Our motivation is to understand string theories in curved spacetimes where the metric component $`g_{00}`$ is non-trivial, of which $`AdS_3`$ is the simplest example. Moreover, it is possible to construct black hole solutions as quotients of $`AdS_3`$ , so understanding string theory on $`AdS_3`$ would lead to an understanding of strings moving near black hole horizons.
In the spectrum of $`SL(2,R)`$ WZW model was studied, using spectral flow to generate new representations from the standard ones. These new representations include states corresponding to long strings , with a continuous energy spectrum, as well as discrete states. The existence of spectral flow as a symmetry of the theory was argued on the basis of classical and semi-classical analysis. Further support was given by the fact that the seemingly arbitrary upper bound on the mass of string states in $`AdS_3`$ was removed, thus recovering the infinite tower of masses one expects from string theory. We would like to verify these results by an explicit calculation of the one-loop partition function. As shown in , the Euclidean black hole background is equivalent to the thermal $`AdS_3`$ background. So we will consider string theory on $`AdS_3`$ at a finite temperature, which is described by strings moving on a Euclidean $`AdS_3`$ background with the Euclidean time identified. The calculation of the partition function for this geometry is a minor variation on the calculation of Gawedzki in . From this we can read off the spectrum of the theory in Lorentzian signature by interpreting the result as the free energy of a gas of strings.
This paper is organized as follows. In section 2 we review the spectrum found in . In section 3 we compute the one-loop partition function on thermal $`AdS_3`$. In section 4 we read off the spectrum from the one-loop calculation. First we present a qualitative analysis, which is then followed by a precise calculation. We explain how the different parts of the spectrum arise from this calculation. We further show how the one-loop result contains information about the $`SL(2,R)`$ and Liouville reflection amplitudes.
## 2 The spectrum
We begin by briefly summarizing the results of , where a concrete proposal for the spectrum of $`AdS_3`$ string theory was made. We consider a critical bosonic string theory on $`AdS_3\times `$. The Hilbert space of the $`SL(2,R)`$ WZW model is generated by the action of the left-moving and right-moving current algebra $`\widehat{SL(2,R)}_L\times \widehat{SL(2,R)}_R`$, and all the states form representations of this algebra. The simplest representations are built by first choosing representations for the zero modes, then regarding them as the primary states annihilated by $`J_{n>0}^{3,\pm }`$. The raising operators $`J_{n<0}^{3,\pm }`$ are then used to generate the representations of the current algebra. From harmonic analysis, i.e. quantum mechanical limit, it is known that the left-right symmetric combinations $`𝒞_{j=1/2+is}^\alpha \times 𝒞_{j=1/2+is}^\alpha `$ and $`𝒟_{j>1/2}^\pm \times 𝒟_{j>1/2}^\pm `$ form a complete basis in $`^2(AdS_3)`$, where $`𝒞_{j=1/2+is}^\alpha `$ is the principal continuous representation and $`𝒟_{j>1/2}^\pm `$ the principal discrete representation of $`SL(2,R)`$. These representations are unitary, but the resulting current algebra representations $`\widehat{𝒞}_{j=1/2+is}^\alpha \times \widehat{𝒞}_{j=1/2+is}^\alpha `$ and $`\widehat{𝒟}_{j>1/2}^\pm \times \widehat{𝒟}_{j>1/2}^\pm `$, constructed as explained above, are not. This is not a surprise, for even in flat Minkowski space it is not until one imposes the Virasoro constraints
$$(L_n+_n\delta _{n,0})|\mathrm{physical}=0,n0$$
(1)
that a unitary spectrum is obtained. Here $`_n`$ is the Virasoro generator for the internal conformal field theory corresponding to $``$. The proposal of is that one should consider not just these representations but also those obtained by the spectral flow
$`J_n^3`$ $``$ $`\stackrel{~}{J}_n^3=J_n^3{\displaystyle \frac{k}{2}}w\delta _{n,0}`$
$`J_n^+`$ $``$ $`\stackrel{~}{J}_n^+=J_{n+w}^+`$ (2)
$`J_n^{}`$ $``$ $`\stackrel{~}{J}_n^{}=J_{nw}^{}.`$
The Virasoro generators, given by the Sugawara form, then become $`\stackrel{~}{L}_n=L_n+wJ_n^3\frac{k}{4}w^2\delta _{n,0}`$. Imposing on $`\widehat{𝒟}_{j>1/2}^\pm \times \widehat{𝒟}_{j>1/2}^\pm `$ the condition (1) with $`\stackrel{~}{L}_n`$ one finds that these states have a discrete energy spectrum
$`E`$ $`=`$ $`J_0^3+\overline{J}_0^3=q+\overline{q}+kw+2\stackrel{~}{j}`$ (3)
$`=`$ $`1+q+\overline{q}+2w+\sqrt{1+4(k2)(N_w+h1{\displaystyle \frac{1}{2}}w(w+1)),}`$
here $`N_w`$ is defined to be the level of the current algebra after spectral flow by amount $`w`$, $`N_w=\stackrel{~}{N}wq`$, and $`\stackrel{~}{N}`$ is the level before spectral flow. The state with energy (3) is obtained from a lowest weight state by acting with the $`SL(2,R)`$ currents $`\stackrel{~}{J}_{n0}^\pm |\stackrel{~}{j},\stackrel{~}{j}`$, with $`q`$ the net number of $`\pm `$ signs in this expression. In other words, $`q`$ is the number of spacetime energy raising operators $`J_a^+`$ minus the number of spacetime energy lowering operators $`J_a^{}`$ that we have to apply to the lowest weight, lowest energy state $`|\stackrel{~}{j},m=\stackrel{~}{j}`$ to get to the state whose spacetime energy is (3). $`\overline{q}`$ is the corresponding quantity for the generators $`\overline{J}_a^\pm `$. We also have a level matching condition of the form
$$N_w+h=\overline{N}_w+\overline{h}$$
(4)
which implies that the angular momentum in $`AdS_3`$, $`\mathrm{}=J_0^3\overline{J}_0^3=q\overline{q}`$, is an integer. We argued in that $`\stackrel{~}{j}`$ is further restricted to the range
$$\frac{1}{2}<\stackrel{~}{j}<\frac{k1}{2},$$
(5)
which implies
$$\frac{k}{4}w^2+\frac{1}{2}w<N_w+h1+\frac{1}{4(k2)}<\frac{k}{4}(w+1)^2\frac{1}{2}(w+1).$$
(6)
A similar analysis on $`\widehat{𝒞}_{j=1/2+is}^\alpha \times \widehat{𝒞}_{j=1/2+is}^\alpha `$ yields a continuous spectrum
$$E=\frac{k}{2}w+\frac{1}{w}\left(\frac{2s^2+\frac{1}{2}}{k2}+\stackrel{~}{N}+h+\stackrel{~}{\overline{N}}+\overline{h}2\right),$$
(7)
where $`s`$ takes values over the real numbers and is interpreted as the momentum in radial direction for the long strings. These states satisfy the level matching condition
$$\stackrel{~}{N}+h=\stackrel{~}{\overline{N}}+\overline{h}+w\times (\mathrm{integer}).$$
(8)
In the rest of the paper we will do an independent calculation which will reproduce this single string spectrum.
## 3 One-loop partition function
In this section we compute the worldsheet one-loop partition function. First we explain the relation between various useful coordinate systems. Then we consider thermal $`AdS_3=H_3/Z`$ and show how the identification of Euclidean time in the global coordinates translates into particular boundary conditions for the target space fields. The partition function is then calculated by an explicit evaluation of the functional integral following .
### 3.1 Coordinates on $`H_3`$ and thermal $`AdS_3`$.
The natural metric on $`H_3`$ is given by
$$ds^2=\frac{k}{y^2}(dy^2+dwd\overline{w}),$$
(9)
which is the Euclidean continuation of the Poincaré metric on $`AdS_3`$. By the coordinate transformation
$`w`$ $`=`$ $`\mathrm{tanh}\rho e^{t+i\theta }`$
$`\overline{w}`$ $`=`$ $`\mathrm{tanh}\rho e^{ti\theta }`$ (10)
$`y`$ $`=`$ $`{\displaystyle \frac{e^t}{\mathrm{cosh}\rho }}`$
we obtain the cylindrical coordinates on Euclidean $`AdS_3`$,
$$\frac{ds^2}{k}=\mathrm{cosh}^2\rho dt^2+d\rho ^2+\mathrm{sinh}^2\rho d\theta ^2.$$
(11)
For the purpose of calculating the partition function, however, it is convenient to use coordinates in which the metric reads
$$\frac{ds^2}{k}=d\varphi ^2+(dv+vd\varphi )(d\overline{v}+\overline{v}d\varphi ),$$
(12)
which corresponds to the parametrization of $`H_3`$ as
$$g=\left[\begin{array}{cc}e^\varphi (1+|v|^2)& v\\ \overline{v}& e^\varphi \end{array}\right].$$
(13)
The coordinate transformation from (11) to (12) is
$`v`$ $`=`$ $`\mathrm{sinh}\rho e^{i\theta }`$
$`\overline{v}`$ $`=`$ $`\mathrm{sinh}\rho e^{i\theta }`$ (14)
$`\varphi `$ $`=`$ $`t\mathrm{log}\mathrm{cosh}\rho .`$
Thermal $`AdS_3`$ is given by the identification
$$t+i\theta t+i\theta +\widehat{\beta },$$
(15)
where $`\widehat{\beta }`$ represents the temperature $`T`$ and the imaginary chemical potential $`i\mu `$ for the angular momentum,
$$\widehat{\beta }=\beta +i\mu \beta =\frac{1}{T}+i\frac{\mu }{T}.$$
(16)
The corresponding identifications in the coordinates (12) are
$`v`$ $``$ $`ve^{i\mu \beta }`$
$`\overline{v}`$ $``$ $`\overline{v}e^{i\mu \beta }`$ (17)
$`\varphi `$ $``$ $`\varphi +\beta ,`$
which is a consistent symmetry of the WZW action,
$$S=\frac{k}{\pi }d^2z\left(\varphi \overline{}\varphi +(\overline{v}+\varphi \overline{v})(\overline{}v+\overline{}\varphi v)\right).$$
(18)
### 3.2 Computation of the partition function on thermal $`AdS_3`$.
In this subsection we compute the partition function for string theory on thermal $`AdS_3`$. We consider a conformal field theory on a worldsheet torus with modular parameter $`\tau `$ ($`zz+2\pi z+2\pi \tau `$). The two-dimensional conformal field theory on the worldsheet is the sum of three parts: the conformal field theory on $`H_3`$, the internal conformal field theory on $``$, and the $`(b,c)`$ ghosts. First we start with the computation of the partition function for the conformal field theory describing the three dimensions of thermal $`AdS_3`$ and then we will multiply the result by the partition function of the ghosts and the internal conformal field theory.
Due to the identification (3.1), the string coordinates now satisfy the following boundary conditions
$`\varphi (z+2\pi )`$ $`=`$ $`\varphi (z)+\beta n,\varphi (z+2\pi \tau )=\varphi (z)+\beta m,`$
$`v(z+2\pi )`$ $`=`$ $`v(z)e^{in\mu \beta },v(z+2\pi \tau )=v(z)e^{im\mu \beta }.`$ (19)
The thermal circle is non-contractible and therefore we get two integers $`(n,m)`$ characterizing topologically nontrivial embeddings of the worldsheet in spacetime. In order to implement these boundary conditions it is convenient to define new fields $`\widehat{\varphi },\widehat{v}`$ such that they are periodic:
$`\varphi `$ $`=`$ $`\widehat{\varphi }+\beta f_{n,m}(z,\overline{z})`$
$`v`$ $`=`$ $`\widehat{v}\mathrm{exp}(i\mu \beta f_{n,m}(z,\overline{z})),`$ (20)
with
$$f_{n,m}(z,\overline{z})=\frac{i}{4\pi \tau _2}\left[z(n\overline{\tau }m)\overline{z}(n\tau m)\right].$$
(21)
When we substitute this into the action (18), we get
$$S=\frac{k\beta ^2}{4\pi \tau _2}|n\tau m|^2+\frac{k}{\pi }d^2z\left(|\widehat{\varphi }|^2+\left|\left(+\frac{1}{2\tau _2}U_{n,m}+\widehat{\varphi }\right)\widehat{\overline{v}}\right|^2\right),$$
(22)
where
$$U_{n,m}(\tau )=\frac{i}{2\pi }(\beta i\mu \beta )(n\overline{\tau }m).$$
(23)
We are interested in the functional integral
$$𝒵(\beta ,\mu ;\tau )=𝒟\varphi 𝒟v𝒟\overline{v}e^S.$$
(24)
This integral is evaluated as explained in . We can first do the integral over $`\widehat{v},\widehat{\overline{v}}`$ which is quadratic, giving the determinant
$$det\left|+\frac{1}{2\tau _2}U_{n,m}+\widehat{\varphi }\right|^2.$$
(25)
We calculate the $`\widehat{\varphi }`$ dependence on the determinants by realizing that we can view (25) as an inverse of two fermion determinants. We can then remove $`\widehat{\varphi }`$ from the determinants by a chiral gauge transformation and using the formulas for chiral anomalies. The result is
$$det\left|+\frac{1}{2\tau _2}U_{n,m}+\widehat{\varphi }\right|^2=e^{\frac{2}{\pi }{\scriptscriptstyle d^2z\widehat{\varphi }\overline{}\widehat{\varphi }}}det\left|+\frac{1}{2\tau _2}U_{n,m}\right|^2.$$
(26)
The remaining integral over $`\widehat{\varphi }`$ gives the usual result for a free boson, except that $`kk2`$ due to (26). The functional integral for the thermal $`AdS_3`$ partition function then gives
$`𝒵(\beta ,\mu ;\tau )`$ (27)
$`=`$ $`{\displaystyle \frac{\beta (k2)^{\frac{1}{2}}}{8\pi \sqrt{\tau _2}}}{\displaystyle \underset{n,m}{}}{\displaystyle \frac{e^{k\beta ^2|mn\tau |^2/4\pi \tau _2+2\pi (\mathrm{Im}U_{n,m})^2/\tau _2}}{|\mathrm{sin}(\pi U_{n,m})|^2|_{r=1}^{\mathrm{}}(1e^{2\pi ir\tau })(1e^{2\pi ir\tau +2\pi iU_{n,m}})(1e^{2\pi ir\tau 2\pi iU_{n,m}})|^2}}`$
$`=`$ $`{\displaystyle \frac{\beta (k2)^{\frac{1}{2}}}{2\pi \sqrt{\tau _2}}}(q\overline{q})^{\frac{3}{24}}{\displaystyle \underset{n,m}{}}{\displaystyle \frac{e^{k\beta ^2|mn\tau |^2/4\pi \tau _2+2\pi (\mathrm{Im}U_{n,m})^2/\tau _2}}{|\vartheta _1(\tau ,U_{n,m})|^2}},`$
where $`\vartheta _1`$ is the elliptic theta function and $`q=e^{2\pi i\tau }`$. The factor $`\beta (k2)^{\frac{1}{2}}`$ comes from the length of the circle in the $`\varphi `$ direction. This partition function is explicitly modular invariant after summing over $`(n,m)`$<sup>2</sup><sup>2</sup>2In our previous paper, there was a puzzle about the apparent lack of modular invariance of the $`SL(2,R)`$ partition functions with $`J^3`$ insertions (see Appendix B of ). Here we have found that, if we introduce the twist by considering the physical set-up of thermal $`AdS_3`$, the result (27) turns out to be manifestly modular invariant. This resolves the puzzle raised in ..
We also need to include the contribution of the $`(b,c)`$ ghosts and the internal CFT. Partition function of the latter will be of the form
$$𝒵_{}=(q\overline{q})^{\frac{c_{int}}{24}}\underset{h,\overline{h}}{}D(h,\overline{h})q^h\overline{q}^{\overline{h}},$$
(28)
where $`D(h,\overline{h})`$ is the degeneracy at left-moving weight $`h`$ and right-moving weight $`\overline{h}`$, and $`c_{int}`$ the central charge of the internal CFT. Modular invariance requires that $`h\overline{h}Z`$, a fact which will be useful in the next section. Vanishing of the total conformal anomaly gives
$$c_{SL(2,R)}+c_{int}=26.$$
(29)
We can calculate now the total contribution to the ground state energy. We found a ground state energy of $`3/24`$ in (27), the ghosts contribute with $`2/24`$ and the internal CFT with $`c_{int}/24=(c_{SL(2,R)}26)/24`$. Using $`c_{SL(2,R)}=3+\frac{6}{k2}`$, we find the overall factor<sup>3</sup><sup>3</sup>3Note that $`c_{int}0`$, $`k>2`$, and (29) imply that there will always be a tachyon in the theory.
$$(q\overline{q})^{(1+c_{int})/24}=e^{4\pi \tau _2\left(1\frac{1}{4(k2)}\right)}.$$
(30)
After multiplying (27) by the $`(b,c)`$ ghosts and the internal CFT partition functions, we should integrate the resulting expression over the fundamental domain $`F_0`$ of the modular parameter $`\tau `$. The computation is much facilitated by the trick invented in to trade the sum over $`n`$ in (27) for the sum over copies of the fundamental domain. See Figure 1. This is possible since $`(n,m)`$ transforms as a doublet under the modular group $`SL(2,Z)`$. If $`(n,m)(0,0)`$, it can be mapped by an $`SL(2,Z)`$ transformation to $`(0,m)`$, $`m>0`$. The $`SL(2,Z)`$ transformation also maps the fundamental domain into the strip $`\mathrm{Im}\tau 0`$, $`|\mathrm{Re}\tau |1/2`$. On the other hand, $`(n,m)=(0,0)`$ is invariant under the $`SL(2,Z)`$ transformation, and the corresponding term still has to be integrated over the fundamental domain $`F_0`$. This term represents the zero temperature contribution to the cosmological constant, or the zero temperature vacuum energy. In addition to the usual tachyon divergence of bosonic string theory at large $`\tau _2`$, it is also divergent due to the $`\mathrm{sin}^1`$ factor in (27); this divergence can be interpreted as coming from the infinite volume of $`AdS_3`$. Since the temperature dependence of this term is trivial we will ignore it from now on. The final result then is that we fix $`n=0`$ in (27) and we integrate over the entire strip shown in Figure 1. The one-loop partition function of bosonic string theory on $`H_3/Z\times `$ is then
$`Z(\beta ,\mu )`$ $`=`$ $`{\displaystyle \frac{\beta (k2)^{\frac{1}{2}}}{8\pi }}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\tau _2}{\tau _2^{3/2}}}{\displaystyle _{1/2}^{1/2}}𝑑\tau _1e^{4\pi \tau _2\left(1\frac{1}{4(k2)}\right)}{\displaystyle \underset{h,\overline{h}}{}}D(h,\overline{h})q^h\overline{q}^{\overline{h}}`$ (31)
$`\times {\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{e^{(k2)m^2\beta ^2/4\pi \tau _2}}{|\mathrm{sinh}(m\widehat{\beta }/2)|^2}}|{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1e^{2\pi in\tau }}{(1e^{m\widehat{\beta }+2\pi in\tau })(1e^{m\widehat{\beta }+2\pi in\tau })}}|^2.`$
## 4 Reading off the spectrum
We will now extract the spectrum of Lorentzian string theory on $`AdS_3`$ by interpreting the one-loop partition function in the spacetime theory. The one-loop partition function is the single particle contribution to the spacetime thermal free energy, $`Z(\beta ,\mu )=\beta F`$. To this each string state makes a contribution $`\beta ^1\mathrm{log}(1e^{\beta (E+i\mu \mathrm{})})`$, where $`E`$ and $`\mathrm{}`$ are the energy and the angular momentum of the state. The total free energy is the sum over all such factors:
$$F(\beta ,\mu )=\frac{1}{\beta }\underset{string}{}\mathrm{log}\left(1e^{\beta (E_{string}+i\mu \mathrm{}_{string})}\right)=\underset{m=1}{\overset{\mathrm{}}{}}f(m\beta ,m\mu ),$$
(32)
where
$$f(\beta ,\mu )=\frac{1}{\beta }\underset{string}{}e^{\beta (E_{string}+i\mu \mathrm{}_{string})}.$$
(33)
Here $``$ is the physical Hilbert space of single string states. In both (31) and (32), we have the sums over functions of $`(m\beta ,m\mu )`$. It is therefore sufficient to compare the $`m=1`$ terms in these expressions. In other words, we want to verify that $`E_{string}`$ and $`\mathrm{}_{string}`$ extracted from the identification,
$`f(\beta ,\mu )={\displaystyle \underset{string}{}}{\displaystyle \frac{1}{\beta }}e^{\beta (E_{string}+i\mu \mathrm{}_{string})}`$ (34)
$`=`$ $`{\displaystyle \frac{(k2)^{\frac{1}{2}}}{8\pi }}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\tau _2}{\tau _2^{3/2}}}{\displaystyle _{1/2}^{1/2}}𝑑\tau _1e^{4\pi \tau _2\left(1\frac{1}{4(k2)}\right)}{\displaystyle \underset{h,\overline{h}}{}}D(h,\overline{h})q^h\overline{q}^{\overline{h}}`$
$`\times {\displaystyle \frac{e^{(k2)\beta ^2/4\pi \tau _2}}{|\mathrm{sinh}(\widehat{\beta }/2)|^2}}\left|{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1e^{2\pi in\tau }}{(1e^{\widehat{\beta }+2\pi in\tau })(1e^{\widehat{\beta }+2\pi in\tau })}}\right|^2`$
agree with the string spectrum found in our previous paper . We will see that the sum over the Hilbert space breaks up into a sum over the discrete states and an integral over the continuous states, with the expressions for the energies that were reviewed in section 2. Since the one-loop calculation presented here is independent of the assumptions made in about strings in Lorentzian $`AdS_3`$, we can regard this as a derivation of the spectrum starting from the well-defined Euclidean path integral.
### 4.1 Qualitative analysis
In this subsection we will analyze (34) in a qualitative way and explain where the different contributions to the spectrum come from. To keep the notation simple, we set $`\mu =0`$ in this subsection, leaving the exact computation for the next subsection.
As expected, in (34) there is an exponential divergence as $`\tau _2\mathrm{}`$, coming from the tachyon. This is just as in the flat space case, where $`(\mathrm{mass})^2<0`$ of the tachyon causes its contribution to be weighted with a positive exponential. We will disregard this divergence<sup>4</sup><sup>4</sup>4 A skeptical reader could think that we are really doing the superstring partition function (the fermions included in the internal CFT, etc.). Then the tachyon divergence will disappear but one would still find the divergences that we discuss below. Of course, the one-loop partition function is non-vanishing even in the supersymmetric case since the thermal boundary conditions break supersymmetry.. However, rather unexpectedly, the expression above has additional divergences at finite values of $`\tau `$. In string theory one might naively expect that divergences come only from the corners of the fundamental domain in the $`\tau `$-plane, but in this case the divergence is coming from points in the interior of the fundamental domain. Overcoming the initial panic, one realizes that these divergences are related to the presence of long strings. In fact, as with any other string divergence, it can be interpreted as an IR effect. This divergence is due to the fact that long strings feel a flat potential as they go to infinity and therefore we get an infinite volume factor. To see this, note that near the pole (see Figure 2)
$$\tau =\tau _{pole}+ϵ,$$
(35)
where
$$\tau _{pole}=\frac{r}{w}+i\frac{\beta }{2\pi w},$$
(36)
we can expand the partition function and replace $`\tau `$ in all terms by its value at the pole, except in the one term that has the pole.
If we integrate (34) near the pole, i.e. in the region $`ϵ<|\tau \tau _{pole}|1`$ , we find that it diverges as $`\mathrm{log}ϵ`$ with coefficient
$$\frac{1}{\sqrt{w\beta ^3}}\mathrm{exp}\left[\beta \left(\frac{k}{2}w+\frac{1}{w}(\stackrel{~}{N}+h+\stackrel{~}{\overline{N}}+\overline{h}2+\frac{1}{2(k2)})\right)+\frac{2\pi ir}{w}(\stackrel{~}{N}+h\stackrel{~}{\overline{N}}\overline{h})\right].$$
(37)
We now sum over $`r`$, with $`|r/w|1/2`$, since these are the ones corresponding to the poles in the strip<sup>5</sup><sup>5</sup>5If some poles are on the boundaries of the strip, $`\tau _1=\pm 1/2`$, then we only count them once. . This sum constrains $`\stackrel{~}{N}+h\stackrel{~}{\overline{N}}\overline{h}`$ to be an integer multiple of $`w`$, as in (8), and it introduces an additional factor of $`w`$ in (37). The log divergence in $`\tau `$-integral can therefore be expressed as
$$f(\beta ,\mu )\frac{1}{\beta }\mathrm{log}ϵ_0^{\mathrm{}}𝑑se^{\beta E(s)}+\mathrm{},$$
(38)
where $`E(s)`$ is the energy spectrum given by (7). Note that the $`s`$-integral and the sum over $`r`$ we mentioned above give the factor $`\sqrt{w/\beta }`$ needed to match the prefactor in (37) to that in (38). This reproduces the expected contribution from the long strings in the left hand side of (34). The logarithmic divergence should be interpreted as a volume factor due to the fact that the long string can be at any radial position. In the next subsections, we will see more precisely that it is indeed associated to the infinite volume in spacetime by relating $`ϵ`$ to a long distance cutoff.
Now we would like to calculate the short string spectrum. Since the long string spectrum gives a divergent result, while the short string spectrum gives a finite one, it might appear at first that extracting the contributions due to the short strings from a divergent expression such as (34) will be problematic. Fortunately we can get around this difficulty since the temperature dependence of the long string free energy is different from that of the short string free energy. In the next subsection we will explain how to do this precisely and reproduce the short string spectrum which agrees with . One of the more puzzling aspects of the short string spectrum found there is that there is a cutoff $`1/2<\stackrel{~}{j}<(k1)/2`$ in the value of the $`SL(2,R)`$ spin $`\stackrel{~}{j}`$. In the remainder of this section we will explain in a qualitative way how this cutoff arises by doing the calculation for large $`k`$.
If we were to evaluate the right hand side of (34) naively (and incorrectly), we would expand the integrand in powers of $`q=e^{2\pi i\tau }`$ and then perform the $`\tau `$ integral. If we did this, we would obtain the short string spectrum with $`w=0`$ and no upper bound on the value of $`\stackrel{~}{j}`$. However this expansion is not correct. How we can expand the integrand in (34) depends on the value of $`\tau _2`$. When we cross the poles at $`\tau _2=\frac{\beta }{2\pi w}`$, a different expansion should be used for the denominator:
$`{\displaystyle \frac{1}{1e^{\beta +2\pi iw\tau }}}`$ $`=`$ $`{\displaystyle \underset{q=0}{\overset{\mathrm{}}{}}}e^{q(\beta +2\pi iw\tau )},\left(\tau _2>{\displaystyle \frac{\beta }{2\pi w}}\right),`$ (39)
$`=`$ $`{\displaystyle \underset{q=0}{\overset{\mathrm{}}{}}}e^{(q+1)(\beta +2\pi iw\tau )},\left(\tau _2<{\displaystyle \frac{\beta }{2\pi w}}\right).`$ (40)
When $`\tau _2`$ is in the range
$$\frac{\beta }{2\pi (w+1)}<\tau _2<\frac{\beta }{2\pi w},$$
(41)
the product over $`n`$ in the first term in the denominator in (34) is broken up into two factors, a product in $`1nw`$ and a product in $`w+1n`$. The first factor is expanded in powers of $`e^{(\beta +2\pi in\tau )}`$ and the second factor is expanded in powers of $`e^{\beta +2\pi in\tau }`$. Combining them together with the terms coming from the expansion of the remaining products in (34), we get an exponent of the form<sup>6</sup><sup>6</sup>6 The first term $`\beta /2`$ comes from expanding $`1/\mathrm{sinh}(\beta /2)`$ in (34).
$$\left(\frac{1}{2}+q+w\right)\beta +2\pi i\tau \left(N_w\frac{1}{2}w(w+1)\right),$$
(42)
for some integers $`q`$ and $`N_w`$. There is a similar term for $`\tau \overline{\tau }`$. We are then to do the $`\tau `$-integral of the form,
$$\frac{d^2\tau }{\tau _2^{3/2}}e^{4\pi \tau _2\left(1\frac{1}{4(k2)}\right)(k2)\frac{\beta ^2}{4\pi \tau _2}\beta (1+q+\overline{q}+2w)+2\pi i\tau \left(N_w+h\frac{1}{2}w(w+1)\right)2\pi i\overline{\tau }\left(\overline{N}_w+\overline{h}\frac{1}{2}w(w+1)\right)},$$
(43)
over the region (41). The integral over $`\tau _1`$ produces the level matching condition (4). Now we evaluate the integral over $`\tau _2`$ using the saddle point approximation. We find that the saddle point is at
$$\tau _{saddle}=\frac{(k2)\beta }{2\pi \sqrt{1+4(k2)(N_w+h1\frac{1}{2}w(w+1))}}$$
(44)
and the integral gives
$$\frac{1}{\beta }\mathrm{exp}\left[\beta \left(1+q+\overline{q}+2w+\sqrt{1+4(k2)(N_w+h1\frac{1}{2}w(w+1))}\right)\right].$$
(45)
This is the correct form of the contributions due to the short strings in the left hand side of (34). Moreover we obtain the bound on $`\stackrel{~}{j}`$ precisely, because $`\tau _{saddle}`$ has to be in the range (41) in order for the saddle point approximation to give a non-zero result. By (44), this condition is the same as the bound on the spectrum (6), which is equivalent to $`1/2<\stackrel{~}{j}<(k1)/2`$. (It is a bit surprising that we get all factors precisely right from the saddle point approximation.) Notice then that the cutoff in $`\stackrel{~}{j}`$ is associated to the fact that we expand the integrand in (34) in different ways depending on the value of $`\tau `$. The value of $`\tau `$ making the biggest contribution to the integral depends on the values of $`N`$ and $`h`$ of the string state.
### 4.2 A precise evaluation of the $`\tau `$-integral
Now let us study the partition function (34) more systematically. In this subsection, we go back to the general case with $`\mu 0`$. From what we saw in the previous subsection, we expect to find the discrete states from the integral over the range (41), and the continuous states from the poles after a suitable regularization.
In order to evaluate the $`\tau `$-integral exactly, it is useful to introduce a new variable $`c`$ by
$$e^{(k2)\frac{\beta ^2}{4\pi \tau _2}}=\frac{8\pi i}{\beta }\left(\frac{\tau _2}{k2}\right)^{\frac{3}{2}}_{\mathrm{}}^{\mathrm{}}𝑑cce^{\frac{4\pi \tau _2}{k2}c^2+2i\beta c}.$$
(46)
Now suppose $`\tau _2`$ is in the range,
$$\frac{\beta }{2\pi (w+1)}<\tau _2<\frac{\beta }{2\pi w},$$
(47)
and expand the integrand in (34) as explained in the previous subsection. The right hand side of (34) becomes a sum of terms like
$`{\displaystyle \frac{4}{\beta (k2)i}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑cc{\displaystyle _{\frac{\beta }{2\pi (w+1)}}^{\frac{\beta }{2\pi w}}}𝑑\tau _2{\displaystyle _{1/2}^{1/2}}𝑑\tau _1`$
$`\times \mathrm{exp}[\widehat{\beta }(q+w+{\displaystyle \frac{1}{2}})\overline{\widehat{\beta }}(\overline{q}+w+{\displaystyle \frac{1}{2}})+2\pi i\tau _1(N_w+h\overline{N}_w\overline{h})`$
$`+2ic\beta 2\pi \tau _2(h+\overline{h}+N_w+\overline{N}_w+{\displaystyle \frac{2c^2+\frac{1}{2}}{k2}}w(w+1)2)].`$ (48)
The integral over $`\tau _1`$ gives a delta function enforcing $`N_w+h=\overline{N}_w+\overline{h}`$, which is the level-matching condition (4). Integrating over $`\tau _2`$ in the range (47) gives
$`{\displaystyle \frac{1}{\beta \pi i}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑cc{\displaystyle \frac{\mathrm{exp}\left[2ic\beta \widehat{\beta }\left(q+w+\frac{1}{2}\right)\overline{\widehat{\beta }}\left(\overline{q}+w+\frac{1}{2}\right)\right]}{c^2+\frac{1}{4}+(k2)\left(N_w+h1\frac{1}{2}w(w+1)\right)}}`$
$`\times \{\mathrm{exp}[{\displaystyle \frac{\beta }{w}}(2N_w+2h2+{\displaystyle \frac{2c^2+\frac{1}{2}}{k2}}w(w+1))]`$
$`+\mathrm{exp}[{\displaystyle \frac{\beta }{w+1}}(2N_w+2h2+{\displaystyle \frac{2c^2+\frac{1}{2}}{k2}}w(w+1))]\}`$ (49)
where we used (4).
Let us first look at the first term (the second line) in (49). We note that the exponent can be expressed in the form of a complete square if we set $`c=s+\frac{i}{2}(k2)w`$. As it will become clear shortly, it is natural to shift the contour of the $`c`$-integral from $`\mathrm{Im}c=0`$ to $`\mathrm{Im}c=\frac{1}{2}(k2)w`$ so that $`s`$ becomes real. During this process the contour crosses some poles in the integrand, picking up the residues of the poles in the range $`0<\mathrm{Im}c<\frac{1}{2}(k2)w`$. See Figure 3. The poles are located at
$$\frac{c^2}{(k2)}=N_w+h\frac{1}{2}w(w+1)1+\frac{1}{4(k2)}<\frac{k2}{4}w^2.$$
(50)
Similarly, for the second exponential term (the third line) in (49) we shift the contour to $`c=s+\frac{i}{2}(k2)(w+1)`$ with $`s`$ real. This picks up the poles at
$$\frac{c^2}{(k2)}=N_w+h\frac{1}{2}w(w+1)1+\frac{1}{4(k2)}<\frac{k2}{4}(w+1)^2.$$
(51)
It is important to note that the residues of these poles have a sign opposite to that of the residues of the poles obeying (50). So the result is that we are left with only those poles in the range
$$\frac{k2}{2}w<\mathrm{Im}c<\frac{k2}{2}(w+1),$$
(52)
with residues
$$\frac{1}{\beta }\mathrm{exp}\left[\widehat{\beta }q\overline{\widehat{\beta }}\overline{q}\beta \left(1+2w+\sqrt{1+4(k2)(N_w+h1\frac{1}{2}w(w+1))}\right)\right].$$
(53)
This is the expected contribution of the short strings to the right hand side of (34), and we see also that (52) translates into the correct bound on $`\stackrel{~}{j}`$ (5).
It remains to examine the resulting integral over $`s`$. Notice that the term coming from just above the pole at $`\tau =\widehat{\beta }/2\pi w`$ has a very similar $`w`$ dependence in the exponent as that coming from just below the pole. In other words, we combine the first term of (49) with the second term of an expression similar to (49) but with $`ww1`$ and we find, after shifting the countours as above,
$`{\displaystyle \frac{1}{2\pi i\beta }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑s\left({\displaystyle \frac{2s}{w(k2)}}+i\right)`$ (54)
$`\times ({\displaystyle \frac{\mathrm{exp}\left[\widehat{\beta }q\overline{\widehat{\beta }}\overline{q}\beta \left(\frac{k}{2}w+\frac{2}{w}\left(\frac{s^2+1/4}{k2}+N_{w1}+h1\right)\right)\right]}{\frac{1}{2}+is\frac{k}{4}w+\frac{1}{w}\left(N_{w1}+h1+\frac{s^2+1/4}{k2}\right)}}`$
$`{\displaystyle \frac{\mathrm{exp}\left[\widehat{\beta }q\overline{\widehat{\beta }}\overline{q}\beta \left(\frac{k}{2}w+\frac{2}{w}\left(\frac{s^2+1/4}{k2}+N_w+h1\right)\right)\right]}{\frac{1}{2}+is\frac{k}{4}w+\frac{1}{w}\left(N_w+h1+\frac{s^2+1/4}{k2}\right)}}).`$
Let us concentrate for now on the third line of (54). We first note that the sum of such terms over all states gives rise to the log divergence. To see this, it is useful to notice that the combinations
$$\stackrel{~}{N}=qw+N_w,\stackrel{~}{\overline{N}}=\overline{q}w+\overline{N}_w$$
(55)
that appear in the exponent of the third line of equation (54) are the levels before spectral flow. Thus, for a given state $`|\psi `$, states of the form $`(\stackrel{~}{J}_0^+\stackrel{~}{\overline{J}_0^+})^n|\psi `$ all have the same value of $`\stackrel{~}{N}`$ and $`\stackrel{~}{\overline{N}}`$. Acting with $`\stackrel{~}{J}_0^+\stackrel{~}{\overline{J}_0^+}`$ on $`|\psi `$ does not change the exponent in (54), but it does change the denominator by one. This implies that when we sum over all the states of this type, we will find a divergent sum of the form
$$\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{An}.$$
This divergence has the same origin as the divergence of the right hand side of (34) at the pole $`\tau _{pole}=\widehat{\beta }/2\pi w`$. In fact, if we regularize the $`\tau `$-integral by removing a small region near the pole as $`|\tau \tau _{pole}|>ϵ`$, we find an additional factor $`e^{nϵ}`$ in the sum. In the next subsection, we will give the spacetime interpretation of this procedure. With this regularization, the sum behaves as $`\mathrm{log}ϵ`$. More precisely we have
$$\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{An}e^{nϵ}=\mathrm{log}ϵ+\frac{d}{dA}\mathrm{log}\mathrm{\Gamma }(A)+𝒪(ϵ)$$
(56)
where
$$A=\frac{1}{2}+is\frac{k}{4}w+\frac{1}{w}\left(\frac{s^2+\frac{1}{4}}{k2}+\stackrel{~}{N}+h1\right).$$
(57)
Here we have assumed that
$$\stackrel{~}{\overline{N}}+\overline{h}\stackrel{~}{N}+h,$$
(58)
but it can be seen that the other case gives the same result.
Now we turn our attention to the second line of (54). In those terms we have one less unit of spectral flow, as compared to the third line in (54) that we analyzed above. In other words, now we will have that $`(w1)q+N_{w1}=\stackrel{~}{N^{}}`$. These states are in the spectral flow image of $`𝒟_j^+`$. Since we want to combine these states with the states coming from the third line in (54) it is convenient to do spectral flow one more time and think of these states as in the spectral flow image of $`𝒟_j^{}`$ under $`w`$ units of spectral flow. In this case we find that $`q+\stackrel{~}{N^{}}=\stackrel{~}{N}`$ where now $`\stackrel{~}{N}`$ is the level of the $`𝒟_j^{}`$ representation before spectral flow. From now on the discussion is very similar to what we had above. The states with $`(\stackrel{~}{J}_0^{}\stackrel{~}{\overline{J}_0^{}})^n|\psi `$ all have the same energies but they will contribute to the denominator of the second line in (54) with
$$\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{B+n}e^{nϵ}=\mathrm{log}ϵ\frac{d}{dB}\mathrm{log}\mathrm{\Gamma }(B)+𝒪(ϵ)$$
(59)
where
$$B=\frac{1}{2}+is\frac{k}{4}w+\frac{1}{w}\left(\frac{s^2+\frac{1}{4}}{k2}+\stackrel{~}{\overline{N}}+\overline{h}1\right),$$
(60)
again assuming (58).
After we perform these two sums, we find that (54) can be written in the form
$$\frac{2}{\beta }_0^{\mathrm{}}𝑑s\rho (s)\mathrm{exp}\left[\beta \left(E(s)+i\frac{\mu }{w}(\stackrel{~}{N}+h\stackrel{~}{\overline{N}}\overline{h})\right)\right]$$
(61)
with $`E(s)`$ the energy of long strings (7) and $`\rho (s)`$ the density of states. We will later see that the physical momentum $`p`$ of a long string in the $`\rho `$ direction is equal to $`p=2s`$. The angular momentum $`\mathrm{}=(\stackrel{~}{N}+h\stackrel{~}{\overline{N}}\overline{h})/w`$ is an integer since the states in (54) were obeying (4) and the definition (55) ensures that (8) is satisfied. The density of states $`\rho (s)`$ derived from this analysis is
$$\rho (s)=\frac{1}{2\pi }2\mathrm{log}ϵ+\frac{1}{2\pi i}\frac{d}{2ds}\mathrm{log}\left(\frac{\mathrm{\Gamma }(\frac{1}{2}is+\stackrel{~}{\overline{m}})\mathrm{\Gamma }(\frac{1}{2}is\stackrel{~}{m})}{\mathrm{\Gamma }(\frac{1}{2}+is+\stackrel{~}{\overline{m}})\mathrm{\Gamma }(\frac{1}{2}+is\stackrel{~}{m})}\right),$$
(62)
where
$$\stackrel{~}{m}=\frac{k}{4}w+\frac{1}{w}\left(\frac{s^2+\frac{1}{4}}{k2}+\stackrel{~}{N}+h1\right),\stackrel{~}{\overline{m}}=\frac{k}{4}w+\frac{1}{w}\left(\frac{s^2+\frac{1}{4}}{k2}+\stackrel{~}{\overline{N}}+\overline{h}1\right).$$
(63)
Note that, despite appearances to the contrary, (62) is actually symmetric under $`\stackrel{~}{m}\stackrel{~}{\overline{m}}`$ since $`\stackrel{~}{m}\stackrel{~}{\overline{m}}=\mathrm{}`$ is an integer. In the next subsection we will show that this density of states (62) is what is expected from the spacetime meaning of the cutoff $`ϵ`$. In going from (54) to (61) we have states which could be interpreted as coming from the spectral flow of the discrete representations $`𝒟_j^+`$ and $`𝒟_j^{}`$, with the zero modes essentially stripped off since they were explicitly summed over in (56) and (59). This implies that the states we have in the end belong to the continuous representation. Note also that the integral over $`s`$ in (61) has only half the range in (54). We rewrote it in this way using the fact that the exponent is invariant under $`ss`$, and that is the reason why we have four Gamma functions in (62). In going from (54) to (61) we have also used that $`\frac{d}{dA}=\frac{1}{\frac{dA(s)}{ds}}\frac{d}{ds}`$ in (57) and similarly in (60).
Combining eqns. (53) and (61), we have finally
$$f(\beta ,\mu )=\frac{1}{\beta }D(h,\overline{h},\stackrel{~}{N},\stackrel{~}{\overline{N}},w)\left[\underset{q,\overline{q}}{}e^{\beta (E+i\mu \mathrm{})}+_0^{\mathrm{}}𝑑s\rho (s)e^{\beta (E(s)+i\mu \mathrm{})}\right]$$
(64)
which is the free energy due to the short strings and the long strings, respectively.
### 4.3 The density of long string states
What remains to be shown is the interpretation of $`\rho (s)`$ given by (62) as the density of long string states. Whenever we have a continuous spectrum the density of states may be calculated by first introducing a long distance cutoff which will make the spectrum discrete, and then removing the cutoff. If the cutoff is related to the volume of the system then the density of states will have a divergent part, proportional to the volume and dependent only on the bulk physics, and a finite part which encodes information about the scattering phase shift and also has some dependence on the precise cutoff procedure. To see this, let us consider a one-dimensional quantum mechanical model on the half line, $`\rho >0`$, with a potential $`V(\rho )`$. We assume that $`V(\rho )`$ vanishes sufficiently fast for large $`\rho `$, and that there is continuous spectrum above a certain energy level. To define the density of states, it is convenient to introduce a long distance cutoff at large $`\rho `$ so that the spectrum becomes discrete. Let us first consider a cutoff by an infinite wall at $`\rho =L`$. If $`L`$ is sufficiently large, an energy eigenfunction $`\psi (\rho )`$ near the wall can be approximated by the plane wave
$$\psi (\rho )e^{ip\rho }+e^{ip\rho +i\delta (p)},$$
(65)
where $`\delta (p)`$ is the phase shift due to the original potential $`V(\rho )`$. Imposing Dirichlet boundary condition $`\psi (L)=0`$ at the wall, we have
$$2pL+\delta (p)=2\pi \left(n+\frac{1}{2}\right)$$
(66)
for some integer $`n`$. If $`L`$ is sufficiently large, there is a unique solution $`p=p(n)`$ to this equation for a given $`n`$. As we remove the cutoff by sending $`L\mathrm{}`$, the spectrum of $`p`$ becomes continuous. We then define the density of states $`\rho (p)`$ by
$$dn=\rho (p)dp.$$
(67)
From (66), we obtain
$$\rho (p)=\frac{1}{2\pi }\left(2L+\frac{d\delta }{dp}\right).$$
(68)
Thus the finite part of the density of states is given by the derivative of the phase shift.
Instead of the infinite wall at $`\rho =L`$, we may consider a more general potential $`V_{wall}(\rho L)`$ which vanishes for $`\rho <L`$ but rises steeply for $`L<\rho `$ to confine the particle. Let us denote by $`\delta _{wall}(p)`$ the phase shift due to scattering from $`V_{wall}(\rho )`$. We then obtain the condition on the allowed values of momenta by matching these two wavefunctions and their derivatives at $`\rho =L`$ as
$$\psi (\rho )e^{ip\rho }+e^{ip\rho +i\delta (p)}A\left[e^{ip(\rho L)}+e^{ip(\rho L)+i\delta _{wall}(p)}\right],(\rho L).$$
(69)
It follows that
$$pL+\delta (p)=pL+\delta _{wall}(p)+2\pi n.$$
(70)
In the limit $`L\mathrm{}`$, the density of states given by $`dn=\rho (p)dp`$ is then
$$\rho (p)=\frac{1}{2\pi }\left(2L+\frac{d\delta }{dp}\frac{d\delta _{wall}}{dp}\right).$$
(71)
When we have the infinite wall, the phase shift due to the wall is independent of $`p`$ ($`\delta _{wall}=\pi `$), and (71) reduces to (68).
In order to apply this observation to our problem, it is useful to first identify the origin of the logarithmic divergence in the one-loop amplitude $`Z(\beta ,\mu )`$ by examining the functional integral (24) near the boundary of $`AdS_3`$. In the cylindrical coordinates (11), the string worldsheet action (18) for large $`\rho `$ takes the form
$$S\frac{k}{\pi }d^2z\left(\rho \overline{}\rho +\frac{1}{4}e^{2\rho }|\overline{}(\theta it)|^2+\mathrm{}\right).$$
(72)
Because of the factor $`e^{2\rho }`$, the functional integral for large $`\rho `$ restricts $`(t,\theta )`$ to be a harmonic map from the worldsheet to the target space. Since $`(t,\theta )`$ are coordinates on the torus,
$$\theta it\theta it+2\pi n+i\widehat{\beta }m,(n,m\mathrm{integers}),$$
(73)
the harmonic map from the torus to the torus is
$`\theta it`$ $`=`$ $`(2\pi w+i\widehat{\beta }m)\sigma ^1+(2\pi r+i\widehat{\beta }n)\sigma ^2`$ (74)
$`=`$ $`\left[(2\pi w+i\widehat{\beta }m)\tau (2\pi r+i\widehat{\beta }n)\right]{\displaystyle \frac{\overline{z}}{2i\tau _2}}`$
$`\left[(2\pi w+i\widehat{\beta }m)\overline{\tau }(2\pi r+i\widehat{\beta }n)\right]{\displaystyle \frac{z}{2i\tau _2}},`$
where $`z=\sigma ^1+\tau \sigma ^2`$ is the worldsheet coordinate and $`(r,w,n,m)`$ are integers. In particular, the map $`(\theta it)`$ with $`(n,m)=(1,0)`$ becomes $`w`$-to-$`1`$ and holomorphic when $`\tau `$ takes the special value
$$\tau _{pole}=\frac{r}{w}+i\frac{\widehat{\beta }}{2\pi w}.$$
(75)
On the other hand, if $`\tau `$ is not at one of these points, $`\overline{}(\theta it)`$ cannot be set to zero<sup>7</sup><sup>7</sup>7 For any $`\tau `$, we also have a trivial holomorphic map $`(t,\theta )=\mathrm{const}`$. The functional integral around such a map gives a result independent of $`\beta `$ and we can neglect it in the following discussion.. This gives rise to an effective potential $`e^{2\rho }`$ for $`\rho `$, which keeps the worldsheet from growing towards the boundary. If $`\tau `$ is near $`\tau _{pole}`$
$$\tau =\tau _{pole}+ϵ,$$
(76)
the harmonic map (74) with $`(n,m)=(1,0)`$ gives
$$|\overline{}(\theta it)|^2\left(\frac{2\pi ^2w^2}{\beta }\right)^2ϵ^2.$$
(77)
Thus the action (72) generates the Liouville potential $`ϵ^2e^{2\rho }`$. When we computed the one-loop amplitude in sections 4.1 and 4.2, we regularized the $`\tau `$-integral by removing a small disk $`|\tau \tau _{pole}|<ϵ`$ around each of these special points. Near $`\tau =\tau _{pole}`$, this is equivalent to adding the infinitesimal Liouville potential $`ϵ^2e^{2\rho }`$ to the worldsheet action. For $`|\tau \tau _{pole}|ϵ`$, the worldsheet can never grow large enough and the effect of the Liouville term is negligible. To be precise, the Gaussian functional integral of $`(t,\theta )`$ shifts $`k(k2)`$ as in (26) and the effective action for $`\rho `$ near $`\tau =\tau _{pole}`$ is
$`S_{Liouville}={\displaystyle \frac{k2}{\pi }}{\displaystyle d^2z\left(\rho \overline{}\rho +ϵ^2e^{2\rho }\right)}.`$ (78)
Therefore, we find that our choice of regularization in (56) and (59) amounts to introducing the Liouville wall which prevents the longs strings from going to very large values of $`\rho `$. By looking at the potential in (78), we see that the effective length of the interval is $`L\mathrm{log}ϵ`$. The central charge of this Liouville theory is such that the $`e^{2\rho }`$ term has conformal weight one,
$$c_{Liouville}=1+6\left(b+\frac{1}{b}\right)^2,b\frac{1}{\sqrt{k2}}.$$
(79)
The finite part of the density of states will be given through (71) by $`\delta (s)`$, the phase shift in the $`SL(2,R)`$ model, and $`\delta _{wall}(s)`$, the corresponding quantity in Liouville theory. The first one was calculated in ,
$$i\delta (s)=\mathrm{log}\left(\frac{\mathrm{\Gamma }(\frac{1}{2}+is\stackrel{~}{m})\mathrm{\Gamma }(\frac{1}{2}+is+\stackrel{~}{\overline{m}})\mathrm{\Gamma }(2is)\mathrm{\Gamma }(\frac{2is}{k2})}{\mathrm{\Gamma }(\frac{1}{2}is\stackrel{~}{m})\mathrm{\Gamma }(\frac{1}{2}is+\stackrel{~}{\overline{m}})\mathrm{\Gamma }(2is)\mathrm{\Gamma }(\frac{2is}{k2})}\right),$$
(80)
while the second one was obtained in <sup>8</sup><sup>8</sup>8 In order to compare with the expressions in we use the value of $`b`$ given in (79) and note that the relevant values of $`\alpha `$ are $`\alpha =Q/2+isb`$.
$$i\delta _{wall}(s)=\mathrm{log}\left(\frac{\mathrm{\Gamma }(2is)\mathrm{\Gamma }(\frac{2is}{k2})}{\mathrm{\Gamma }(2is)\mathrm{\Gamma }(\frac{2is}{k2})}\right).$$
(81)
Using these two formulas we can check that indeed the density of states (62) is given by (71). We can view this as an independent calculation of (80) or as an overall consistency check. Notice that the physical momentum $`p`$ of a long string along the $`\rho `$ direction is $`p=2s`$. This can be seen by comparing the energy of a long string (7) with the energy expected from (78) with spacetime momentum $`p`$ along the radial direction, $`p=(k2)w\dot{\rho }`$. We have chosen the variable $`s`$ since it is conventional to denote by $`j=1/2+is`$ the $`SL(2,R)`$ spin of a continuous representation.
## Acknowledgements
H.O. would like to thank J. Schwarz and the theory group at Caltech for the kind hospitality while this work was carried out.
The research of J.M. was supported in part by DOE grant DE-FGO2-91ER40654, NSF grant PHY-9513835, the Sloan Foundation and the David and Lucile Packard Foundations. The research of H.O. was supported in part by NSF grant PHY-95-14797, DOE grant DE-AC03-76SF00098, and the Caltech Discovery Fund.
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# Vortex melting and decoupling transitions in YBa2Cu4O8 single crystals
## I Introduction
The vortex dynamics in the mixed state of high T<sub>c</sub> superconductors (HTSCs) remains a topic subject to intense investigation because of its importance for both the fundamental research and future applications. It is now well established that the thermal fluctuation induces a second order transition from a high temperature vortex liquid phase to a low temperature vortex glass phase for the vortex matter in a superconductor with a strong disorder. In a clean superconductor the second order transition is replaced by a first order transition from a vortex liquid to an Abrikosov vortex lattice.
Flux melting in various high quality high T<sub>c</sub> superconducting single crystals has been observed by different experimental techniques. It demonstrates itself as a sharp resistive drop or a jump in the magnetization. An important issue concerning the melting transition is that if the vortices lose their coherence in the c direction during the melting transition or they will keep the c-axis correlations intact and then lose them at a still higher temperature. This issue has been widely pursued by doing transport measurements with the dc flux transformer configuration. So far the results remain controversial. Doyle et al. and Fuchs et al. found a coincidence of the melting transition and a decoupling one on Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8</sub> (Bi2212) single crystals. However, Wan et al. and Keener et al. did similar measurements on Bi2212 single crystals. They concluded that the melting transition took place in a two-stage fashion, i.e., the rigid lattice at low temperature first melts into a three dimensional (3D) liquid followed by a decoupling transition upon an increase in the temperature. Recently Blasius et al. performed $`\mu `$-spin rotation measurements on Bi2212 single crystals with different oxygen contents and they obtained evidence for a two stage transition of the vortex matter as a function of temperature under equilibrium conditions. On the other hand, in the YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> (Y123) system a coincidence of vortex melting and loss of vortex correlation along the c-axis was observed by López et al. A decoupling transition was not observed in the optimally doped Y123 system since the vortex in this system is always in the 3D region.
Underdoped YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub> (Y124) is naturally stoichiometric and untwinned, it has a moderate anisotropy parameter $`\gamma `$ as compared to those of Y123 and Bi2212. The weakness of the pinning strength in this material reduces the contribution from spurious effects in the melting transition in the vortex system. Thus it is a good prototype material for the investigation of the interplay between the vortex interaction and dimensionality on the vortex dynamics. Recently, we observed a first order melting transition in the Y124 single crystals. It is of great interest to extend our previous work to examine the vortex dynamics along the c-axis. Here we report the transport data obtained in simultaneous measurements for both the primary (V<sub>top</sub>) and secondary (V<sub>bot</sub>) voltages using the dc transformer configuration, in order to investigate the vortex correlation along the c-axis in the Y124 system. It is demonstrated that in Y124 single crystals, vortex lattice (VL) melts into a 3D vortex liquid via a first order transitionand and that the decoupling transition is seperated from the melting one. Our results support the two stage transition scenario that the vortices first undergo a melting transition, followed by a disappearance of correlation in the transverse direction at a higher temperature.
## II Experimental and Results
YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub> single crystals were grown by the high-pressure flux method as described previously. These single crystals were needle-like with typical dimensions of 1.2$`\times `$0.4$`\times `$0.01 mm<sup>3</sup>. The zero field critical transition temperature T<sub>c</sub>(0) of the single crystals was about 78 K. Two single crystals were used for these measurements. Each crystal was carefully cleaved to obtain optically flat surfaces with the c-axis normal to the sample surface. Gold wires were attached to the top and bottom surfaces of crystal by using Platinum epoxy. Then the crystal was heated in air at 100 C for 1 hour, the resultant contact resistance was typically below 0.5 $`\mathrm{\Omega }`$. The electrical contact geometry for the measurement is shown in the inset of Fig. 1. The secondary voltage contact pair (5 and 6) on the bottom ab face was placed directly beneath the top face primary voltage pair (2 and 3). The distance between each voltage pair was about 0.5 mm. The resistance were measured using a low frequency (17 Hz) ac lock-in technique with an excitation current of 0.1 mA injected from contact 1 to 4 while measuring simultaneously the voltage drops V<sub>23</sub> and V<sub>56</sub> with the help of a Keithley 228 scanner. The magnetic field was generated by a 15 Tesla Oxford superconducting magnet and was applied parallel to the c-axis of the crystal throughout the measurements.
Shown in Fig. 1 are the representative R<sub>top</sub>(T) and R<sub>bot</sub>(T) (defined as V<sub>23</sub>/I<sub>14</sub> and V<sub>56</sub>/I<sub>14</sub>, respectively) curves taken at different applied magnetic fields up to 12 Telsa for one of the Y124 single crystals. Similar results were obtained for the other crystal. The resistivity was measured by first cooling down the sample below T<sub>c</sub> and then collecting the data while warming up the sample. While the applied magnetic fields induce a negligible depression in the superconducting onset temperature for the in plane superconductivity, they appear to considerably suppress the temperature where R<sub>bot</sub> deviates from its normal state value. A sharp jump in R<sub>top</sub> with a magnitude of R/R<sub>n</sub> $``$ 10% (R<sub>n</sub> being the normal state resistance) is clearly observed at each magnetic field. We attribute it to the occurrence of the vortex melting transition. We define the melting transition temperature T<sub>m</sub> as the one where a sharp peak in dR/dT exists. The temperature T<sub>m</sub> defined by R<sub>top</sub>(T) and R<sub>bot</sub>(T) are equal, i.e., T$`{}_{}{}^{top}{}_{m}{}^{}`$=T$`{}_{}{}^{bot}{}_{m}{}^{}`$. This is demonstrated more clearly by plotting the derivatives of R<sub>top</sub> and R<sub>bot</sub> with respect to T together, as shown in Fig. 2a. Below the melting temperature, R<sub>top</sub>=R<sub>bot</sub>. Immediately above the melting temperature T<sub>m</sub>, R<sub>top</sub> becomes larger than R<sub>bot</sub>. As the temperature increases, both R<sub>top</sub> and R<sub>bot</sub> increase and the difference between R<sub>top</sub> and R<sub>bot</sub> becomes larger and larger. Finally near a characteristic temperature T<sub>d</sub>, R<sub>bot</sub> reaches its maximum and at the same temperature, a small kink develops in R<sub>top</sub>, as can be seen in a separated plot in Fig. 2b. It is interesting to notice that above T<sub>m</sub>, dR<sub>top</sub>/dT decreases with the increasing temperature, reaches a minimum at T<sub>d</sub> before a second local maximum appears at T<sub>p</sub>. At exactly the same temperature T<sub>p</sub>, dR<sub>bot</sub>/dT shows a local minimal. This behavior persists for the dR/dT curves at all applied magnetic fields. By plotting T<sub>m</sub> and T<sub>d</sub> as a function of the magnetic field, we obtain a phase diagram for this Y124 single crystal as shown in Fig. 3. Previously we have found that T<sub>m</sub>(B) could be well described by the anisotropic Ginzburg-Landau theory and the melting line could be fitted by an empirical formula B<sub>m</sub>(T)=31.4(1-T/T<sub>c</sub>)<sup>1.44</sup> with T<sub>c</sub>=78.6 K. The fitting result is shown as the solid line in Fig.3. In the mean time, angular dependence of the melting transition gives for Lindermann criterium c<sub>L</sub>=0.14 and the anisotropy parameter $`\gamma `$=12.4.
## III Discussions
The dc flux transformer configuration is useful means for probing the dimensionality and longitudinal correlation of vortices in the mixed state. It measures the velocity correlation for vortices running along the CuO planes of Y124 since the voltage drops V<sub>23</sub> and V<sub>56</sub> (and thus R<sub>top</sub> and R<sub>bot</sub>) in the mixed state is induced by the motion of vortices whose velocity is determined by the integrated Lorentz force on the vortex segments of length l<sub>c</sub> piercing the top and bottom surfaces, respectively. The relation between R<sub>top</sub> and R<sub>bot</sub> is determined by the velocity of the vortices in the top and bottom surfaces. When the vortex correlation length along the c-axis is longer than the sample thickness, the voltage drops at the top and bottom surfaces would be of the same magnitude. The fact that immediately above the melting temperature R<sub>top</sub> is larger than R<sub>bot</sub> suggests that the l<sub>c</sub> is smaller than the sample thickness t above the melting transition.
Above T<sub>c</sub> in the normal state, R<sub>top</sub> is larger than R<sub>bot</sub> because the sample geometrical aspect, electrode arrangement and anisotropy result in a nonuniform current distribution, i.e., more current will pass through the top surface than on the bottom surface. The fact that in the normal state the ratio R<sub>top</sub>/R<sub>bot</sub> is nearly constant suggests that the anisotropy of Y124 is nearly temperature independent in this temperature region. Upon entering the superconducting state, the superconductivity will make an easy path for electrons to run across the CuO planes and therefore reduces the anisotropy. As we will discuss below, the vortices at this temperature region is the 2D pancakes, the voltage drop will be mainly determined by the anisotropy. Since we are measuring the resistance by passing a constant current, the decrease in anisotropy means that more current will pass through the CuO layers and that the current flowing in the bottom layers will increase. The decreasing anisotropy with decreasing temperature is reflecting in the reduction of dR<sub>bot</sub>/dT as shown in Fig.2a.
In the present measurement configuration, R<sub>top</sub> is equivalent to the in-plane resistance R<sub>ab</sub> commonly measured in a four terminal measurement. Our previous work on the in-plane resistivity demonstrated that for R<sub>top</sub>(T), above the melting temperature in the liquid state, there are two distinct parts where lnR<sub>top</sub> shows linear dependence upon 1/T with different slopes, which is a typical characteristic for thermally activated flux flow (TAFF) behavior. We have identified the temperature T<sub>d</sub> where the crossover between the two TAFF regions occurs as the decoupling transition through an activation energy analysis. The decoupling temperature is described by
$$B_{cr}=\frac{\mathrm{\Phi }_0^3}{16\pi ^3k_BTse\lambda _{ab}^2\gamma ^2}\text{,}$$
(1)
where s is the interlayer spacing between the neighboring CuO planes. With a $`\gamma `$ value of 12.4 a satisfied fit is obtained by Eq. (1) for the decoupling line shown as the dashed line in Fig. 3.
We notice that the temperature where R<sub>bot</sub> shows a maximum corresponds to exactly the same temperature where a crossover between two different TAFF regions occurs as we have identified before. Such a peak resembles that of the out-of-plane resistance R<sub>c</sub> often appears in single crystals with large anisotropy such as Bi2212. However it can be seldom seen in single crystal with moderate anisotropy such as fully oxidized Y123. Although there is no consensus on the occurrence of the peak in R<sub>bot</sub> and R<sub>c</sub> yet, one possible reason for it could be the tunneling of vortices over the c-axis barrier in the liquid state following by the complete coupling of CuO planes. Recently, there have been considerable efforts devoted to the study of interlayer charge dynamics by far infrared reflectivity. Studies on the interlayer charge dynamics with the electric field component E polarized along the c-axis reveals that above a characteristic temperature T<sub>pl</sub>, the superconductor behaves like a poor metal or an ionic insulator. The carriers are confined inside the CuO planes and the c-axis resistivity is controlled by the interlayer tunneling. The c-axis coherence is signaled by the appearance of a reflectivity edge associated with the Josephson plasma from the intrinsic Josephson junctions. Above T<sub>pl</sub>, the electron transport along the c-axis is incoherent. Below T<sub>pl</sub>, coherence is built up for the interlayer carrier transport. Therefore, given the high transition temperature and weak interlayer magnetic coupling between the pancake vortices in the neighboring CuO planes, T<sub>d</sub> would quite possibly correspond to T<sub>pl</sub> (the Josephson plasma frequency) at which the phase coherence and thus the superconductivity is built up for the entire system along the c-axis. The suppression of the onset temperature for c-axis superconductivity by an applied magnetic field is due to the breakdown of the phase coherence along the c-axis by a magnetic field. This is consistent with the model proposed by Briceño et al. and later developed by Suzuki et al. who suggested that phase fluctuation induced dissipation controlled the dissipation and could be well described by the Ambegaokar-Baratoff theory.
Although the picture of vortex melting has been widely accepted, the mechanism of the vortex melting has yet not well understood. Recently, Nonomura et al. carried out Monte Carlo simulations of the three dimensional frustrated XY model, in which the melting temperature was determined by the helicity modulus along the c-axis. They found that the melting transition is propagated through vortex entanglement. At the melting temperature, the percentage of entangled vortices abruptly changes. Upon increasing the temperature, the entanglement length becomes smaller than the sample thickness and decreases rapidly. The dissipation is governed by the cutting of vortices above the melting temperature. As a consequence of the entanglement mechanism of the FLL melting, T<sub>m</sub> is scaled by the inverse of the system size along the c axis. Our results can be explained by their simulation results very well since we indeed observed a melting transition into 3D line liquid with its correlation length l<sub>c</sub> smaller than the sample thickness t.
The vortex melting scenario has been questioned by Moore who argued that the melting transition could just be a second order crossover of the vortex matter from a three dimensional behavior to a two dimensional one when the l<sub>c</sub> in the vortex liquid becomes compatible to the sample thickness t. Because l<sub>c</sub> grows very rapidly as the temperature is lowered, the crossover region appears narrow enough to be misinterpreted as a first order melting transition. For a cutting length of the sample thickness t, the cutting temperature can be calculated according to
$$B_{cut}=\frac{\mathrm{\Phi }_0\epsilon _0}{\gamma ^2Lk_BT}\text{,}$$
(2)
by substituting L with t$``$0.01 mm. Here $`\varphi _0`$=2.07$`\times `$10<sup>-7</sup> Gauss cm<sup>-2</sup> is the flux quantum, $`\epsilon _0=\mathrm{\Phi }_0^2\mathrm{ln}\kappa /16\pi ^2\lambda _{ab}^2(T)`$, $`\kappa `$ being the Ginzburg-Landau parameter and $`\lambda _{ab}(0)2000`$Å, the penetration depth in the ab plane. The obtained vortex cutting line lies slightly below the melting line B<sub>m</sub>(T). Although our result can also be qualitatively explained with this argument, it is not clear why such a sharp crossover is not observed in thin films (strong pinning disorder) since in any sample with a finite thickness larger than the distance between neighboring CuO planes, a crossover temperature is expected. One reason could be that the relatively smaller numbers of pinning centers in the single crystals (weak pinning) than in the thin films result in a larger Larkin correlation length L<sub>p</sub> inside which an ordered VL can persist. In the case of thin films, the transverse correlation length is smaller than the average lattice spacing, only a second order glass transition can be observed. This in turn suggests that a regular VL with a relatively large radius is a necessity for the occurrence of a melting transition. This is consistent with the fact that artificially introduced pinning centers will induce a small Larkin length and smear out the melting transition. Further work on the study of the relatiohship between pinning strength and Larkin length need to be done in order to clarify this point.
Our results suggest the following physical scenario: in a clean Y124 single crystal, below the melting temperature T<sub>m</sub>, the vortices form a rigid VL with l$`{}_{c}{}^{}t`$. As the superconductor is warmed up, at T<sub>m</sub>, VL melting occurs. At this stage, the vortices lose their translational coherence while keeping the correlations along the c-axis. Thus the vortices in the temperature region just above the melting transition are in a 3D line liquid state. At this temperature region, the vortex form a kind of entangled vortices with their entanglement length l<sub>c</sub> smaller than the sample thickness t but much larger than the distance s between the CuO planes. The dissipation is governed by the vortex cutting and recombination. As the temperature increases, the correlation length in the c-axis becomes smaller and smaller and the difference between R<sub>top</sub> and R<sub>bot</sub> gets larger and larger. Finally at the temperature T<sub>d</sub>, inter-plane decoupling occurs (l<sub>c</sub> $`<`$s) and the vortices lose the coherence in the c direction and the vortices dissolve into 2D pancake vortices that move independently in the individual CuO planes.
## IV Conclusions
In conclusion, we have observed the melting transition in Y124 single crystals. The vortices following the melting transition are in the 3D line liquid state with a correlation length in the c- direction l<sub>c</sub> smaller than the sample thickness but larger than the distance between the CuO planes of Y124. The temperature at which a resistive peak exists at R<sub>bot</sub> is found to correspond to the interlayer decoupling transition temperature. Above the decoupling temperature, the vortices lose their coherence in the c-direction and the dissipation would be governed by tunneling of 2D vortices across the CuO planes. Our results also suggest that the interlayer decoupling transition is a continuous crossover rather than a sharp transition.
## V Acknowledgments
We thank E. Rossel and P. Wagner for their help during the measurements and Y. Bruyneseraede for helpful discussions. This research has been supported by the ESF Programme ”VORTEX”, the Belgian IUAP and Flemish GOA and FWO Programmes.
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# X–ray Emission of Mkn 421: New Clues From Its Spectral Evolution: II. Spectral Analysis and Physical Constraints
## 1. Introduction
Blazars are radio–loud AGNs characterized by strong variability, large and variable polarization, and high luminosity. Radio spectra smoothly join the infrared-optical-UV ones. These properties are successfully interpreted in terms of synchrotron radiation produced in relativistic jets and beamed into our direction due to plasma moving relativistically close to the line of sight (e.g. Urry & Padovani up95, 1995). Many blazars are also strong and variable sources of GeV $`\gamma `$–rays, and in a few objects, the spectrum extends up to TeV energies. The hard X to $`\gamma `$–ray radiation forms a separate spectral component, with the luminosity peak located in the MeV–TeV range. The emission up to X–rays is thought to be due to synchrotron radiation from high energy electrons in the jet, while it is likely that $`\gamma `$-rays derive from the same electrons via inverse Compton (IC) scattering of soft (IR–UV) photons –synchrotron or ambient soft photons (e.g. Sikora, Begelman & Rees sbr94, 1994; Ghisellini & Madau gg\_madau\_96, 1996; Ghisellini et al. gg\_sed98, 1998). The contributions of these two mechanisms characterize the average blazar spectral energy distribution (SED), which typically shows two broad peaks in a $`\nu F_\nu `$ representation (e.g. Fossati et al. 1998a, ). The energies at which the peaks occur and their relative intensity provide a powerful diagnostic tool to investigate the properties of the emitting plasma, such as electron energies and magnetic field (e.g. Ghisellini et al. gg\_sed98, 1998). In X–ray bright BL Lacs (HBL, from High-energy-peak-BL Lacs, Padovani & Giommi pg95, 1995) the synchrotron maximum occurs in the soft-X–ray band.
Variability studies constitute the most effective means to constrain the emission mechanisms taking place in these sources as well as the geometry and modality of the energy dissipation.
The quality and amount of X–ray data on the brightest sources start to allow thorough temporal analysis as function of energy and the characterization of the spectral evolution with good temporal resolution.
Mkn 421 is the brightest HBL at X–ray and UV wavelengths and thus it is the best available target to study in detail the properties of the variability of the highest frequency portion of the synchrotron component, which traces the changes in the energy range of the electron distribution which is most critically affected by the details of the acceleration and cooling processes.
This paper is the second of two, which present the uniform analysis of the X–ray variability and spectral properties from BeppoSAX observations of Mkn 421 performed in 1997 and 1998. In Paper I (Fossati et al. fossatiI, 2000) we presented the data reduction and the timing analysis of the data, which revealed a remarkably complex phenomenology. The study of the characteristics of the flux variability in different energy bands shows that significant spectral variability is accompanying the pronounced changes in brightness. In particular, a more detailed analysis of the remarkable flare observed in 1998 revealed that: i) the medium energy X–rays lag the soft ones, ii) the post–flare evolution is achromatic, and iii) the light curve is symmetric in the softest X–ray band, and it becomes increasingly asymmetric at higher energies, with the decay being progressively slower that the rise.
The general guidelines which we followed for the data reduction and filtering are described in Paper I (in particular in §2 and §3.2). Here we will only report on details of the treatment of the data specific to the spectral fitting.
The paper is organized as follows. In Sections §2 and §3 we briefly summarize the basic information on BeppoSAX and the 1997 and 1998 observations. The results and discussion relative to the spectral analysis are the content of §4 and §5. In particular, the observed variability behavior strongly constrains any possible time dependent particle acceleration prescription. We will therefore consider these results together with those of the temporal analysis, discuss which constraints are provided to current models and present a possible scenario to interpret the complex spectral and temporal findings (§5.4). Finally, we draw our conclusions in §6.
## 2. BeppoSAX overview
For an exhaustive description of the Italian/Dutch BeppoSAX mission we refer to Boella et al. (boella97, 1997) and references therein. The results discussed in this paper are based on the data obtained with the Low and Medium Energy Concentrator Spectrometers (LECS and MECS) and the Phoswich Detector System (PDS). The LECS and MECS have imaging capabilities in the 0.1–10 keV and 1.3–10 keV energy band, respectively, with energy resolution of 8% at 6 keV. The PDS covers the range 13–300 keV.
The present analysis is based on the SAXDAS linearized event files for the LECS and the MECS experiments, together with appropriate background event files, as produced at the BeppoSAX Science Data Center (rev 0.2, 1.1 and 2.0). The PDS data reduction was performed using the XAS software (Chiappetti & Dal Fiume chiappetti\_dalfiume, 1997) according to the procedure described in Chiappetti et al. (chiappetti\_2155, 1999).
## 3. Observations
Mkn 421 has been observed by BeppoSAX in the springs of 1997 and 1998. For reference, the BeppoSAX light curves for the 4–6 keV band are reported in Fig. 1. Most of the forthcoming analysis is focused on the spectral variability observed during the flare of 1998 April 21<sup>st</sup>, clearly visible in the bottom panel. The journal of observations is given in Table 1 of Paper I.
## 4. Spectral Analysis
In order to unveil the spectral variability during and after the flare, i.e. perform a time resolved spectral analysis, we subdivided the whole dataset in sub–intervals corresponding to single BeppoSAX orbits (42 for the 1997 dataset, 16 on 1998 April 21<sup>st</sup>, 19 on 1998 April 23<sup>rd</sup>) or a grouping of them sufficient to reach the same statistics (these are reported in Table 1, together with the median time \[UTC\] of each sub–interval).
LECS and MECS spectra have been accumulated as described in Paper I. LECS data have been considered only in the range 0.12–3 keV due to calibration problems (a spurious hardening) at higher energies (Guainazzi 1997, private comm.)<sup>1</sup><sup>1</sup>1A word of caution is necessary about localized features probably related with calibration problems yet to be solved (e.g. E$``$ 0.29 keV, $``$ 2 keV, in correspondence to the Carbon edge and the Gold features due to the optics, respectively).
The nominal full resolution spectra (i.e. channels #11–285 for 0.1–3 keV in the LECS, and #36–220 for 1.6–10 keV in the MECS) have been then rebinned using the grouping templates available at BeppoSAX–SDC<sup>2</sup><sup>2</sup>2The energy resolution of the LECS and MECS detectors is lower than the energy spacing between instrumental calibrated channels ($``$ 10 eV for LECS and 45 eV for MECS). Therefore it is compelling to rebin the spectra even in cases (like ours) where the statistics in each low energy channel is good, in order to avoid to overweight (as $`\chi _\mathrm{a}^2+\chi _\mathrm{b}^2>\chi _{\mathrm{a}+\mathrm{b}}^2`$) information that is not truly independent. BeppoSAX–SDC templates are downloadable at ftp://www.sdc.asi.it/pub/sax/cal/responses/grouping/. Due to the very steep spectral shape, in order to maintain a good statistics in each new bin, we altered the templates for MECS above $`7`$ keV, increasing the grouping of the original PI channels. The background has been evaluated from the blank fields provided by the BeppoSAX–SDC<sup>3</sup><sup>3</sup>3Blank fields event files were accumulated on five different pointings of empty fields and are available at the anonymous ftp: ftp://www.sdc.asi.it/pub/sax/cal/bgd/., using an extraction region similar in size and position to the source extraction region. The Nov. 1998 release of public calibration files, matrices and effective areas was used.
For the PDS spectra we applied the improved screening, implementing the temperature and energy dependence of the pulse rise time (so–called PSA correction method).
Due to a slight mis–match in the cross–calibration among the different detectors, it has been necessary to include in the fitting models multiplicative factors of the normalization (LECS/MECS and PDS/MECS ratios). The correct absolute flux normalization is provided by the MECS (the agreement between MECS units is within the 2–3% limit of the systematics). The expected value of these constant factors is now well known and does not constitute a major/additional source of uncertainty: according to Fiore, Guainazzi & Grandi (cookbook, 1999) the acceptable range for LECS/MECS is 0.7–1.0, with some dependence on the source position in the detectors, while the PDS/MECS ratio can be constrained between 0.77 and 0.95. The latter parameter is indeed crucial in sources like Mkn 421 where there is a significant possibility that a different component arises in the PDS range, whose recognition critically depends on the capability/sensitivity to reject the hypothesis that the PDS counts can be accounted for by the extrapolation of the LECS–MECS spectrum (see Section 4.3.3).
### 4.1. Single and Broken Power Laws
We do not discuss in detail any power law spectral models. For each LECS$`+`$MECS spectrum we fit the data with the single and the broken power law models<sup>4</sup><sup>4</sup>4In Appendix C we enclose a Table reporting the best fit values for the the relevant parameters for the broken power law model., both with free and fixed<sup>5</sup><sup>5</sup>5In the direction of Mkn 421 the Galactic equivalent absorbing column is N$`{}_{\mathrm{H}}{}^{}=(1.61\pm 0.1)\times 10^{20}`$ cm<sup>-2</sup> (Lockman & Savage lockman\_savage95, 1995), with the estimated uncertainty for high Galactic latitudes. absorbing column density.
Here we just show that they are an inadequate description of the downward curved spectra<sup>6</sup><sup>6</sup>6In Fossati et al. (1998b, ) we presented a preliminary analysis of 1997 data, where LECS and MECS spectra have been fit with a variety of combinations of single and broken power law models to better describe the spectral curvature. We refer the interested reader to that work.. To illustrate the discrepancy, we took the data/model ratio for the best fit continuum model, for each of the individual LECS$`+`$MECS spectra (16 for 1997, 12 for 1998), and we summed all of them, propagating the errors accordingly. The results are shown in Figure 2 for the 1998 data only, while in Table 2 we report the total $`\chi ^2`$ value for each trial model, together with the (simple) average value of N<sub>H</sub> for fits with free absorption.
Neither the single nor the broken power laws provide an adequate representation of the data, not even when the value of the absorbing column density is left free to vary<sup>7</sup><sup>7</sup>7Although the plot shown in Figure 2d could appear as a good fit, the significance of the small deviations above a few keV is indeed high, yielding an unsatisfactory fit. –a common workaround used when the spectrum shows either a soft excess or a marked steepening toward higher energies (e.g. Takahashi et al. takahashi96, 1996).
Indeed the description of the curvature by means of soft X–ray absorption is not only un–physical, but it hinders the possibility of extracting all the information from the data, by leaving as meaningful parameter only the higher energy spectral index and accounting for all the other effects by N<sub>H</sub>. Furthermore this yields spectral indices that are only rough estimates, as it is clear that the best fit single power law does not describe the observed data in any (even narrow) energy range (see Figure 2a,b). Finally, in the case of Mkn 421 there is no reason to postulate any intrinsic absorbing component responsible for the observed spectral curvature (see also §4.3.2).
### 4.2. The Curved Model
Motivated by the failure of the simplest power law models, we developed a spectral model which is intrinsically curved<sup>8</sup><sup>8</sup>8In Appendix A we discuss an important caveat concerning the definition of curved models., with the aim of extracting all the information contained in the data, and in particular estimate the position of the peak of the synchrotron component, one of the crucial quantities in blazars modeling.
A spectral model providing an improved description of the continuum shape also allows a more sensitive study of discrete spectral features such as absorption edges, whose presence has been long searched in the soft X–ray blazars spectra (e.g. Canizares & Kruper canizares84, 1984; Sambruna et al. sambruna\_1426, 1997; Sambruna & Mushotzky sambruna\_0548, 1998).
We started from the following general description of a continuously curved shape (see e.g. Inoue & Takahara inoue\_takahara96, 1996, Tavecchio, Maraschi & Ghisellini tavecchio\_etal98, 1998):
$$F(E)=KE^\alpha _{\mathrm{}}\left(1+\left(\frac{E}{E_\mathrm{B}}\right)^f\right)^{\frac{\alpha _{\mathrm{}}\alpha _+\mathrm{}}{f}}$$
(1)
where $`\alpha _{\mathrm{}}`$ and $`\alpha _+\mathrm{}`$ are the asymptotic values of spectral indices for $`EE_\mathrm{B}`$ and $`EE_\mathrm{B}`$ respectively, while $`E_\mathrm{B}`$ and $`f`$ determine the scale length of the curvature.
We re–expressed this function in terms of the spectral indices at finite values of $`E`$, which characterize the local shape of the spectrum, instead of the asymptotic ones, which do not have any direct reference to the observed portion of the spectrum.
The spectral model is then expressed in a form such that the available parameters are $`(E_1,\alpha _1,E_2,\alpha _2,E_\mathrm{B},f)`$ instead of ($`\alpha _{\mathrm{}}`$, $`\alpha _+\mathrm{}`$, $`E_\mathrm{B}`$, $`f`$) (for more details see Appendix B). As we have two extra parameters, for a meaningful use of this spectral description we have to fix one for each of the pairs $`(E_1,\alpha _1)`$ and $`(E_2,\alpha _2)`$. Eventually this degeneracy turns out to be a powerful property of this model, because it allows us to derive the spectral index at selected energies (setting $`E_\mathrm{i}`$ at the preferred values), or even more interestingly to estimate the energy at which a certain spectral index is obtained (setting $`\alpha _\mathrm{i}`$ at the desired value). The most relevant example of this latter possibility is the determination of the position of the peak (as seen in $`\nu F_\nu `$ representation) of the synchrotron component $`E_{\mathrm{peak}}`$ –if it falls within the observed energy band– and the estimate of the associated error. This can be obtained by setting one spectral index, i.e. $`\alpha _1=1`$, and leaving the corresponding energy $`E_1`$ free to vary in the fit: the best fit value of $`E(\alpha =1)`$ gives $`E_{\mathrm{peak}}`$.
### 4.3. Results
Most of the spectral analysis with the curved model has been performed by keeping the absorbing column fixed to the Galactic value, to remove the degeneracy between the effects of a variable absorption and intrinsic spectral curvature. We will briefly comment on the fits with a variable absorbing column in the next section.
In order to determine the values for the really interesting parameters, we tried a set of values for the parameter $`f`$, ranging from 0.5 to 3 (with a step of 0.5). All the other parameters were left free to vary during these step–fits. We then selected the value of $`f`$ yielding the minimum total $`\chi ^2`$. The best fit value for 1997 is $`f=1`$, while for 1998 $`f=2`$, as expected from the stronger curvature of the spectra. In Table 1 we report the spectral parameters, together with the 0.2–1, 2–10 and 0.1–10 keV fluxes (de–reddened) and $`\chi ^2`$ values, for each of the time sliced spectra. The global Data/Model ratio is shown in Figure 2e, and the total $`\chi ^2`$ is reported in Table 2. The curved model is strongly preferred also from the purely statistical standpoint: in fact the model with Galactic N<sub>H</sub> has a significantly better $`\chi ^2`$ than the broken power law with free N<sub>H</sub>, although it has one less adjustable parameter (i.e. one more degree of freedom).
Each spectrum has been fitted a few times in order to derive spectral indices at several energies (0.5, 1, 5, 10 keV), and also E<sub>peak</sub>. For consistency we checked each time that not only the value of the $`\chi ^2`$ remained the same, but also all the other untouched parameters took the same value and confidence intervals. The fits are indeed very robust in this sense.
The main new result is that we were able to determine the energy of the peak of the synchrotron component, and to assign an error to it. This has been possible with reasonable accuracy for both 1997 and 1998 spectra, with a couple of cases yielding only an upper limit for the peak energy<sup>9</sup><sup>9</sup>9These are the cases “35–37” and “41–42” of 1997. Actually the formal fitting yields a confidence interval, but with the lower extreme falling outside of the observed energy range..
An example of the remarkable spectral evolution during the 1998 flare is shown by the deconvolved spectra in Figure 3, which also illustrates the strong convexity of the spectrum and the well defined peak.
In 1997 the source was in a lower brightness state, with an average X–ray spectrum softer (at all energies $`\mathrm{\Delta }\alpha _{9798}0.4`$), and a peak energy 0.5 keV lower.
More globally, the analysis shows that there is a clear relation between the flux variability and the spectral parameters, both in 1998 and, albeit less strikingly, in 1997. Particularly important is the correlation between changes in the brightness (even the small ones) and shifts of the peak position, as the latter carries direct information on the source physical properties. In Figure 4a,b, $`\alpha `$ at 5 keV and E<sub>peak</sub> for 1997 and 1998 are plotted versus the 0.1–10 keV flux. The source reveals a coherent spectral behavior between 1997 and 1998 and through a large flux variability (a factor 5 in the 0.1–10 keV band). Both the peak energies and the spectral indices show a tight relation with flux, the latter being described by E$`{}_{\mathrm{peak}}{}^{}\mathrm{F}^ϵ`$, with $`ϵ=0.55\pm 0.05`$.
#### 4.3.1 Hard Lag in 1998 spectra
One further important finding of the time resolved spectral study is the signature of the hard lag which strengthens what already inferred from the timing analysis (see Paper I). In fact the local spectral slope at different energies indicates that the flare starts in soft X–rays and then extends to higher energies.
A blow up of the 1998 flare interval is shown in Figure 5: spectral indices at 1 and 5 keV, and E<sub>peak</sub> are plotted versus time in the left hand plots, and versus flux in the right hand ones. The synchrotron peak shifts toward higher energies during the rise and then decreases as soon as the flare is over. The spectral index at 1 keV reflects exactly the same behavior, as expected being computed at the energy around which the peak is moving (and in fact it moves in a narrow range around $`\alpha =1`$). On the contrary, the spectral shape at 5 keV does not vary until a few ks after the peak, and only then –while the flux is decaying and the peak is already receding– there is a response with a significant hardening of the spectrum. The right side diagrams illustrate this behavior directly in terms of flux: the evolution of the peak energy and the 1 keV spectral shape follow closely that of the flare, while at 5 keV the spectrum does not change shape until after the peak.
Thus the spectral evolution at higher energies develops during the decay phase of the flare, tracing a loop in the $`\alpha `$ vs. Flux diagram, with a significant hardening of the spectrum. Due to the presence of the hard lag, the loop is traced in the opposite way with respect the evolution that is commonly observed during flares in HBL, that is a spectral hardening during the rise of the outburst, followed by a softening during the decay (see however Sembay et al. sembay\_2155, 1993, Catanese & Sambruna catanese\_sambruna\_00, 2000).
This behavior –usually accompanied by a soft lag– produces a clockwise loop in $`\alpha `$ vs. Flux.
#### 4.3.2 Absorption Features: N<sub>H</sub> and edges
We also fit the 1997 and 1998 spectra with the curved spectral model leaving the N<sub>H</sub> free to vary. The only constraint that we imposed is on the exponent $`f`$, which is set at the best fit values of $`f=1`$ and $`f=2`$ for 1997 and 1998, respectively. The improvement of the global fit<sup>10</sup><sup>10</sup>10$`\mathrm{\Delta }\chi ^2=17.9`$ for 16 additional d.o.f. for 1997 and $`\mathrm{\Delta }\chi ^2=15.9`$ for 12 additional d.o.f. for 1998, corresponding to an average $`\mathrm{\Delta }\chi ^21.2`$ per spectrum for a change in d.o.f. from 57 to 56. is not statistically significant. The average best fit absorbing column is indeed very close to and consistent with the Galactic one, leading us to conclude that there is no requirement whatsoever for an additional absorbing column to model the data.
We searched for the presence of absorption edges at low energies (E $``$ 1 keV), both in 1997 and 1998 data. We added to the curved continuum an absorption edge, with a first guess fit at E$`0.40.5`$ keV, as even in the case of the best curved model with free N<sub>H</sub> there are residues around this energy (see Figure 2f). The most delicate issue is the calibration problem related with the instrumental Carbon absorption edge at 0.29 keV. Even the smallest uncertainties ($`12`$%) in the knowledge of the instrument response around this energy can induce significant spurious features, typically at slightly higher energy (see Fossati & Haardt fg\_fh\_lecs, 1997).
There is not convincing evidence for discrete absorption features. In 5/28 cases the fit routine has not found an edge, while there are indeed a few spectra (8/28) in which the fit definitely put the edge at 0.27–0.30 keV, to account for systematic deviation due to calibration uncertainties. There are then 15/28 cases where the edge energy is around 0.4–0.6 keV, but again only in a handful of cases the best fit energy E<sub>edge</sub> is significantly different from 0.29 keV.
Summarizing, there are only three individual spectra (out of 28) for which the inclusion of the edge is formally statistically significant ($`>`$ 95%), according to the F–test. These three cases do not stand out in the sample for any other property, such as brightness or curvature of the continuum. Best fit energies and optical depths are typically E$`{}_{\mathrm{edge}}{}^{}500_{100}^{+50}`$ eV, and $`\tau 0.4_{0.2}^{+0.30.6}`$. However, considering that the sample comprises 28 spectra, with the significance threshold set at 95%, three positive detections can not be regarded as a compelling evidence for the presence of absorption edges.
#### 4.3.3 PDS data
Covering the range above $`12`$ keV, the PDS instrument could detect the IC component, which might start to dominate the emission in this band. It has always proven to be very difficult to detect/constrain this component, although it is expected to have a very hard spectrum, that should make it easier to disentangle from the very steep tail of the synchrotron component. There are a few objects, the so–called intermediate BL Lacs, in which the cross–over between synchrotron and IC occurs in the 0.1–10 keV band, for which it has then been possible to clearly detect the hard IC power law, with typical $`\alpha 0.30.7`$ (e.g. Tagliaferri et al. tagliaferri\_on231, 2000 for ON 231; Giommi et al. giommi\_0716, 1999 for S5 0716$`+`$714; Sambruna et al. sambruna\_bllac\_99, 1999 for BL Lac). The good sensitivity of the PDS yielded a convincing evidence of IC emission in the case of PKS 2155$``$304 (Giommi et al. giommi\_2155, 1998). The preliminary analysis of the Mkn 421 1997 data (Fossati et al. 1998b, ) suggested the presence of a hard component. We therefore focused on the search for the IC component for both of 1997 and 1998 PDS datasets.
For 1998, when the source was brighter, we accumulated PDS spectra according to the same partition of the light curves used for the 0.1–10 keV data. In most of the individual spectra there is a significant ($``$3-$`\sigma `$) detection only up to 40 keV. The light curve for the 12–26 keV band is shown in Figure 3 of Paper I.
We then restricted the analysis to only one spectrum for 1997, integrating over the whole campaign, and two for 1998, one for April 21<sup>st</sup>, and one for April 23<sup>rd</sup>. The integration times are 54.8, 21.1 and 25.2 ks respectively (see Paper I). We grouped the data in 6 channels spectra with boundaries at 12, 18, 27, 40, 60, 90 and 130 keV.
In the 1997 spectrum there is a positive detection up to 90 keV. For April 21<sup>st</sup> –corresponding to the synchrotron flare– there is a strong signal up to 60 keV, while, surprisingly, in the April 23<sup>rd</sup> spectrum there is a $``$3-$`\sigma `$ detection in each individual bin up to 130 keV. From 12 to 40 keV the April 21<sup>st</sup> count rates are significantly higher than those of April 23<sup>rd</sup>, but while the former continue to steeply decrease with increasing energy, the latter flatten above 40 keV, with a cross over in the 40–60 keV bin. As the integration times for these two spectra are similar (and in fact statistical uncertainties on the count rate are of the same order), this cannot account for this difference<sup>11</sup><sup>11</sup>11Also, we are not aware of any intra–observation “background fluctuations” to be taken into account in the standard PDS data reduction (which has been performed by following carefully the prescriptions given by the instrument team, for all the details refer to Chiappetti et al. chiappetti\_2155, 1999), nor of any long term “background fluctuations” that could be responsible for the observed variation. No such thing was reported neither by the PDS instrument team, nor by the BeppoSAX Science Data Center..
The results of single power law fits are reported in Table 3: clearly the April 23<sup>rd</sup> spectrum is significantly harder. Moreover, the spectral indices for the first two cases are consistent with those derived from the LECS and MECS datasets and the general picture of a continuously steepening spectrum, while this does not hold for the third case.
While this is already interesting, there is one additional, somehow unexpected, finding: apparently the April 23<sup>rd</sup> harder spectrum cannot be simply interpreted as the detection of the IC component when the synchrotron tail recedes. In fact a hard power law at this flux level –if present– would have been detected also on April 21<sup>st</sup>, while the count rates in the higher energy bins for April 21<sup>st</sup> set a tight upper limit. This behavior might be ascribed to either a different independent origin or a delayed response of the IC component with respect to the synchrotron one (indeed some evidence of a possible delayed decay of the PDS with respect to the MECS emission can be seen in the light curves).
Finally we tried to quantify the presence of the IC power law in the three datasets by fitting MECS and PDS data together. We proceeded in the following way:
* We accumulated LECS and MECS spectra over the same intervals of the PDS ones, i.e. one cumulative spectrum for each day of 1998, and one for the whole 1997.
* We fit the LECS$`+`$MECS spectra in order to constrain the average synchrotron spectrum.
* We then fixed all the parameters of the curved model to the LECS$`+`$MECS best fit values, and added to the model a hard power law component<sup>12</sup><sup>12</sup>12We used the pegpwrlw (power law with pegged normalization) model of XSPEC that provides a robust measure of the normalization. In fact this is defined over a selectable energy band, instead of the monochromatic value at 1 keV that can be strongly correlated with other parameters if 1 keV does not fall close to the logarithmic median of the band spanned by the data. with fixed spectral index (we tried the values $`\alpha _{\mathrm{IC}}=0.4,0.5,0.6,0.7`$). We also added an exponential cut–off to the synchrotron component setting its two parameters, i.e. cut–off and e–folding energies, to 10 and 40 times the best fit value of E<sub>peak</sub>.
* We used this partially constrained model to fit the MECS (above E$`5`$ keV) and PDS spectra. The two 1998 datasets have been fit jointly with PDS/MECS intercalibration factor free to vary, but set to be the same for both, and of course independent normalization.
The results are reported in Table 4, for the $`\alpha _{\mathrm{IC}}=0.5`$ case (the values relative to the other $`\alpha _{\mathrm{IC}}`$ are not significantly different). The best fit PDS/MECS normalization ratio is 0.88.
There seems to be a significant change in the flux level of the expected IC component between 1998 April 21<sup>st</sup> and 23<sup>rd</sup>. The best fit value for the normalization for the first day is zero, while for the second day a non–zero IC flux is definitely required. In order to test the likelihood of the apparent variation, we fit the data constraining the normalization to be the same for the two datasets (third case in Table 4) and the F–test probability to obtain by chance the observed change in $`\chi ^2`$ is $``$ 0.01 ($`\mathrm{}\chi ^2=6.08`$, d.o.f. from 30 to 31).
The basic result for 1997 is that the IC component is definitely detected. Its brightness level is comparable with the average of 1998, while the synchrotron component shows a significant variation.
It should be noted that the synchrotron fluxes reported in Table 4 are not the best reference because at those energies, well above the synchrotron peak, the spectrum is very steep and even the smallest shift of the peak energy translates in a big change in the flux. Nevertheless even when considering the energy range 0.1–10 keV, which includes the synchrotron peak, the variation between 1997 and 1998 is of the order of a factor of 2–3 (e.g. Figure 5). The sampling of the 1997 light curve is not good enough to enable us to investigate any possible variation of the PDS spectrum with the source brightness.
## 5. Discussion
Before focusing on the modeling of the temporal and spectral behavior illustrated in this paper and in Paper I, let us consider a few issues which we consider of particular relevance for the interpretation of the origin of blazar variability.
### 5.1. Spectral Variability, Synchrotron Peak and IC component
Soft and hard X–ray bands show a different behavior. This might be attributed to the contribution from both the synchrotron and inverse Compton components at these energies. As electrons emitting X–rays through the two processes would have different energies, the two components are not expected to vary on the same timescales and in phase. Indeed on one hand the constraints provided by the analysis of the PDS data suggest that the synchrotron component constitutes the dominant contribution to the flux in the 12–26 keV band on 1998 April 21<sup>st</sup> (during the flare) with a light curve somewhat different from that of the softer synchrotron X–rays. On the other hand, the evidence supporting the possibility of a substantial variation in the IC component during the 1998 observations is extremely interesting. We can rule out the possibility that the flux level of the hard component is approximately constant between the two observations of 1998 (with a 99% confidence), and thus that its detectability only depends on the the brightness level of the tail of the synchrotron emission.
### 5.2. The Flaring + Steady Components Hypothesis
A further interesting issue raised in Paper I in connection with the temporal analysis is the presence and role of quasi–stationary emission. Our results support the view that the short–timescale, large–amplitude variability events could be attributed to the development of new individual flaring components, giving rise to a spectrum out–shining a more slowly variable contribution.
Indeed there is growing evidence that the overall blazar emission comprises a component that is (possibly) changing only on long timescales and that does not take part in the flare events (see for instance Mkn 501, Pian et al. pian\_mkn501\_98, 1998, and Pian et al. pian\_3c279, 1999 for the case of the UV variability of 3C 279).
The deconvolution of the SED into different contributions would allow the study of the evolution of the flaring component possibly clarifying the nature of the dissipation events occurring in relativistic jets –in particular the particle acceleration mechanism– and the modality and temporal characteristics of the initial release of plasma and energy through these collimated structures.
The deconvolution is however still difficult to achieve with the available temporal and spectral information. In paper I we discussed the results of our attempt to measure directly the relative contributions of what we identify as steady and flaring components (parameter $``$), from the characteristics of the flare decay. The result is very interesting, although it does not provide a useful handle on the evolving spectral properties of the flare itself. A more feasible approach is to assume specific forms of the flare evolution and test them against the data. One of the simplest working hypothesis is to assume that the flare evolution can be reproduced by the time dependent flux and energy shifts of a peaked spectral component with fixed spectral shape (e.g. Krawczynski et al. krawczynski\_mkn501\_2000, 2000). In this context it is relevant to remind that the maximum of the synchrotron emission corresponds to the energy at which most of the energy in particles is located. Its value and temporal behavior thus trace the evolution of the bulk of the energy deposition into particles by the dissipation mechanism.
It is therefore possible that the variability characteristics (both temporal and spectral) might chiefly depend on the (particles/photons) energy relative to the synchrotron peak. In particular, the results on T<sub>short</sub> (Paper I) can be either related to the dominance of the light crossing time over the cooling timescales, or alternatively to the different position of the sampled energies with respect to the synchrotron peak as above the peak the timescales are shorter. In any case the fact that timescales are similar at the different energies possibly indicates the evolution of a (observed) fixed–shape flaring component.
### 5.3. Sign of the lag
One of the main new results of the temporal analysis presented in Paper I is the significant detection of a hard lag, opposite to the behavior commonly detected in HBL. Interestingly, a recent study by Zhang et al. (zhang\_2155, 1999) on the source PKS 2155–304 has shown the presence of an apparent inverse trend between the source brightness level and the duration of the (soft) lag. And intriguingly, the 1998 flare represents the brightest state ever observed from Mkn 421, thus possibly suggesting that an extrapolation of the above phenomenological trend might even account for a hard lag. This possibility is currently under study (Zhang et al., in preparation).
Again the relative contribution of two components might be responsible for this trend. During the most intense events the flaring component would completely dominate over the quasi–stationary one, progressively shifting to higher energies (hard lag), while in weaker flares the varying component would exceed the steady one only at energies higher than the peak, where the spectrum steepens, producing an observed soft lag behavior.
However, one should also consider that in different brightness states the different observed energies correspond to different positions with respect to E<sub>peak</sub> (e.g. comparing the 1997 and 1998 datasets we showed that its value was typically a factor 2 different), requiring a more subtle analysis before deriving inferences on the lag–brightness relationship.
### 5.4. Physical Interpretation: the signature of particle acceleration
Let us now focus on the interpretation of the two main and robust results of this work, namely the hard lag and the evolution of the synchrotron peak.
The occurrence of the peak at different times for different energies is most likely related to the particle acceleration/heating process. Although models to reproduce the temporal evolution of a spectral distribution have been developed (e.g. Chiaberge & Ghisellini chiab\_gg, 1999, Georganopoulos & Marscher markos\_98, 1996; Kirk, Rieger & Mastichiadis kirk\_etal\_98, 1998), they mostly do not consider the role of particle acceleration (see however Kirk et al. kirk\_etal\_98, 1998).
In order to account for the above results, we thus introduced a (parametric) acceleration term in the particle kinetic equation within of the time dependent model studied by Chiaberge & Ghisellini (chiab\_gg, 1999). In particular, this model takes into account the cooling and escape in the evolution of the particle distribution and the role of delays in the received photons due to the light crossing time of different parts of the emitting region.
The model includes the presence of a quiescent spectrum, which is assumed to be represented by and thus fitted to the broad band spectral distribution observed in 1994 (Macomb et al. macomb95, 1995). To this a flaring component is added and this is constrained by the observed spectral and temporal evolution (the parameters for both the stationary and variable spectra are reported in Table 5). The values of the model parameters are very close to those derived in Maraschi et al. (maraschi\_letter, 1999) in a somewhat independent way, i.e. trying only to fulfill the constraints provided by the X–ray and TeV spectra, simultaneous but averaged over the whole flare. On this basis it is very likely that the model discussed here satisfies the TeV constraints, which are however not directly taken into account (the computation of the TeV light curve in this complex model would require to treat the non–locality of the IC process, which is beyond the scope of this paper).
Clearly a parametric prescription does not reproduce a priori a specific acceleration process, but we rather tried to constrain its form from the observed evolution. The main constraints on the form of the acceleration term are the following:
* Particles have to be injected at progressively higher energies on the flare rise timescale to produce the hard lag.
* Globally the range of energies over which the injection occurs has to be narrow, to give rise to a peaked spectral component.
* A quasi–monochromatic injection function for the flux at the highest energies to reach its maximum after that at lower ones (as opposed to e.g. a power law with increasing maximum electron Lorentz factor).
* The emission in the LECS band from the particles which have been accelerated to the highest energies (i.e. those radiating initially in the MECS band) should not exceed that from the lower energy ones, as after the peak no further increase of the (LECS) flux is observed; this in particular requires for the injection to stop after reaching the highest energies.
* The total decay timescale might be dominated by the achromatic crossing time effects, although the initial fading might be determined by different cooling timescales.
The acceleration term has then been described as a Gaussian distribution in energy, centered at a typical particle Lorentz factor $`\gamma _\mathrm{c}(\mathrm{t})`$ and with width $`\sigma =0.01\gamma _\mathrm{c}(\mathrm{t})`$, exponentially evolving with time: $`\gamma _\mathrm{c}(\mathrm{t})e^{(\mathrm{t}_{\mathrm{max}}\mathrm{t})}`$, where t<sub>max</sub> corresponds to the end of the acceleration phase, when the maximum $`\gamma _{\mathrm{c},\mathrm{max}}`$ is reached. The injected luminosity is assumed to be constant in time.
The duration of the injection, assumed to correspond to the light crossing time of the emitting region, is such that it mimics the passage of a shock front (i.e. a moving surface passing through the region in the same time interval, see Chiaberge & Ghisellini chiab\_gg, 1999).
It should be also noted that, within this scenario, the symmetry between the rise and decay in the soft energy light curve seems to suggest that if the energy where most of the power is released is determined by the balance between the acceleration and cooling rates, at this very same energy the latter timescales are comparable to the light crossing time of the region.
The predictions of the model are shown in Figures 7 and 7, in the form of light curves and spectra at different times, respectively. In particular, in Figure 7 the light curves (normalized over the stationary component) at the two centroid energies of the soft and hard bands considered in the data analysis are reported. The presence of a hard lag is clearly visible as well as the larger variability amplitude at higher frequencies. Also the spectral evolution, reported in Figure 7 in $`\nu F_\nu `$ in the observed energy range, seems to be at least in qualitative agreement with what observed (compare with Fig. 3).
A posteriori it is important to notice that the inclusion of a quasi–stationary component and of high energy emission (from synchrotron self–Compton) are indeed crucial in order to reproduce the retarded spectral variability at energies above a few keV. This constitute a further difference with respect to the discussion of Maraschi et al. (maraschi\_letter, 1999), where a simplified description of the time decay was used, based on evidence for an energy dependence of the flare decay timescale (obtained modeling the decay with no baseline).
## 6. Conclusions
BeppoSAX has observed Mkn 421 in 1997 and 1998. We analyzed and interpreted the combined spectral and temporal evolution in the X–ray range. During these observations the source has shown a large variety of behaviors, both concerning the X–ray band itself and its variability properties with respect to the $`\gamma `$–ray one, providing us with a great wealth of information, but at the same time revealing a richer than expected phenomenology.
Several important results follow from this work:
* The X–ray and 2 TeV light curves peak simultaneously within one hour, although the halving time in the TeV band seems shorter than those at LECS and MECS energies (see Maraschi et al. maraschi\_letter, 1999).
* The detailed comparison of 0.1–1.5 keV and 3.5–10 keV band light curves shows that the higher energy band lags the softer one, with a delay of the order of 2–3 ks. This finding is opposite to what has been commonly observed in HBL X–ray spectra (see Paper I).
* Moreover, extracting LECS$`+`$MECS spectra for $``$ 5 ks intervals, we were able to follow in detail the spectral evolution during the flare. For the first time we could quantitatively track the shift of the peak of the synchrotron component moving to higher energy during the rising phase of the flare, and then receding.
* An energy dependence of the shape of the light curve during the flare has been revealed: at low energies the shape is consistent with being symmetric, while at higher energies is clearly asymmetric (faster rise) (see Paper I).
* Evidence has been found for the presence of the IC component, and more importantly for its substantial variability, which is possibly delayed with respect to the synchrotron one.
These findings provide several temporal and spectral constraints on any model. In particular, they seem to reveal the first direct signature of the ongoing particle acceleration mechanism, progressively “pumping” electrons from lower to higher energies. The measure of the delay between the peaks of the light curves at the corresponding emitted frequencies thus provides a tight constraint on the timescale of the acceleration process.
Indeed, within a single emission region scenario, we have been able to reproduce the sign and amount of lag by postulating that particle acceleration follows a simple exponential law in time, stops at the highest particle energies, and lasts for an interval comparable to the light crossing time of the emitting region. If this timescale is intrinsically linked to the typical source size, we indeed expect the observed light curve to be symmetric at the energies where the bulk of power is concentrated and an almost achromatic decay. The same model can account for the spectral evolution (shift of the synchrotron peak) during the flare.
The other very important clue derived from the analysis is the presence of a quasi–stationary contribution to the emission, which seems to be dominated by a highly variable peaked spectrum, possibly maintaining a quasi–rigid shape during flares. The decomposition of the observed spectrum into these two components might allow us to determine the nature and modality of the energy dissipation in relativistic jets.
We are grateful to the BeppoSAX Science Data Center (SDC) for their invaluable work and for providing standardized product data archive, and to the RossiXTE ASM Team. We thank Gianpiero Tagliaferri and Paola Grandi for their contribution to our successful BeppoSAX program, and for useful comments and the anonymous referee for suggestions which have improved the clarity of the paper. AC, MC and YHZ acknowledge the Italian MURST for financial support. This research was supported in part by the National Science Foundation under Grant No. PHY94–07194 (AC). Finally, GF thanks Cecilia Clementi for providing tireless stimulus.
## Appendix A A. A caveat on continuously curved spectral models
The increasing quality of the available X–ray spectral data, both in terms of signal–to–noise ratio and energy resolution, has made necessary to consider more complex fitting models. One of the most interesting continuum feature in the X–ray spectra of blazars is the curvature of the synchrotron component, as good quality data could enable us to estimate the energy of the emission peak. As clearly the traditional single or broken power law models are unsatisfactory, we developed the curved spectral model presented in this paper: a simple analytic expression representing a continuous curvature, two “pivoting” points –useful for analysis purposes– asymptotically joining two power law branches, whose slope is completely determined by the behavior of the function at the pivoting energies.
Another increasingly popular model has been introduced by Giommi et al. (giommi\_2155, 1998) (to reproduce BeppoSAX data of PKS 2155$``$304, and then used also for Mkn 421 by Guainazzi et al. guainazzi\_mkn421, 1999 and Malizia et al. malizia\_mkn421, 1999), but unfortunately this is affected by a problem that can lead to misleading results.
This more generally occurs for all curved continuum models defined as:
$$F(E)=E^{\stackrel{~}{\alpha }(E)}$$
(A1)
where the function $`\alpha (E)`$ is suitably defined to smoothly change between two asymptotic values $`\alpha _1`$ and $`\alpha _2`$. The specific choice by Giommi et al. (giommi\_2155, 1998) is:
$$\stackrel{~}{\alpha }(E)=f(E)\alpha _1+(1f(E))\alpha _2$$
(A2)
with
$$f(E)=\left(1e^{E/E_0}\right)^\beta .$$
(A3)
Thus $`\stackrel{~}{\alpha }\alpha _1`$ for $`EE_0`$, and $`\stackrel{~}{\alpha }\alpha _2`$ for $`EE_0`$. While this exactly the behavior of the spectral index function, this does not correspond to the real spectral index of the function $`F(E)`$, conventionally defined as:
$$\alpha (E)\frac{d\mathrm{ln}F(E)}{d\mathrm{ln}E}.$$
(A4)
In fact for $`F(E)E^{\stackrel{~}{\alpha }(E)}`$ this yields
$$\alpha (E)=\stackrel{~}{\alpha }(E)+\mathrm{ln}E\frac{d\stackrel{~}{\alpha }(E)}{d\mathrm{ln}E}=\stackrel{~}{\alpha }(E)+\delta _{\stackrel{~}{\alpha }}(E)$$
(A5)
thus including an unavoidable term $`\delta _{\stackrel{~}{\alpha }}(E)`$, present because of the not vanishing (by definition) derivative of $`\stackrel{~}{\alpha }(E)`$.
For the above specific choice of $`\stackrel{~}{\alpha }(E)`$:
$$\delta _{\stackrel{~}{\alpha }}(E)=\beta \frac{E}{E_0}\mathrm{ln}E(\alpha _1\alpha _2)e^{E/E_0}\left(1e^{E/E_0}\right)^{\beta 1}.$$
(A6)
The real spectral index then takes a value (sensibly) different from the one expected over a large range of energies.
The effect can be better illustrated with an example. Let us consider a set of spectral parameters that would well reproduce the spectrum of the peak of the 1998 flare of Mkn 421: $`\alpha _1=0.75,\alpha _2=1.5,E_0=3.5`$ keV, $`\beta =1.0`$. In Figure 8 we show the spectrum together with $`\stackrel{~}{\alpha }(E)`$ and $`\alpha (E)`$ as derived by the model over an energy interval broader than that effectively covered by data. For reference we also plotted the values of $`\alpha `$ derived from our curved model fits, to show that the real spectral index $`\alpha (E)`$ matches them. A few features emerge:
1. according to $`\stackrel{~}{\alpha }(E)`$ between $`E_0`$ and 10 keV there is an apparent wiggle in the spectrum, which steepens beyond the supposedly asymptotic value $`\alpha _2`$ before reaching it at about 30 keV.
2. within the 0.1–10 keV range $`\alpha (E)`$ and $`\stackrel{~}{\alpha }(E)`$ never match each other except for $`E=1`$ keV, when $`\delta _{\stackrel{~}{\alpha }}(E)`$ vanishes.
3. the range spanned by the true spectral index is broader, and in particular it is always steeper than the value expected on the basis of $`\stackrel{~}{\alpha }(E)`$ in the whole band 1–20 keV. $`\mathrm{\Delta }\alpha `$ can be as much as $`+`$0.42 (at $`E5.5`$ keV).
4. at no energy in the 0.1–10 keV band the actual spectrum has a slope equal to either $`\alpha _1`$ or $`\alpha _2`$, and thus these two parameters do not have any real descriptive meaning (this is strikingly illustrated by the displacement of the real data points from the expected curve).
5. the total change of slope in the 0.1–10 keV interval would be expected $`\mathrm{\Delta }\stackrel{~}{\alpha }0.68`$ while it actually turns out to be 1.04 (and even more dramatically between 0.5 and 5.0 keV, $`\mathrm{\Delta }\stackrel{~}{\alpha }0.46`$ while for the data it is 0.86; see Table 1).
Summarizing, there is no way to define a curved continuum model by means of a suitable function describing the power law exponent.
## Appendix B B. More details on the adopted curved model
The relationship between $`\alpha (E)`$ and $`\alpha _{\mathrm{}}`$ and $`\alpha _+\mathrm{}`$ is:
$$\alpha (E)\frac{d\mathrm{ln}F(E)}{d\mathrm{ln}(E)}=\alpha _{\mathrm{}}(\alpha _{\mathrm{}}\alpha _+\mathrm{})\frac{\left(E/E_\mathrm{B}\right)^f}{1+\left(E/E_\mathrm{B}\right)^f}$$
(B1)
yielding:
$`\alpha _{\mathrm{}}`$ $`=`$ $`{\displaystyle \frac{B(1+A)\alpha _1A(1+B)\alpha _2}{BA}}`$ (B2)
$`\alpha _+\mathrm{}`$ $`=`$ $`{\displaystyle \frac{(1+B)\alpha _2(1+A)\alpha _1}{BA}}`$ (B3)
where $`A=\left(E_1/E_\mathrm{B}\right)^f`$, and $`B=\left(E_2/E_\mathrm{B}\right)^f`$.
## Appendix C C. More details on the broken power law fits.
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# Rotating Band Pion Production Targets for Muon Colliders and Neutrino Factories aafootnote aPresented at the ICFA/ECFA Workshop ”Neutrino Factories based on Muon Storage Rings” (𝜈FACT’99), Lyon, France, 5–9 July, 1999.
## 1 Introduction and Motivation
Conceptual design studies for muon colliders that have taken place since the mid-1990’s have motivated the design of pion production targets that can operate and survive with megawatt-scale pulsed proton beams. Over the past few months, the design timescale and potential learning curve for such targets has effectively been abbreviated by the expanded interest in neutrino factories – muon storage rings dedicated to producing neutrino beams that will require similarly large muon currents to muon colliders and that have prospects for being built on a shorter timescale – perhaps to be ready within the next decade.
We present an update on a previous report describing a conceptual design for a cupronickel rotating band pion production target for muon colliders that was proposed as a relatively conservative extrapolation from existing targets. A more detailed write-up on this target is in progress . (Note that another rotating band design based on reference is presented elsewhere in these proceedings .)
The rotating band design is readily adaptable to neutrino factories and, indeed, its conservative nature makes its development particularly well matched to the shorter timescales for this application. Alternative, more exotic, targetry options that were originally proposed for muon colliders – such as pulsed mercury jets – may involve extensive multi-year exploratory experimental R&D programs that do not appear to be particularly compatible with the shorter timescales envisaged for neutrino factories.
This paper is laid out as follows. General design considerations and strategies for high power solid targets form the topic of the following section. Sections 3 through section 5 then focus in on a review of the specific cupronickel band design that was proposed in reference . Specifically, section 3 gives an overview of the conceptual design and strawman specifications, section 4 summarizes the results of computer simulations predicting its pion yield performance and characterizing the heating effects from the beam, and section 5 discusses stress and durability issues. Section 6 steps back from the parameters in to examine the technology options and parameter optimizations available to the rotating band target concept for the spectrum of possible proton driver scenarios at both muon colliders and neutrino factories. The paper concludes with comments on the potential for the rotating band target concept and with an outlook on the R&D program required to bring this conceptual design to practical fruition.
## 2 General Design Strategies for High Power Production Targets
This section gives an overview of the general design goals that were considered important when proposing the target design of reference , and on the strategies employed to achieve these design goals. The intention was to design a target station that:
1. can be designed quickly and with relatively modest resources, so it will assuredly not hold up the overall development of neutrino factories or muon colliders
2. will clearly survive any beam-induced stresses it might be subjected to and have an acceptable lifetime
3. is relatively straightforward and affordable to build and maintain (including target replacement and disposal)
4. has pion yields per proton and phase space densities that are as good as, or not much inferior to, what could be achieved with the more idealized targets that could be designed for operating with low beam powers.
The first item implies that the R&D program for such a target can be conducted mainly through paper studies, computer simulations and engineering computer-assisted design studies, perhaps augmented by a modest amount of mechanical prototyping and/or beam tests if convenient. Such an R&D program appears plausible for the strawman design scenario of reference . Items 3 and 4 also appear to be met by this scenario, as can be judged from the information presented in section 3 of this paper and elsewhere .
The remainder of this section addresses item 2 in the list, since the single most difficult design constraint for rotating bands and other solid targets for muon colliders is the requirement of survivability in the face of instantaneous beam energies per proton pulse of order 100 kJ and megawatt-scale beam powers, i.e., comparable to or larger than the largest existing proton facilities. A sound design strategy to satisfy item 2 is to choose beam and target parameters such that the maximum material stresses and radiation exposures do not go beyond what has been explicitly achieved in existing targets or targetry studies. It is clear that this can always be achieved in principle, even for the highest beam powers under consideration, by:
1. choosing a target material with appropriate mechanical properties and moderate thermal stresses, consistent with optimizing the pion yield. In practice this might mean choosing from materials with medium atomic numbers in the range from titanium ($`\mathrm{Z}=22`$) through nickel ($`\mathrm{Z}=28`$), as discussed further in sections 4 and 6, and
2. sufficiently spreading out the spot size to cope with higher energy beam pulses, and
3. rotating or otherwise moving the target to continually expose new areas of the target to each beam pulse. This limits both the instantaneous local thermal stresses and the lifetime radiation exposure of the target material.
Concerning the final item of strategy, the considerable potential for designs that continually expose the beam to new target material is amply illustrated by the tungsten target design for the proposed Accelerator Production of Tritium (APT) project at Los Alamos National Laboratory. A 1.7 GeV, 100 mA beam is continually rastered across a 19 x 190 cm area of the APT tungsten target to limit the local heating and thermal shock stresses. Impressively, the projected 170 MW beam power onto the APT target is approximately two orders of magnitude larger than the applications considered in this paper.
As a detailed difference from the approach for the APT, the band target uses the more traditional solution for high power targets of moving the target material rather than the beam spot. The idea of rotating or scanning high-power targets has already found successful application in several existing facilities including, for example, the trolled tungsten-rhenium target for the SLC positron source and the rotating nickel target at the BNL “g-minus-2” experiment . Target rotation can limit temperature rises and the consequent thermal stresses on timescales of a second or less and, in the longer term, it spreads out the radiation load over much more material. The Fermilab antiproton source target is rotated more slowly than the preceding examples in order to spread out the radiation dose over the circumference of the target disk, as is discussed further in section 5.
## 3 A Conceptual Design and Straw-man Specifications for a Cupronickel Rotating Band Target
The discussion presented in this section and the two that follow are specific to the cupronickel rotating band target design proposed in reference , beginning with the overview of that conceptual design given in this section.
Figure 1 gives a schematic overview of the target concept presented in reference and figure 2 zooms in on the production region. It should be emphasized that such details as the rollers and cooling setup are shown only schematically and that no concerted effort has been put into their design or layout. The cupronickel target band is enclosed in a 20 Tesla solenoidal magnetic pion capture magnet whose general design has previously been studied by the Muon Collider Collaboration. The pion secondaries spiral along the solenoidal magnetic channel before decaying into the muon bunches needed for cooling, acceleration and injection into the collider ring. The radius of the solenoidal channel, 7.5 cm, and the magnetic field strength, 20 Tesla, are those commonly assumed for recent muon collider studies . The design modification specific to this particular geometry concerns the provision of entry and exit ports for the target band. Example coil geometries show that these ports can be provided with little modification to the design of the solenoidal channel.
Table 1 gives some relevant parameters for the cupronickel band and figure 3 illustrates the trajectory of the proton beam into the target band. The band parameters correspond to the proton beam parameters for a muon collider that are given in table 2. Different parameter values might be appropriate for alternative beam scenarios at muon colliders or for neutrino factories. The geometry of the band is chosen to approximately maximize the pion yield. The general requirements are that the proton pathlength through the target material should be approximately 2 nuclear interaction lengths and that the band should be thin enough to allow most of the pions to escape the target. To optimize the pion yield , the trajectory of the beam through a chord of the target band is at a tilt angle of 150 milliradians to the axis of the solenoid.
The cupronickel band is guided and powered by several sets of rollers that can be connected by driveshafts to remotely housed motors outside any radiation shielding. This scenario has been taken from the Zenzimmer mills that are used for pressing metal sheets and has the attraction of being mechanically very simple in the high radiation area surrounding the production region. Bennett even suggests the total elimination of moving parts other than the target band by using electromagnetic guidance and rotation of the target band by linear motors. Procedures for installation and extraction of the target band are proposed in reference .
The rotation rate of 3 m/s given in table 1 corresponds, for the 15 Hz beam frequency, to a target advance per pulse by 1/3 of the chord spanned by the proton beam. The three overlapping proton pulses in any part of the band imply a maximum total temperature rise approximately double the instantaneous rise from each individual pulse. Temperature rises and stresses will be further discussed in section 5.
A competing concern to the target heating stress that limits the acceptable rotation rate of the target is the eddy currents induced by the rapidly rotating band in a strong magnetic field. The eddy current power is proportional to the conductivity and to the square of the band velocity and a very approximate analytic calculation predicts that the power dissipated will be of order several kW. This is more than an order of magnitude below the beam heating and so is clearly a manageable heat source. The power will have to be supplied by the electric motor driving the rotation of the target and the requirements on the drive mechanism are within typical operating parameters for Zenzimmer steel mills.
Cupronickel alloys are preferred over both copper and nickel for the particular targetry application considered in this paper because their lower electrical conductivity will reduce the eddy currents from rotation through the magnetic field of the solenoidal capture magnet. Cupronickel alloys such as, for example, alloy 715 produced by Olin Brass have essentially the same density and interaction lengths as copper and nickel and have very similar mechanical properties. However, Olin alloy 715 has an electrical conductivity at 20 degrees centigrade of only 2.6 MS/m, compared to 58 MS/m for copper and 14 MS/m for nickel.
The hot portion of the band is rapidly carried away from the production region into a water cooling channel, as is the case for the nickel production target at the BNL g-2 experiment . Although both the peak temperature and power are much larger than in the g-2 experiment, the surface area of the band has been chosen large enough to give heat transfer rates of approximately 30 $`\mathrm{W}/\mathrm{cm}^2`$ that are comfortable even in the presence of a steam boundary layer. It has been suggested that the target and capture channel should be in a helium atmosphere to allow the easy distillation of water vapor and any other impurities. This will have negligible effect on pion production due to helium’s low density and low atomic number.
## 4 Pion Yield Predictions and Thermal Parameters of the Cupronickel Band Target
Detailed MARS tracking and showering Monte Carlo simulations were performed to obtain the pion yields per proton for the beam and target parameters of that are reproduced in tables 1 and 2. The graphical results in figure 4 correspond to yields of $`Y_+`$ = 0.622 and $`Y_{}`$ = 0.612 positive and negative pions per proton for the momentum range 0.05$`<`$p$`<`$0.80 GeV/c. The peak energy deposition density was found to be 68.6 J/g per pulse, corresponding to a temperature rise of $`\mathrm{\Delta }T`$=151C and a total power dissipation in the target of 0.324 MW. These predictions are summarized in table 3.
These pion yields and the predicted phase space densities are almost identical to the best predicted yields for the exotic liquid mercury jet targets that are also under consideration for muon colliders and neutrino factories. Optimization studies for both mercury jet and band targets further suggest that such yields are rather close to the optimum that could be obtained even for a low power proton beam, thus satisfying the fourth of the design goals stated in section 2.
## 5 Stress and Durability Issues for the Cupronickel Band
As an application of the design strategy presented in section 2, this section examines the possibilities for evaluating the survivability of the cupronickel band target through benchmarking to existing targets. The Fermilab antiproton source target appears to be one of the most suitable targets for benchmarking the band, so we begin by summarizing its design and operating parameters.
The Fermilab antiproton target consists of a vertical stack of 3 nickel target disks plus one copper target disk, each approximately 1 cm thick and 4.7 cm in radius, and interspersed with copper cooling disks. The target is cooled by forcing air up the vertical axis and through channels in the copper cooling disks. The stack of disks is enclosed in a titanium can that would contain target the material in case of failure, but the can is not in contact with the disks and so is irrelevant for considerations of mechanical survivability. A 120 GeV proton beam passes through a chord of the selected target disk, with an intensity of 1.6 to 2.1$`\times `$10<sup>12</sup> protons per 1.6 $`\mu `$sec pulse, which is incident every 2.4 seconds.
The energy per pulse of the antiproton source, up to 40 kJ, is almost an order of magnitude below the 256 kJ muon collider specification of table 2. Despite this, the round beam spot has a gaussian sigma at entry of only 140 microns , and this exposes the antiproton target to local energy deposits and temperature rises much larger than MARS predictions for the cupronickel band, which assume a much larger elliptical spot of dimensions $`\sigma _x=1.5`$ mm, $`\sigma _y=10`$ mm. As shown in table 1, the cupronickel band is predicted to sustain a maximum instantaneous energy deposition of approximately 70 J/g per proton pulse, compared to the 500-600 J/g maximum depositions at the the Fermilab antiproton source. The impressive peak temperature rise in the antiproton target is 1100$`{}_{}{}^{\mathrm{o}}\mathrm{C}`$ over 1.6 $`\mu `$sec, to be compared to 150$`{}_{}{}^{\mathrm{o}}\mathrm{C}`$ instantaneously and approximately 300$`{}_{}{}^{\mathrm{o}}\mathrm{C}`$ over a fraction of a second for the band.
While the comparison of the preceding paragraph is suggestive that the instantaneous heat stresses on the cupronickel band might be acceptable, there are several issues to be resolved before one has confidence in the benchmarking comparison with the Fermilab antiproton target. Issues include the effects of the different target geometries and the question of how closely the 1.6 $`\mu `$sec timescale for energy deposition in the antiproton target approximates the instantaneous energy deposition in the band.
To further understand the shock heating stresses on the cupronickel band, finite element computer simulations have been performed using ANSYS, a commercial package that is very widely used for stress and thermal calculations. Energy density distributions for the simulations were generated using the MARS particle production Monte Carlo package. The very preliminary simulations show periodic returns to maximum stress (i.e. “ringing”) but with little or no amplification beyond the initial stress. This is encouraging, and seems plausible given that the oscillations occur in the band dimension that is much shorter than the other two so the geometry is quasi-one dimensional and with little potential for focusing. Further ANSYS simulations of band target geometries are commencing , and these are intended to include explicit benchmarking simulations on the Fermilab antiproton target.
The other threat to the survival of the target band comes from exposure to radiation. This can change the material properties of target materials by transmuting some of the target atoms to new isotopes or elements and by causing dislocations in the atomic lattice. Face-centered cubic metal lattices such as copper and nickel are known from experience to survive radiation damage better than body-centered cubic metals like iron or tungsten. Target damage studies at Los Alamos National Laboratory and elsewhere predict a risk of failure for copper and nickel targets after integrated doses somewhere in the range $`10^{21}`$ to $`10^{22}`$ minimum ionizing particles per square centimeter, which corresponds to approximately 0.3 – 3 GJ/g of deposited energy.
It is straightforward to obtain a rough estimate of the integrated doses on the cupronickel by noting that the parameters of tables 1 and 2 correspond to doses, along the centerline of the band and for an accelerator year of $`10^7`$ seconds, that accumulate to:
$$\mathrm{summed}\mathrm{energy}\mathrm{deposition}140\mathrm{J}/\mathrm{g}\times \frac{3\mathrm{m}/\mathrm{s}}{2\pi \times 2.5\mathrm{m}}\times 10^7\mathrm{s}0.3\mathrm{J}/\mathrm{g},$$
(1)
where 140 J/g is the energy deposited per rotation and the second term is the target rotation frequency. Since 0.3 GJ/g is the lower limit for predictions of target failure, the initial conclusion to be drawn is that it may well be acceptable to replace the target band after each year’s running. More detailed studies are obviously needed to check and refine this first simple estimate.
To recap, this section has provided a first look at benchmarking the cupronickel band target to the Fermilab antiproton target and other existing data from operating targets. The two indications from these initial comparisons are that:
1. the much greater temperature rises in the operating Fermilab target give some initial confidence in the short-term survivability of the target but the comparison has not yet been made rigorous
2. the simple calculation of equation 1 suggests that potentially damaging radiation doses would accumulate over an acceptably long timescale for the cupronickel band target, despite the 3.8 MW beam power, because the radiation dose is spread over the entire circumference of the band rather than being concentrated in one region.
More detailed studies are beginning to check and clarify these very preliminary findings.
## 6 Technology Options and Parameter Optimizations
The parameters of the cupronickel band target presented in reference and section 3 represent no more than an educated guess at a relatively optimal configuration for the 4 MW proton driver parameters of table 2. The R&D required to verify and refine this design is just beginning and it is expected that other parameter values and design refinements will likely turn out to be better suited for this and other proton driver scenarios.
Entry-level neutrino factories have been discussed with proton driver powers of 1 MW or less, and this will clearly allow some relaxation of target design parameters. At the other end of the scale, proton drivers for neutrino factories of up to 20 MW have also been discussed at CERN. In some neutrino factory scenarios, the proton beam is partitioned into smaller bunches than is feasible for muon colliders and this will generally also allow for a relaxation of target parameters.
Some examples of parameters that need to be optimized depending on the specific proton driver scenario are:
1. the beam spot size and the cross-sectional area of the band. These will tend to become larger for increasing proton pulse energy
2. the target circumference and rotation rate: these will both tend to increase with increasing average beam power; the first to increase the surface area available for cooling the target and the second to moderate the localized target heating.
These parameter values and the those of the proton driver will also determine the optimal technology decisions for several design options, including:
1. the cooling technology. Helium gas cooling might perhaps be technically simpler than water cooling but the latter can provide larger heat transfer rates per unit area. Radiative cooling is a third possible option for refractory target materials such as graphite , tantalum or tungsten
2. the band drive and guidance. The optimal choice might depend on the level of frictional drag from eddy currents in the magnetic field of the capture solenoid and this depends in turn on the target band cross section, the target material and the target rotation velocity. The Zenzimmer-type rollers presented in reference can comfortably handle the several kilowatt frictional drag for the cupronickel band and default parameter set. Smaller frictional loads would instead allow for more modest electromagnetic guidance and a linear electric drive, as suggested in reference
3. the choice of target band material, as will now be discussed further.
Evaluations of several targets for pion yield and for heating and shock stresses were performed using MARS simulations. It should be noted that these comparisons between elements depended on target densities and on the specific targetry scenario used, so they provide only approximate guidance. However, the trend in heating stress was clear, with the stresses best for low atomic number (Z) elements and becoming rapidly worse with increasing Z. To balance this, pion yields were predicted to be lower for low-Z materials than for those with medium or high atomic number. Instead of a steady rise in yield, the yield was found to plateau somewhere between Al (Z=13) and Cu (Z=29) (elements in between were not investigated) and then to remain constant to within the accuracy of the study all the way out to the high-Z elements tungsten (Z=74) and mercury (Z=80).
This comparative study of target elements, along with the outstanding track record of both copper and nickel as target materials, was part of the basis for the choice of cupronickel as the target band material proposed in reference , with nickel (Z=27) and copper (Z=28) both towards the low-Z end of the plateau in pion yield. It would clearly be helpful to repeat the yield study for elements between Al and Cu to pin down the exact fall-off position of the yield plateau, especially since some of the elements in between and their alloys are known to have excellent mechanical and thermal properties for targetry applications, particularly Ti (Z=22), V (Z=23), Cr (Z=24) and Mn (Z=25) .
While initial studies suggest the suitability of these medium-Z elements for a 4 MW proton driver for muon colliders, lower-Z elements such as graphite can always be considered as options to extend the rotating band design to larger safety margins or to even more demanding beam specifications. The penalty to be paid is that graphite, for example, appears to have only about 2/3 the pion yield of the elements on the yield plateau. Despite this, such bands might anyway be used in the same target station as uses medium-Z bands – either as insurance against problems with the medium-Z band or as an entry-level band to be eventually replaced by one with a higher yield.
## 7 Conclusions and Outlook
This paper has reviewed and enlarged on a previously proposed design for a high power pion production target station based around a rotating cupronickel band target. The design scenario is mechanically rather straightforward and can be readily extrapolated from, and benchmarked to, existing targets. Initial computer simulations have predicted relatively optimal pion yields and the initial comparisons, made in section 5, with the operating Fermilab antiproton target suggest that the cupronickel band target parameters of will survive the proton beams for at least some of the neutrino factory and muon collider scenarios that are under discussion. More generally, it was argued in sections 2 and 6 that viable rotating band target designs almost certainly exist for any of the proton driver scenarios that have been seriously discussed for either muon colliders or neutrino factories. The question is one of design optimization rather than feasibility, and whether the demands of target survivability force any significant compromises on the pion yield or phase space density that can be supplied by the target.
The evolutionary nature and relative simplicity of the target design and concept appears to make it compatible with a few-year R&D program to first explore the options and parameter space and then, if all goes well, to refine and develop the design towards construction at a neutrino factory or muon collider. Such an R&D program might develop something like this:
* in 2000: A) Paper studies to further clarify the design options and issues, including developing databases of target material properties, existing targets, past experimental targetry studies and contact information on targetry experts and contacts, B) stress simulations and optimization studies using ANSYS or a similar finite element analysis package. C) beginning engineering studies on the target layout, mechanical issues and cooling options, D) develop a conceptual design for the beam dump, E) further yield optimization studies on target geometry and the band material, using MARS or similar particle production codes, F) further particle tracking studies to explore the integration of the target design with the beam dump and capture and phase rotation channels.
* in 2001: More detailed design studies for specific scenarios. Detailed assessments of, and comparisons with, other target options.
* 2002 and beyond: continuing design studies might lead to mechanical prototyping if this is found to be necessary and to beam tests if these are convenient and would add significantly to the existing pool of experimental knowledge.
* by approximately 2004: ready to begin constructing a rotating band target for a neutrino factory or muon collider.
The R&D issues for successive years have been spelled out in progressively less detail but will necessarily involve a ramp-up in manpower, beginning from perhaps 2 full-time equivalent people in 2000. Such an R&D program looks to be of modest extent compared to the efforts required for the more exotic targets needed for, e.g., neutron spallation sources. Indeed, the overall program looks compatible with installing a rotating band target station at a neutrino factory or muon collider facility for potential operation as early as 2006 or 2007.
## 8 Acknowledgments
Studies on the cupronickel band target design have been conducted in collaboration with R.J. Weggel, N.V. Mokhov and S.S. Moser. This work has also benefitted from discussions and techical advice from G. Bunce and C. Pai. The organizers and secretariat of NUFACT99 are to be commended for a well-organized and stimulating workshop.
This work was performed under the auspices of the U.S. Department of Energy under contract no. DE-AC02-98CH10886.
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# 1 Introduction
## 1 Introduction
The 2D sine-Gordon model, defined by the action:
$$S=\frac{\pi }{\gamma }d^2x,=(_\nu \varphi )^2m^2(cos(2\varphi )1),$$
(1.1)
where $`\gamma `$ is the coupling constant and $`m`$ is related to the mass scale, is one of the simplest Integrable Quantum Field Theories. It possesses an infinite number of conserved charges $`I_{2n+1}`$, $`nZ`$, in involution. At present, much is known about the corresponding scattering theory. It contains solitons, antisolitons and a number of bound states called “breathers”. The mass spectrum and the S-matrix have been known for about 20 years . Despite this on-shell information, the off-shell Quantum Field Theory is much less developed. In particular, the computation of the corresponding correlation functions is still an important open problem. Actually, some progress towards this direction has been made recently. For instance, the exact Form-Factors ( FF’s) of the exponential fields $`<0|\mathrm{exp}\varphi (0)|\beta _1,\mathrm{},\beta _n>`$ were computed . This allows one to make predictions about the long-distance behaviour of the corresponding correlation functions. On the other hand, some efforts have been made to estimate the short distance behaviour of the theory in the context of the so-called Conformal Perturbation Theory (CPT) . By combining the previuos techniques (FF’s and CPT), it has been possible to extimate several interesting physical quantities ( and references therein). In addition, the exact expression for the Vacuum Expectation Values (VEV’s) of the exponential fields ( and some descendents ) were proposed in . The VEV’s provide a highly non-trivial non-perturbative information about the short distance expansion of the two-point correlation functions. The FF’s and CPT approaches permit to make predictions about the approximate behaviour of the correlation functions of the sine-Gordon theory in the infrared and ultraviolet regions correspondingly. What remains still unclear however is the explicit form of the correlation functions, in particular their intermediate behaviour and analytic properties, a question of primary importance from the field-theoretical point of view. Some exact results exist only at the so-called free-fermion point $`\gamma =\frac{\pi }{2}`$ .
There exists another approach to the sine-Gordon theory. It consists in searchig for additional infinite-dimentional symmetries and is inspired by the success the latter had in the 2D CFT. In fact, it has been shown in that the sine-Gordon theory possesses an infinite dimensional symmetry provided by the $`\widehat{s}l(2)_q`$ algebra. However, this symmetry connects the correlation functions of the fields in the same multiplet without giving a sufficient information about the functions themselves.
It is known to some extent that there should be another kind of symmetry present in the sine-Gordon theory. Actually, it is known that it can be obtained as a particular scaling limit of the so-called XYZ – spin chain . The latter is known to possess an infinite symmetry obeying the so-called Deformed Virasoro Algebra (DVA). It is natural to suppose that in the scaling limit, represented by SG, there should be present some infinite dimensional symmtery, a particular limit of DVA. At present, a lot is known about the mathematical structure of DVA, in particular the highest weight representations and the screening charges have been constructed . What remains unclear is how the corresponding symmetry is realized in the sine-Gordon field theory, for example what is the action of the corresponding generators on the exponential fields, what kind of restrictions it imposes on their correlation functions etc.
In this paper we present a construction of a Virasoro symmetry directly in the sine-Gordon theory. Although we are of course interested in the quantum theory, we restrict ourself to the classical picture in this paper. Also, though it will be clear how to implement it in the general field theory, we are mainly concerned here with the construction of this symmetry in the case of the N-soliton solutions. One of the reasons for this is that the symmetry in this case is much simpler realized ( in particular it becomes local contrary to the field theory realization ). We were also inspired by the work of Babelon,Bernard and Smirnov . It was shown there that certain form-factors can be directly reconstructed by a suitable quantization of the N-soliton solutions. It was also shown in that certain null-vector constraints arise in sine-Gordon theory leading to integral equations for the corresponding form-factors. It is an intriguing question of what is the symmetry structure lying behind it and in particular its relation to the Virasoro symmetry we present here.
This paper is organised as follows. In the next Section we recall the construction of the so-called additional non-isospectral Virasoro symmetry of KdV theory. It is done in the context of the so-called algebraic approach and is a generalization of the well known dressing symmetries of integrable models . In Section 3 we explain how one can restrict this symmetry to the case of N-soliton solutions of KdV. It happens that it becomes local in this case, contrary to the field theory realization where it is quasi-local. In Section 4 we extend the Virasoro symmetry to the sine-Gordon soliton solutions. This is acheeved by introducing an additional time dynamics in the KdV theory. In such a way we obtain “negative” Virasoro flows which complete the “positive” ones of KdV to the whole Virasoro algebra in the sine-Gordon theory. Finally, in Section 5 we summarise the results obtained in the paper. We give some hints about the quantization of the classical picture presented here and discuss some important open problems.
## 2 Virasoro Symmetry of mKdV theory
### 2.1 Dressing Symmetries
Let us recall the construction of the Virasoro symmetry in the context of (m)KdV theory. It was shown in , following the so called algebraic approach, that it appears as a generalization of the ordinary dressing transformations of integrable models. Here we briefly recall the main results of this article. Being integrable, the mKdV system admits a zero-curvature representation:
$$[_tA_t,_xA_x]=0,$$
(2.1)
where the Lax connections $`A_x`$ , $`A_t`$ belong to $`A_1^{(1)}`$ loop algebra:
$`A_x`$ $`=`$ $`vh+(e_0+e_1),`$
$`A_t`$ $`=`$ $`\lambda ^2(e_0+e_1vh){\displaystyle \frac{1}{2}}[(v^2v^{})e_0+(v^2+v^{})e_1]{\displaystyle \frac{1}{2}}({\displaystyle \frac{v^{\prime \prime }}{2}}v^3)h`$ (2.2)
( $`v`$ is connected to the mKdV field $`\varphi `$: $`v=\varphi ^{}`$ ) and
$$e_0=\left(\begin{array}{cc}0& \lambda \\ 0& 0\end{array}\right)=\lambda E,e_1=\left(\begin{array}{cc}0& 0\\ \lambda & 0\end{array}\right)=\lambda F,h=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)=H$$
(2.3)
are the corresponding generators in the fundamental representation. The usual KdV variable $`u`$ is connected to the mKdV field $`\varphi `$ by the Miura transformation:
$$u=\frac{1}{2}(\varphi ^{})^2+\frac{1}{2}\varphi ^{\prime \prime }$$
(2.4)
(we denote by prime the derivative with respect to the space variable $`x`$ of KdV). Of great importance in our construction is the solution $`T(x,\lambda )`$ to the so called associated linear problem:
$$T(x,\lambda )(_xA_x(x,\lambda ))T(x,\lambda )=0$$
(2.5)
which is usually referred to (with suitable normalization) as a transfer matrix. A formal solution to (2.5) can be easily found:
$`T_{reg}(x,\lambda )=e^{H\varphi (x)}𝒫\mathrm{exp}\left(\lambda {\displaystyle _0^x}𝑑y(e^{2\varphi (y)}E+e^{2\varphi (y)}F)\right)=`$
$`=\left(\begin{array}{cc}A& B\\ C& D\end{array}\right).`$ (2.8)
It is obvious that this solution defines $`T(x,\lambda )`$ as an infinite series in positive powers of $`\lambda `$ with an infinite radius of convergence (we shall often refer to (2.8) as regular expansion). For further reference we present also the expansion of the corresponding matrix elements:
$`A=e^\varphi (1+{\displaystyle \underset{1}{\overset{\mathrm{}}{}}}\lambda ^{2n}A_{2n}),B=e^\varphi {\displaystyle \underset{0}{\overset{\mathrm{}}{}}}\lambda ^{2n+1}B_{2n+1},`$
$`C[\varphi ]=B[\varphi ],D[\varphi ]=A[\varphi ].`$ (2.9)
It is also clear from (2.8) that $`T(x,\lambda )`$ possesses an essential singularity at infinity where it is governed by the corresponding asymptotic expansion.
Obviously, the zero-curvature form (2.1) is invariant under the gauge transformation:
$$\delta _nA_x(x,\lambda )=[\theta _n(x,\lambda ),]$$
(2.10)
for $`A_x`$, and a similar one for $`A_t`$. A suitable choice for the gauge parameter $`\theta _n`$ goes through the construction of the following object:
$$Z^X(x,\lambda )=T(x,\lambda )XT(x,\lambda )^1,X=E,F,H$$
(2.11)
essentially the dressed generators of the underlying $`A_1^{(1)}`$ algebra. It is obvious by construction that it satisfies:
$$[,Z^X(x,\lambda )]=0,$$
(2.12)
i.e. it is a resolvent of the Lax operator L. As we shall see, this property is important for the construction of a consistent gauge parameter. In fact, let us insert $`T_{reg}`$ as defined in (2.8) in (2.11) and then construct
$$\theta _n^X(x,\lambda )(\lambda ^nZ^X(x,\lambda ))_{}=\underset{k=0}{\overset{n1}{}}\lambda ^{kn}Z_k^X(x),$$
(2.13)
where the subscript – (+) means that we restrict the series only to negative (non-negative) powers of $`\lambda `$. One can show that due to (2.12) the r.h.s. of (2.10) is of degree zero in $`\lambda `$ and therefore $`\theta _n^X`$ so constructed is a good candidate for a consistent gauge parameter. There is one more consistency condition we have to impose due to the explicit form of $`A_x`$ (2.2), namely $`\delta A_x`$ should be diagonal:
$$\delta _n^XA_x=H\delta _n^X\varphi ^{}.$$
(2.14)
This implies restrictions on the indices of the transformations: it happens that one must take even ones for $`X=H`$ ( $`\theta _{2n}^H`$ ) and odd ones for $`X=E`$ or $`F`$ ( $`\theta _{2n1}^{E,F}`$ ). The first transformations read explicitly:
$`\delta _1^E\varphi ^{}(x)`$ $`=`$ $`e^{2\varphi (x)}`$
$`\delta _1^F\varphi ^{}(x)`$ $`=`$ $`e^{2\varphi (x)}`$
$`\delta _2^H\varphi ^{}(x)`$ $`=`$ $`e^{2\varphi (x)}{\displaystyle _0^x}𝑑ye^{2\varphi (y)}+e^{2\varphi (x)}{\displaystyle _0^x}𝑑ye^{2\varphi (y)}.`$ (2.15)
Note that they are essentially non-local (this is true also for the higher ones). The algebra these transformations close is the (twisted) Borel subalgebra of $`A_1^{(1)}`$. Therefore the remaining ones can be found from(2.15) by commutation.
At this point we want to make an important observation. Consider the KdV variable $`x`$ as a space direction $`x_{}`$ of some more general system (and $`_{}_x`$ as a space derivative). Introduce the time variable $`x_+`$and the corresponding evolution defining:
$$_+(\delta _1^E+\delta _1^F).$$
(2.16)
It is then obvious from (2.15) that the equation of motion for $`\varphi `$ becomes:
$$_+_{}\varphi =2\mathrm{sinh}(2\varphi ),(or2sin(2\varphi )if\varphi i\varphi )$$
(2.17)
i.e. the sine-Gordon equation! We consider this observation very important since it provides a global introduction of sine-Gordon dynamics in the KdV system - a fact which was not known before.
The construction of the gauge parameter in the asymptotic case goes along the same line as above . Explicitly, the asymptotic expansion of the transfer matrix is given by:
$$T(x,\lambda )_{asy}=KG(x,\lambda )e^{_0^x𝑑yD(y)},$$
(2.18)
where $`K=\frac{\sqrt{2}}{2}\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right)`$ and
$$D(x,\lambda )=d(x,\lambda )H,d(x,\lambda )=\underset{k=1}{\overset{\mathrm{}}{}}\lambda ^kd_k(x)$$
(2.19)
($`d_{2k+1}`$ are the conserved densities). $`G`$ is given by:
$$G(x,\lambda )=\text{1}+\underset{j=1}{\overset{\mathrm{}}{}}\lambda ^jG_j(x),$$
(2.20)
where $`G_j(x)`$ are off-diagonal matrices with entries $`(G_j(x))_{12}=g_j(x)`$ and $`(G_j(x))_{21}=(1)^{j+1}g_j(x)`$ ( see ).
The resolvent (2.11) where now $`T=T_{asy}`$ and $`X=H`$ is an infinite series in negative powers of $`\lambda `$:
$$Z(x,\lambda )=\underset{k=0}{\overset{\mathrm{}}{}}\lambda ^kZ_k(x)$$
(2.21)
and obviously satisfies (2.12) by construction. The coefficients in (2.21) are given by:
$$Z_{2k}(x)=b_{2k}(x)E+c_{2k}(x)F,Z_{2k+1}(x)=a_{2k+1}(x)H.$$
(2.22)
A suitable gauge parameter in this case is constructed as:
$$\theta _n(x,\lambda )=(\lambda ^nZ(x;\lambda ))_+=\underset{j=0}{\overset{n}{}}\lambda ^{nj}Z_j(x)$$
(2.23)
and the additional consistency condition (2.14) implies that now the indices should be odd ( $`\theta _{2n+1}`$ ). It happens that these transformations coincide exactly with the commuting higher mKdV flows (or mKdV hierarchy):
$$\delta _{2k+1}\varphi ^{}(x)=a_{2k+1}(x)$$
(2.24)
and are therefore local in contrast with the regular ones. It turns out that the other entries of the resolvent $`b_{2n}(x)`$ are exactly the conserved densities , namely:
$$\delta _{2k+1}\varphi ^{}(x)=\{I_{2k+1},\varphi ^{}(x)\},I_{2k1}=_0^Ldxb_{2k}(x).$$
(2.25)
They differ from $`d_{2k+1}`$ (2.19) by a total derivative. For example:
$`b_2`$ $`=`$ $`d_1+\frac{1}{2}\varphi ^{\prime \prime },`$
$`b_4`$ $`=`$ $`{\displaystyle \frac{3}{4}}d_3+({\displaystyle \frac{7}{32}}\varphi ^{\prime \prime }\varphi ^{}+{\displaystyle \frac{1}{16}}(\varphi ^{})^3+{\displaystyle \frac{1}{16}}\varphi ^{\prime \prime \prime })etc.`$ (2.26)
Finally, let us note that it can be shown that these two kind of symmetries (regular and asymptotic) commute with each other. In this sence the non-local regular transformations provide a true symmetry of the KdV hierarchy.
### 2.2 Generalization - Virasoro Symmetry
Now, let us explain how the Virasoro symmetry appears in the KdV system . The main idea is that one can dress not only the generators of the underlying $`A_1^{(1)}`$ algebra but also an arbitrary differential operator in the spectral parameter. We take for example $`\lambda ^{m+1}_\lambda `$ which, as it is well known, are the vector fields of the diffeomorphisms of the unit circumference and close a Virasoro algebra. Then we proceed in the same way as above.
The analog of our basic object (2.11) now is:
$$Z^V(x,\lambda )=T(x,\lambda )_\lambda T(x,\lambda )^1.$$
(2.27)
It is clear that $`Z^V`$ has again the property of being a resolvent for the Lax operator, i.e. it satisfies (2.12), which was one of the requirements for constructing a good gauge parameter. Let us consider first the regular case, i.e. take $`T=T_{reg}`$ in (2.27):
$$Z_{reg}^V(x,\lambda )=\underset{n=0}{\overset{\mathrm{}}{}}\lambda ^nZ_n(x)+_\lambda .$$
(2.28)
It is clear that $`Z^V(x,\lambda )`$ is a differential operator in $`\lambda `$ in this case. Following the same reasoning as before we construct the gauge parameter as:
$$\theta _m^V(x,\lambda )=(\lambda ^mZ_{reg}^V(x,\lambda ))_{},m>0.$$
(2.29)
Then the additional condition (2.14) imposes that the indices of the transformation should be even $`m=2n`$. The first nontrivial examples are given by:
$`\delta _2^V\varphi ^{}`$ $`=`$ $`e^{2\varphi (x)}{\displaystyle _0^x}𝑑ye^{2\varphi (y)}e^{2\varphi (x)}{\displaystyle _0^x}𝑑ye^{2\varphi (y)}=e^{2\varphi (x)}B_1e^{2\varphi (x)}C_1`$
$`\delta _4^V\varphi ^{}`$ $`=`$ $`e^{2\varphi (x)}(3B_3(x)A_2(x)B_1(x))e^{2\varphi (x)}(3C_3(x)D_2(x)C_1(x))`$
$`\delta _6^V\varphi ^{}`$ $`=`$ $`e^{2\varphi (x)}(5B_5(x)3A_4(x)B_1(x)+A_2(x)B_3(x))`$ (2.30)
$``$ $`e^{2\varphi (x)}(5C_5(x)3D_4(x)C_1(x)+D_2(x)C_3(x)),`$
where $`A_i`$, $`B_i`$, $`C_i`$, $`D_i`$ were defined in (2.9). They have a form very similar to that of $`\delta _{2n}^H`$ but nevertheless it can be shown that they indeed close a (half) Virasoro algebra. Note that these transformations are essentially non-local, so we obtained a very nontrivial realisation of the Virasoro algebra in terms of vertex operators.
Let us now consider the asymptotic case, i.e. take $`T=T_{asy}`$ in (2.27):
$$Z_{asy}^V(x,\lambda )=\underset{n=0}{\overset{\mathrm{}}{}}\lambda ^nZ_n(x)+_\lambda .$$
(2.31)
The coefficients of the above expansion have the general form:
$$Z_{2n}=\beta _{2n}E+\gamma _{2n}F,Z_{2n+1}=\alpha _{2n+1}H,$$
where for example $`\beta _0=x=\gamma _0`$, $`\alpha _1=2xg_1`$, $`\beta _2=xb_2g_1+_0^xd_1`$, $`\gamma _2=xc_2+g_1+_0^xd_1`$ etc. . Again, the suitable gauge parameter is defined by:
$$\theta _m^V(x,\lambda )=(\lambda ^mZ_{asy}^V)_+=\underset{n=0}{\overset{m+1}{}}\lambda ^{m+1n}Z_n+_\lambda ,m0$$
(2.32)
and the self-consistency condition implies that the indices should be even in this case too. Actually, the first transformation:
$$\delta _0^V\varphi ^{}(x)=(x+1)\varphi ^{}(x)$$
(2.33)
is exactly the scale transformation - it counts the dimension (or level). The first non-trivial examples are:
$`\delta _2^V\varphi ^{}`$ $`=`$ $`2xa_3^{}(\varphi ^{})^3+{\displaystyle \frac{3}{4}}\varphi ^{\prime \prime \prime }+2a_1^{}{\displaystyle _0^x}d_1,`$
$`\delta _4^V\varphi ^{}`$ $`=`$ $`2xa_5^{}+(\varphi ^{})^5{\displaystyle \frac{5}{2}}\varphi ^{\prime \prime \prime }(\varphi ^{})^2{\displaystyle \frac{27}{8}}(\varphi ^{\prime \prime })^2\varphi ^{}+{\displaystyle \frac{5}{16}}\varphi ^V+`$ (2.34)
$`+`$ $`2a_3^{}{\displaystyle _0^x}d_1+6a_1^{}{\displaystyle _0^x}d_3.`$
We note that these depend explicitly on $`x`$ and are quasi-local (they contain some indefinite integrals). For further reference we presented the integrands in (2.34) explicitly in terms of the entries of the basic objects $`T(x,\lambda )`$ and $`Z(x,\lambda )`$, defined in (2.19,2.22). Furthermore, one can find the transformation of the resolvent and therefore the transformation of the conserved densities $`\delta _{2k}b_{2n}(x)`$. In particular the first nontrivial transformations of the KdV variable $`u=b_2`$ read:
$`\delta _2^Vb_2`$ $`=`$ $`\delta _2^Vu=2xb_4^{}+u^{\prime \prime }2u^2{\displaystyle \frac{1}{2}}u^{}{\displaystyle _0^x}u,`$
$`\delta _4^Vb_2`$ $`=`$ $`\delta _4^Vu=2xb_6^{}+2u^3+3uu^{\prime \prime }+{\displaystyle \frac{17}{8}}(u^{})^2+{\displaystyle \frac{3}{8}}u^{IV}+`$ (2.35)
$`+`$ $`u^{}{\displaystyle _0^x}b_4+b_4^{}{\displaystyle _0^x}u.`$
It can be easily shown that the asymptotic transformations also close (half) Virasoro algebra. A very non-trivial question concerns the commutation relations between the “negative” and “positive” parts so constructed in view of their completely different nature. Nevertheless it can be shown that, contrary to what happened between the proper dressing transformations and the m-KdV hierarchy, in this case they close a whole Virasoro algebra:
$$[\delta _{2m}^V,\delta _{2n}^V]=(2m2n)\delta _{2m+2n}^V,m,nZ.$$
(2.36)
We want to stress that the Virasoro symmetry just cconstructed is a dynamical one and has nothing to do with the space-time Virasoro symmetry of CFT. Actually, it is well known that one can consider the KdV system as a classical limit of the latter. So we expect that after quantization this dynamical symmetry should be present in CFT. It is interesting to investigate its significance, in particular if CFT could be solved by using this symmetry alone.
## 3 Soliton Solutions of (m)KdV theory
### 3.1 (m)KdV Solitons
We would like now to restrict the Virasoro symmetry to the soliton solutions of the (m)KdV theory. One can expect that in this case it simplifies considerably. There is also another reason for this restriction. It was shown in that one can reconstruct certain form-factors of sine-Gordon theory by directly quantizing the soliton solutions. Moreover, it happens that a kind of null-vectors appear in the theory , leading to integral equations for the form-factors. It is intriguing to understand the rôle that the Virasoro symmetry just described is playing in all these constructions.
We start with a brief description of the well known soliton solutions of (m)KdV. They are best expressed in terms of the so-called tau-function. In the case of N-soliton solution of (m)KdV it has the form:
$$\tau (X_1,\mathrm{},X_N|B_1,\mathrm{},B_N)=det(1+V)$$
(3.1)
where $`V`$ is a matrix:
$$V_{ij}=2\frac{B_iX_i(x)}{B_i+B_j},i,j=1,\mathrm{},N.$$
(3.2)
The m-KdV field is then expressed as:
$$e^\varphi =\frac{\tau _{}}{\tau _+},$$
(3.3)
where:
$$\tau _\pm (x)=\tau (\pm X(x)|B)$$
(3.4)
and $`X_i(x)`$ is simply given by:
$$X_i(x)=X_i\mathrm{exp}(2B_ix).$$
(3.5)
The variables $`B_i`$ and $`X_i`$ are the parameters describing the solitons: $`\beta _i=\mathrm{log}B_i`$ are the so-called rapidities and $`X_i`$ are related to the positions. The integrals of motion, restricted to the N-soliton solutions have the form
$$I_{2n+1}=\underset{i=1}{\overset{N}{}}B_i^{2n+1},n0.$$
(3.6)
It is well known that (m)KdV admits a non-degenerate symplectic structure. One can find the corresponding Poisson brackets between the basic variables $`B_i`$ and $`X_i`$ . The (m)KdV flows are then generated by (3.6) via
$$\delta _{2n+1}=\{\underset{i=1}{\overset{N}{}}B^{2n+1},\},n0.$$
(3.7)
### 3.2 Analytical Variables
Our final goal is the quantization of solitons and of the Virasoro symmetry. It was argued in that this is best performed in another set of variables $`\{A_i,B_i\}`$. The latter are the soliton limit of certain variables describing the more general quasi-periodic finite-zone solutions of (m)KdV. In that context $`B_i`$ are the branch points (i.e. define the complex structure) of the hyperelliptic Riemann surface describing the solution and $`A_i`$ are the zeroes of the so-called Baker-Akhiezer function defined on it. In view of the nice geometrical meaning of these variables they were called analytical variables in .
Explicitly, the change of variables is given by:
$$X_j\underset{kj}{}\frac{B_jB_k}{B_j+B_k}=\underset{k=1}{\overset{N}{}}\frac{B_jA_k}{B_j+A_k},j=1,\mathrm{},N.$$
(3.8)
The non-vanishing Poisson brackets expressed in terms of these new variables take the form:
$$\{A_i,B_j\}=\frac{\underset{ki}{}(B_j^2A_k^2)\underset{kj}{}(A_i^2B_k^2)}{_{ki}(A_i^2A_k^2)_{kj}(B_j^2B_k^2)}(A_i^2B_j^2).$$
(3.9)
The corresponding tau-functions have also a very compact form in terms of the analytical variables:
$`\tau _+`$ $`=`$ $`2^N{\displaystyle \underset{j=1}{\overset{N}{}}}B_j\{{\displaystyle \frac{\underset{i<j}{}(A_i+A_j)\underset{i<j}{}(B_i+B_j)}{_{i,j}(B_i+A_j)}}\}`$
$`\tau _{}`$ $`=`$ $`2^N{\displaystyle \underset{j=1}{\overset{N}{}}}A_j\{{\displaystyle \frac{\underset{i<j}{}(A_i+A_j)\underset{i<j}{}(B_i+B_j)}{_{i,j}(B_i+A_j)}}\}.`$ (3.10)
Therefore, from the explicit form of the m-KdV field in terms of the tau-functions (3.3) we obtain the following very simple expression:
$$e^\varphi \frac{\tau _{}}{\tau _+}=\underset{j=1}{\overset{N}{}}\frac{A_j}{B_j}.$$
(3.11)
The equation of motion of the $`A_i`$ variable is given by:
$$_xA_i\delta _1A_i=\{I_1,A_i\}=\underset{j=1}{\overset{N}{}}(A_i^2B_j^2)\underset{ji}{}\frac{1}{(A_i^2A_j^2)}.$$
(3.12)
One can verify that, as a consequence, the usual KdV variable $`u`$ is expressed as:
$$b_2u=\frac{1}{2}(\varphi ^{})^2+\frac{1}{2}\varphi ^{\prime \prime }=\underset{j=1}{\overset{N}{}}A_j^2\underset{j=1}{\overset{N}{}}B_j^2.$$
(3.13)
One can restrict also the higher KdV flows to the soliton solutions. For example it is clear from (3.7) that
$$\delta _{2n+1}B_i=0,n0.$$
(3.14)
This is a reminiscence of the fact that the KdV flows do not change the complex structure of the hyperelliptic surface describing the finite-zone solution . The variation of the $`A_i`$ variables can be easily computed as :
$$\delta _{2n+1}A_i=\{I_{2n+1},A_i\},n0$$
(3.15)
using the Poisson brackets (3.9).
### 3.3 Virasoro Symmetry of the Soliton Solutions
Now, we want to restrict the Virasoro symmetry of (m)KdV constructed above to the case of soliton solutions. In this section we shall be only interested in the positive part of the latter. The transformation of the rapidities can be easily deduced as a soliton limit of the Virasoro action on the finite-zone solutions described in :
$$\delta _{2n}B_i=B_i^{2n+1},n0,$$
(3.16)
i.e. the Virasoro action changes the complex structure (because of that it’s often called non-isospectral symmetry). What remains is to obtain the transformations of the $`A_i`$ variables. We found it quite difficult to deduce them as a soliton limit of the corresponding transformations of . Instead, we propose here another approach. Namely, we use the transformation of the fields $`\delta _{2n}\varphi `$, $`\delta _{2n}\varphi ^{}`$, $`\delta _{2n}u`$ etc. which we found before, restricted to the soliton solutions using (3.11,3.13). The problem is simplified by the fact that the Virasoro algebra is freely generated , i.e. we need to compute only the $`\delta _0`$, $`\delta _2`$ and $`\delta _4`$ transformations, the remaining ones are then obtained by commutation. In practice, we perform the computation for the first few cases of $`N=1,2,3`$ solitons and then proceed by induction.
Let us make an important observation. As we have stressed, the transformation of the basic objects in the field theory of (m)KdV are quasi-local – they contain certain indefinite integrals. It happens that the corresponding integrands become total derivatives when restricted to the soliton solutions. For example:
$`b_2`$ $``$ $`u=_x{\displaystyle \underset{i=1}{\overset{N}{}}}A_i(x),`$
$`b_4`$ $`=`$ $`_x{\displaystyle \underset{i=1}{\overset{N}{}}}A_i^3\frac{1}{2}u^{}_x[{\displaystyle \underset{i=1}{\overset{N}{}}}(A_i^3\frac{1}{2}_xA_i)].`$ (3.17)
Therefore the Virasoro transformations become local when restricted to the soliton solutions! The calculation is straightforward but quite tedious so we present here only the final result:
$`\delta _0A_i`$ $`=`$ $`(x_x+1)A_i,`$
$`\delta _2A_i`$ $`=`$ $`{\displaystyle \frac{1}{3}}x\delta _3A_i+A_i^3({\displaystyle \underset{j=1}{\overset{N}{}}}A_j)_xA_i,`$
$`\delta _4A_i`$ $`=`$ $`{\displaystyle \frac{1}{5}}x\delta _5A_i+A_i^5\{{\displaystyle \underset{ji}{}}A_i(A_i^2A_j^2)+{\displaystyle \underset{j=1}{\overset{N}{}}}A_j{\displaystyle \underset{k=1}{\overset{N}{}}}B_k^2\}_xA_i,`$ (3.18)
where the KdV flows read explicitly:
$`{\displaystyle \frac{1}{3}}\delta _3A_i`$ $`=`$ $`({\displaystyle \underset{j=1}{\overset{N}{}}}B_j^2{\displaystyle \underset{ki}{}}A_k^2)_xA_i,`$
$`{\displaystyle \frac{1}{5}}\delta _5A_i`$ $`=`$ $`({\displaystyle \underset{j=1}{\overset{N}{}}}B_j^4{\displaystyle \underset{ki}{}}A_k^4)_xA_i{\displaystyle \underset{ji}{}}(A_i^2A_j^2)_xA_i_xA_j.`$ (3.19)
As we already mentioned, the remaining transformations can be obtained by commutation, for example:
$$2\delta _6A_i=[\delta _4,\delta _2]A_i,etc.$$
(3.20)
## 4 Virasoro Symmetry of Sine-Gordon Solitons
### 4.1 From (m)KdV to Sine-Gordon Solitons
Now we pass to the most important part of our paper. We would like to extend the construction presented above in (m)KdV theory to the case of sine-Gordon. For this purpose one has to find a way of extending the mKdV dynamics up to the sine-Gordon one. It is to some extent known how this can be done in the case of the soliton solutions . The idea is close to what we proposed before directly in the field theory of (m)KdV. Namely, let us consider the KdV variable $`x`$ as a space variable of some more general system and call it $`x_{}`$ ( and $`_{}_x`$ correspondingly ). We would like to introduce a new time variable $`x_+`$ and the corresponding time dynamics. In the case of the N - soliton solutions the latter is generated by the Hamiltonian:
$$I_1=\underset{i=1}{\overset{N}{}}B_i^1$$
(4.1)
( essentially the inverse power of the momentum ) so that the time flow is given by:
$$_+=\delta _1=\{I_1,\}$$
(4.2)
using again the Poisson brackets (3.9). In particular:
$$_+A_i=\underset{j=1}{\overset{N}{}}\frac{A_i^2B_j^2)}{B_j^2}\underset{ji}{}\frac{A_j^2}{(A_i^2A_j^2)}.$$
(4.3)
One can check, using (3.11, 3.12), that with this definition the resulting equation for the field $`\varphi `$ is
$$_+_{}\varphi =2\mathrm{sinh}(2\varphi )$$
(4.4)
or under the change $`\varphi i\varphi `$:
$$_+_{}\varphi =2\mathrm{sin}(2\varphi )$$
(4.5)
i.e. the sine-Gordon equation. We were not able to establish at the moment a direct relation between the way of extending the mKdV dynamics to the sine-Gordon one directly in the field theory (2.16) and in the case of N-soliton solution just described (4.2). We hope to answer this important question elsewhere . In a similar manner one can introduce the rest of the sine-Gordon Hamiltonians:
$$I_{2n1}=\underset{i=1}{\overset{N}{}}B_i^{2n1},n0.$$
(4.6)
They generate the “negative KdV flows” via the Poisson brackets (3.9):
$`\delta _{2n1}B_i`$ $`=`$ $`0,`$
$`\delta _{2n1}A_i`$ $`=`$ $`\{I_{2n1},A_i\},n0.`$ (4.7)
### 4.2 Negative Virasoro flows
Now we arrive at the main conjecture of this paper. Having in mind the symmetric role the derivatives $`_{}`$ and $`_+`$ are playing in the sine-Gordon equation we would like to suppose that one can obtain another half Virasoro algebra by using the same construction as above but with $`_{}`$ interchanged with $`_+`$!
So let us define as before:
$$\delta _{2n}B_i=B_i^{2n+1},n0$$
(4.8)
(note the additional – sign in the r.h.s. which is needed for the self-consistency of the construction). Following our conjecture we construct the negative flows of the $`A_i`$ variable in the same way as before but with the change $`_{}_+`$. We have for example:
$`\delta _2\varphi `$ $`=`$ $`x_+(2a_3^+)+b_2^+2a_1^+{\displaystyle _0^{x_+}}b_2^+,`$
$`\delta _2b_2^{}`$ $``$ $`\delta _2u={\displaystyle \frac{1}{3}}x_+\delta _3u+(_+\varphi {\displaystyle _0^{x_+}}b_2^+)_{}e^{2\varphi },`$ (4.9)
where $`\delta _3u\{I_3,u\}`$ etc. In (4.9) the + subscript means that we take the same objects as defined in (2.19, 2.20, 2.22) but with $`_{}`$ changed by $`_+`$. For example:
$`b_2^+`$ $`=`$ $`\frac{1}{2}(_+\varphi )^2+\frac{1}{2}_+^2\varphi ,`$
$`a_3^+`$ $`=`$ $`{\displaystyle \frac{1}{4}}(_+\varphi )^3+{\displaystyle \frac{1}{8}}_+^3\varphi etc.`$ (4.10)
At this point we want to make an important remark. Very non-trivially, it happens again that the integrands in the expressions (4.9) and similar become total derivatives when restricted to the N-soliton solutions. So that again the (negative) Virasoro symmetry is local in the case of solitons! We present below the first examples of this phenomenon:
$`b_2^+`$ $`=`$ $`_+\{{\displaystyle \underset{i,j=1}{\overset{N}{}}}{\displaystyle \frac{A_iA_j}{B_i^2B_j^2}}{\displaystyle \underset{i=1}{\overset{N}{}}}A_i_{}{\displaystyle \underset{i,j=1}{\overset{N}{}}}{\displaystyle \frac{A_iA_j}{B_i^2B_j^2}}\},`$
$`b_4^+`$ $`=`$ $`_+\{{\displaystyle \underset{i,j=1}{\overset{N}{}}}{\displaystyle \frac{A_iA_j}{B_i^4B_j^4}}{\displaystyle \underset{i=1}{\overset{N}{}}}A_i^3+b_2^{}_{}{\displaystyle \underset{i,j=1}{\overset{N}{}}}{\displaystyle \frac{A_iA_j}{B_i^4B_j^4}}`$ (4.11)
$``$ $`_{}b_2^{}{\displaystyle \underset{i,j=1}{\overset{N}{}}}{\displaystyle \frac{A_iA_j}{B_i^4B_j^4}}\}etc.`$
We then proceed as in the case of the positive Virasoro flows, i.e. we restrict the transformations of the fields thus obtained to the soliton solutions. As we explained, it is enough to find only the first transformations $`\delta _2A_i`$ and $`\delta _4A_i`$ and the remaining ones are found by commutation. Following our approach we do the computation explicitly in the case of $`N=1,2,3`$ solitons and then proceed by induction. The exact calculation will be presented in a forthcoming paper , here we give the final results only:
$`\delta _2A_i`$ $`=`$ $`{\displaystyle \frac{1}{3}}x_+\delta _3A_iA_i^1({\displaystyle \underset{j=1}{\overset{N}{}}}A_j^1)_+A_i,`$
$`\delta _4A_i`$ $`=`$ $`{\displaystyle \frac{1}{5}}x_+\delta _5A_iA_i^3`$ (4.12)
$``$ $`\{{\displaystyle \underset{ji}{\overset{N}{}}}{\displaystyle \frac{1}{A_i}}({\displaystyle \frac{1}{A_i^2}}{\displaystyle \frac{1}{A_j^2}})+{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{1}{A_j}}{\displaystyle \underset{k=1}{\overset{N}{}}}{\displaystyle \frac{1}{B_k^2}}\}_+A_i,`$
where as before $`\delta _3A_i=\{_{j=1}^NB_j^3,A_i\}`$ etc. As stated above, we then can compute $`2\delta _6A_i=[\delta _2,\delta _4]A_i`$ etc.
### 4.3 The Algebra
Now, we come to the important problem of the commutation relations between the two half Virasoro algebras so constructed. This is a non-trivial question in view of the different way we obtained them. In fact, it is clear that, by construction, the positive (negative) Virasoro flows commute with the corresponding $`_{}`$ ( $`_+`$ ) derivatives:
$`[\delta _{2n},_{}]A_i`$ $`=`$ $`0,`$
$`[\delta _{2n},_+]A_i`$ $`=`$ $`0,n0.`$ (4.13)
It is easy to see that this is not true for the “cross commutators”. Actually, one finds in this case:
$`[\delta _{2n},_+]A_i`$ $`=`$ $`\delta _{2n1}A_i,`$
$`[\delta _{2n},_{}]A_i`$ $`=`$ $`\delta _{2n+1}A_i,n0.`$ (4.14)
It is clear that we are interested in a true symmetry of the sine-Gordon theory. We must therefore obtain transformations that commute with the $`_{}`$ and $`_+`$ flows and as a consequence with the corresponding Hamiltonians. It is obvious from (4.13,4.14) that this is acheeved by a simple modification of the flows, i.e. let us define:
$`\delta _{2n}^{}`$ $`=`$ $`\delta _{2n}x_+\delta _{2n1},`$
$`\delta _{2n}^{}`$ $`=`$ $`\delta _{2n}x_{}\delta _{2n+1},n0.`$ (4.15)
Then, for the modified transformation we obtain:
$$[\delta _{2n}^{},_\pm ]A_i=0,nZ.$$
(4.16)
Finally, one can show that, with this modification, the commutation relations between the positive and negative parts of the transformations close exactly the whole Virasoro algebra:
$$[\delta _{2n}^{},\delta _{2m}^{}]A_i=(2n2m)\delta _{2n+2m}^{}A_i,n,mZ.$$
(4.17)
## 5 Conclusions and discussion
To summarise, we presented in this paper a construction of a Virasoro symmetry of the sine-Gordon theory. This is acheeved by extending the corresponding symmetry of the (m)KdV theory . Actually, we found it easier to work rather with the N-soliton solutions of the latter. Then, following , we introduced a time coordinate $`x_+`$ and the corresponding flow $`_+`$ in addtion to the space variable $`x_{}`$ of KdV which leads to the sine-Gordon equation. The main idea is to make the same kind of construction as for KdV but interchanging the $`_{}`$ with the $`_+`$ derivative. This change results in what we called “negative” Virasoro flows which complete the “positive” ones coming from the original KdV to the whole Virasoro algebra. We showed also that after a certain modification, needed to obtain a true symmetry of the theory, they close the whole Virasoro algebra.
Actually, in this paper we obtained the infinitesimal transformations of the variables describing the N-soliton solutions. It is intriguing to find the corresponding conserved charges $`J_{2n}(A_i,B_i)`$. This is important in view of the quantization of the classical constructions presented here. Besides, in their commutation relations one can obtain a possible existence of a central extension which cannot be discovered in the case of the infinitesimal transformations.
As already mentioned, we are interested of course in the quantum sine-Gordon theory. In the case of solitons there is a standard procedure, a kind of canonical quantization of the N-soliton solutions. In fact, let us introduce, following , the canonically conjugated variables to the analytical variables $`A_i`$:
$$P_j=\underset{k=1}{\overset{N}{}}\frac{B_kA_j}{B_k+A_j},j=1\mathrm{},N$$
(5.1)
(in the variables $`\{P_j,A_j\}`$ the corresponding simplectic structure is diagonal). In these variables one can perform a kind of canonical quantization of the N-soliton system introducing the deformed commutation relations between the operators $`A_i`$ and $`P_i`$ :
$`P_jA_j`$ $`=`$ $`qA_jP_j,`$
$`P_kA_j`$ $`=`$ $`A_jP_kforkj,`$ (5.2)
where $`q`$ is related to the sine-Gordon coupling constant: $`\mathrm{exp}(i\xi )`$, $`\xi =\frac{\pi \gamma }{\pi \gamma }`$. It is very intriguing to understand how the Virasoro symmetry is deformed after the quantization!
Another important problem is the construction of the Virasoro symmetry directly in the sine-Gordon field theory. We explained above how this can be done using the proper dressing transformations. It happens that the positive Virasoro flows commute with the time sine-Gordon flow $`x_+`$ introduced in this way. We expect that the negative Virasoro symmetry can be constructed following the same approach we presented in this paper for the solitons.
Finally, it will be very interesting to understand the rôle this Virasoro symmetry is playing in the sine-Gordon theory. As we mentioned, it was shown in that certain form-factors can be reconstructed by a suitable quantization of the N-soliton solutions. One can expect that the Virasoro symmetry imposes some constraints leading to certain equations for the form-factors (or correlation functions in the case of field theory). We would like to note in this respect the article where a kind of null-vector constraints were derived in the sine-Gordon theory. The corresponding construction is closely related to the finite-zone solutions. As we mentioned the Virasoro symmetry we presented here has a natural action on such solutions. It is natural to expect that some relation exists between this symmetry and the null-vectors of . We will return to all these problems in a forthcoming paper .
Acknowledgments \- We are indebted to E. Corrigan, V. Fateev, G. Mussardo, G. Sotkov and Al. Zamolodchikov for discussions and interest in this work. D.F. thanks the I.N.F.N.–S.I.S.S.A. for financial support. M.S. acknowledges S.I.S.S.A. and the University of Montpellier -2 for the warm hospitality over part of this work. This work has been realized through partial financial support of TMR Contract ERBFMRXCT960012.
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# 1 Introduction
## 1 Introduction
Orbifold theory$`^{\text{References}\text{References}}`$ has a long history, yet until recently, orbifolds have been studied primarily at the level of examples. This situation has now changed due to a recent synthesis of the principles of orbifold theory with the principles of current-algebraic conformal field theory, and we may now view at a glance the panorama of all current-algebraic orbifolds.$`^{\text{References}}`$
In particular, Ref. References gave a construction<sup>a</sup><sup>a</sup>aThis construction drew heavily on recent advances in the theory of cyclic permutation orbifolds, including the discovery of orbifold affine algebra$`^{\text{References}}`$ and the orbifold Virasoro master equation.$`^{\text{References}}`$ of the twisted currents $`\widehat{J}(\sigma )`$ and stress tensors $`\widehat{T}_\sigma `$
$$T=L_H^{ab}:J_aJ_b:\begin{array}{c}\hfill \\ ^\sigma \hfill \end{array}\widehat{T}_\sigma =^{n(r)\mu ;n(r),\nu }(L_H;\sigma ):\widehat{J}(\sigma )_{n(r)\mu }\widehat{J}(\sigma )_{n(r),\nu }:$$
(1.1a)
$$\sigma =0,\mathrm{},N_c1$$
(1.1b)
of all sectors $`\sigma `$ of any current-algebraic orbifold $`A(H)/H`$. Here $`A(H)`$, described by the stress tensor $`T`$, is any current-algebraic conformal field theory with a finite symmetry group $`H`$. Technically, $`A(H)`$ is a member of the class of $`H`$-invariant CFT’s$`^{\text{References}\text{References},\text{References}}`$ on $`g`$, which includes all the CFT’s with a symmetry $`HAut(g)`$ in the general affine-Virasoro construction.$`^{\text{References},\text{References},\text{References}}`$ The number of sectors $`N_c`$ of the orbifold $`A(H)/H`$ is the number of conjugacy classes of $`H`$ and the construction (a) is shown schematically in Fig. 1.
The orbifold duality transformation indicated in (a)
$$L_H\begin{array}{c}\hfill \\ ^\sigma \hfill \end{array}(L_H;\sigma )$$
(1.2)
gives the twisted inverse inertia tensor $``$ of each sector $`\sigma `$ in terms of the $`H`$-invariant inverse inertia tensor $`L_H`$ of the $`H`$-invariant CFT $`A(H)`$. Other orbifold duality transformations exist for other twisted tensors of the orbifold, and the explicit form of the generic orbifold duality transformation is a discrete Fourier transform.
The central ingredients underlying the breadth and depth of this construction are
$``$ local formulation of the theory in terms of currents, OPE’s and OPE isomorphisms
$``$ the $`L_H^{ab}`$ formulation of the $`H`$-invariant CFT’s in the general affine-Virasoro construction.
The local formulation (as opposed to a mode formulation) is the key to the orbifold duality transformations, while the $`L_H^{ab}`$ formulation allows us to study all current-algebraic orbifolds at once, or any example.
The organization of this paper is best summarized by the names of its sections:
2 Local Formulation of Current-Algebraic Orbifolds
3 The Mode Formulation of Orbifold Theory
4 The Orbifold Adjoint Operation
5 The $`\widehat{T}\widehat{J}`$ OPE’s of $`A(H)/H`$
6 The Orbifolds of the ($`H`$ and Lie $`h`$)-invariant CFT’s
7 About Permutation Orbifolds
8 The Permutation Orbifolds $`A(S_N)/S_N`$
9 The Inner-Automorphic Orbifolds $`A(H(d))/H(d)`$.
In particular, Sec. 2 reviews and extends the local formulation$`^{\text{References}}`$ of current-algebraic orbifold theory in a form which exhibits the natural grading of the orbifold operator products. In Secs. 35, we work out the implied mode formulation and other consequences of the local formulation. We mention in particular four results obtained in the mode formulation for all sectors $`\sigma `$ of all orbifolds $`A(H)/H`$: the general twisted current algebra in Eq. (3.1), the orbifold Virasoro generators in Eq. (3.12), the orbifold adjoint operation in Eqs. (4.2) and (4.7) and the commutator (5.9) of the Virasoro generators with the modes of the twisted currents.
The ($`H`$ and Lie $`h`$)-invariant CFT’s $`A(\text{Lie}h(H))`$ are those “doubly-invariant” CFT’s with both a finite symmetry and a Lie symmetry$`^{\text{References}\text{References},\text{References},\text{References}}`$ although we mod out only by the finite symmetry to obtain (see Sec. 6) the orbifolds $`A(\text{Lie}h(H))/H`$. The general WZW orbifold and the general coset orbifold are discussed as special cases of $`A(\text{Lie}h(H))/H`$ in Subsec. 6.4. The twisted current-current correlators and ground state conformal weights of the general permutation orbifold are computed in Sec. 7. We also work out two new large examples in further detail, including the general $`S_N`$ permutation orbifold in Sec. 8 and the general inner-automorphic orbifold in Sec. 9. The story of the inner-automorphic orbifolds $`A(H(d))/H(d)`$ is particularly interesting, not least because of their overlap with the orbifolds $`A(\text{Lie}h(H))/H`$. Indeed, we will argue that this overlap contains almost all the inner-automorphic orbifolds which can be equivalently described by stress-tensor spectral flow$`^{\text{References},\text{References},\text{References},\text{References}}`$ whereas the generic inner-automorphic orbifold apparently can not be described in this way.
The seminal case$`^{\text{References}}`$ of the cyclic permutation orbifolds $`A(_\lambda )/_\lambda `$ is used as an example in various sections, and the appendices include the setups for the permutation orbifolds $`A(𝔻_\lambda )/𝔻_\lambda `$ and the outer-automorphic orbifolds.
This is as far as we have worked out the consequences of the “orbifold program”, but much remains to be done, including other large examples and analysis of the situation when $`A(H)`$ has a larger chiral algebra.
## 2 Local Formulation of Current-Algebraic Orbifolds
In this section we review and extend the local formulation of the general current-algebraic orbifold$`^{\text{References}}`$. The extension includes a new periodic notation and new “selection rules”, both of which are necessary to exhibit the natural grading of the orbifold operator products. As an illustration, the seminal case of the general cyclic permutation orbifold$`^{\text{References}}`$ is recalled in Subsec. 2.6.
### 2.1 The current-algebraic CFT’s
The study of current-algebraic CFT’s begins with the general affine Lie algebra$`^{\text{References},\text{References},\text{References},\text{References}}`$
$$J_a(z)J_b(w)=\frac{G_{ab}}{(zw)^2}+\frac{if_{ab}^cJ_c(w)}{zw}+O((zw)^0),J_a(ze^{2\pi i})=J_a(z)$$
(2.1a)
$$[J_a(m),J_b(n)]=if_{ab}^cJ_c(m+n)+mG_{ab}\delta _{m+n,0}$$
(2.1b)
$$g=_Ig^I,G_{ab}=_Ik_I\eta _{ab}^I,a,b,c=1,\mathrm{},\text{dim}g,m,n$$
(2.1c)
where $`G_{ab}`$ and $`f_{ab}^c`$ are the metric and structure constants of semisimple<sup>b</sup><sup>b</sup>bNon-compact versions of $`g`$ are included and abelian components of $`g`$ are easily included$`^{\text{References}}`$ as well (see Subsec. 3.3). Lie $`g`$. The desired real form of the affine algebra is specified by the adjoint operation
$$J_a(m)^{}=\rho _a^bJ_b(m),\rho _a^c\rho _c^b=\delta _a^b,\rho ^{}\rho =1$$
(2.2a)
$$G_{ab}^{}=\rho _a^c\rho _b^dG_{cd},f_{ab}^c=\rho _a^d\rho _b^ef_{de}^f\rho _f^c$$
(2.2b)
where the conjugation matrix $`\rho `$ is unity in any Cartesian basis, and corresponding forms follow in other bases. For integer level $`x_I=2k_I/\psi _I^2`$ of compact $`g^I`$ with highest root $`\psi _I`$, the adjoint operation in (2.2a) guarantees unitarity of the affine Hilbert space.
The general affine-Virasoro construction$`^{\text{References},\text{References},\text{References},\text{References},\text{References}}`$ is
$$T(z)T(w)=\frac{c/2}{(zw)^4}+(\frac{2}{(zw)^2}+\frac{_w}{zw})T(w)+O((zw)^0)$$
(2.3a)
$$T(z)=L^{ab}:J_a(z)J_b(z):,L^{ab}=L^{ba}$$
(2.3b)
$$L^{ab}=2L^{ac}G_{cd}L^{db}L^{cd}L^{ef}f_{ce}^af_{df}^bL^{cd}f_{ce}^ff_{df}^{(a}L^{b)e},c=2G_{ab}L^{ab}$$
(2.3c)
where we have chosen<sup>c</sup><sup>c</sup>cBecause of the symmetry $`L^{ab}=L^{ba}`$, the stress tensor $`T`$ in (2.3b) can equivalently be described by various normal ordering prescriptions. OPE normal ordering however comes equipped with an efficient calculus$`^{\text{References}}`$ for computing OPE’s of products of ordered products and this prescription has been extended$`^{\text{References},\text{References}}`$ to compute corresponding OPE’s in the twisted sectors of orbifolds (see Subsec. 2.5). OPE normal ordering
$$:J_a(w)J_b(w):=_w\frac{dz}{2\pi i}\frac{J_a(z)J_b(w)}{zw}$$
(2.4)
and the inverse inertia tensor $`L^{ab}`$ satisfies the Virasoro master equation in (2.3c). The $`L^{}`$ relation given below$`^{\text{References}}`$
$$L^{ab}=L^{cd}\rho _c^a\rho _d^bL(m)^{}=L(m)$$
(2.5)
guarantees the indicated adjoint of the Virasoro generators, and this adjoint operation guarantees the unitarity of the CFT, given the unitarity of the affine Hilbert space.
### 2.2 The automorphism group $`H`$ and the $`H`$-invariant CFT’s
We consider any $`H`$-covariant algebra $`g`$, that is, any algebra $`g`$ with a finite-order<sup>d</sup><sup>d</sup>dAutomorphism groups of infinite order (including Lie groups) are also included formally. See our remarks in Subsecs. 6.1 and 9.1. automorphism group $`HAut(g)`$. The automorphism group $`H`$ acts on the $`g`$ currents $`J`$ in matrix representation $`\omega `$,
$$J_a(z)^{}=\omega (h_\sigma )_a^bJ_b(z),\omega (h_\sigma )HAut(g)$$
(2.6a)
$$\omega (h_\sigma )_a^c\omega ^{}(h_\sigma )_c^b=\delta _a^b$$
(2.6b)
$$\omega (h_\sigma )_a^c\omega (h_\sigma )_b^dG_{cd}=G_{ab},\omega (h_\sigma )_a^d\omega (h_\sigma )_b^ef_{de}^f\omega ^{}(h_\sigma )_f^c=f_{ab}^c$$
(2.6c)
$$\omega (h_\sigma )^{}\rho \omega ^{}(h_\sigma )=\rho J_a(m)^{}{}_{}{}^{}=J_a(m)^{}^{}$$
(2.6d)
where the relations in (d) hold for all $`h_\sigma H`$. The relations (2.6c) – which express the $`H`$-invariance of $`G_{ab}`$ and $`f_{ab}^c`$ – guarantee that $`J^{}`$ satisfies the same algebra as $`J`$, so that $`HAut(g)`$ is also an automorphism group of affine $`g`$. Consequently, we often refer to the affine algebra as an $`H`$-covariant affine algebra on $`g`$. The $`H`$-covariance of the conjugation matrix $`\rho `$ in (2.6d) is equivalent to the statement that the automorphism group preserves the real form of the affine algebra (2.1b).
The $`H`$-invariant CFT’s$`^{\text{References}\text{References},\text{References}}`$ $`A(H)`$ are those CFT’s for which the inverse inertia tensor $`L_H`$ is invariant under $`H`$
$$T=L_H^{ab}:J_aJ_b:,L_H^{cd}\omega (h_\sigma )_c^a\omega (h_\sigma )_d^b=L_H^{ab},h_\sigma HAut(g)$$
(2.7a)
$$T^{}=L_H^{ab}:J_a^{}J_b^{}:=L_H^{ab}:J_aJ_b:=T.$$
(2.7b)
These are the CFT’s whose stress tensors also describe the untwisted sectors of the general current-algebraic orbifold $`A(H)/H`$. For each symmetry group $`H`$, it is known$`^{\text{References}}`$ that the inverse inertia tensors of the $`H`$-invariant CFT’s satisfy a consistent reduced Virasoro master equation (with an equal number of equations and unknowns).
The simplest $`H`$-invariant CFT’s include $`A_g(H)`$, $`HAut(g)`$, whose stress tensor is the affine-Sugawara construction$`^{\text{References},\text{References}\text{References},\text{References}}`$ on $`g`$, and the general $`H`$-invariant coset construction$`^{\text{References},\text{References}}`$ $`\frac{g}{h}(H)`$. We will return to these popular examples as illustrations below (see Subsecs. 6.16.4 and 9.3).
In what follows all inverse inertia tensors $`L_H`$ are $`H`$-invariant, so we may drop the subscript $`L_H^{ab}L^{ab}`$ without confusion.
### 2.3 The $`H`$-eigenvalue problem
Choosing one representative $`h_\sigma `$ from each conjugacy class $`\sigma `$ of $`H`$, the $`H`$-eigenvalue problem is defined as
$$\omega (h_\sigma )_a^bU^{}(\sigma )_b^{n(r)\mu }=U^{}(\sigma )_a^{n(r)\mu }E_{n(r)}(\sigma ),\sigma =0,\mathrm{},N_c1$$
(2.8a)
$$U^{}(\sigma )U(\sigma )=1,E_{n(r)}(\sigma )=e^{\frac{2\pi in(r)}{\rho (\sigma )}},n(r).$$
(2.8b)
Here $`N_c`$ is the number of conjugacy classes of $`H`$ and the integer $`r`$ runs over the distinct eigenvalues of the problem at each $`\sigma `$. The eigenvalues $`E_{n(r)}(\sigma )`$ are labeled by the spectral index $`n(r)n(r;\sigma )`$ and the degeneracy index $`\mu \mu (n(r;\sigma ))`$ labels the degenerate eigenstates with eigenvalue $`E_{n(r)}(\sigma )`$. The quantity $`\rho (\sigma )^+`$ is the order of the automorphism $`h_\sigma `$, defined as the smallest positive integer satisfying
$$\omega (h_\sigma )^{\rho (\sigma )}=1.$$
(2.9)
As a convention, we assign $`\sigma =0`$ to the trivial automorphism
$$\omega (h_0)_a^b=\delta _a^b,U(0)=U^{}(0)=1,E_{n(r)}(0)=\rho (0)=1.$$
(2.10)
Other useful forms of the eigenvalue problem include
$$\omega =U^{}EU,\omega ^{}=U^{}E^{}U$$
(2.11a)
$$\omega U^{}=U^{}E,U\omega =EU,U^{}\omega ^{}=E^{}U^{},\omega ^{}U^{}=U^{}E^{},U\omega ^{}=E^{}U$$
(2.11b)
where $`\text{diag}(E)=\{E_{n(r)}\}`$. As seen here, we often suppress the label $`\sigma `$ for brevity.
The eigenvalues $`E_{n(r)}`$ are invariant under the transformation $`n(r)n(r)\pm \rho (\sigma )`$, so we may require this periodicity for the eigenvectors as well,
$$E_{n(r)\pm \rho (\sigma )}(\sigma )=E_{n(r)}(\sigma )$$
(2.12a)
$$U^{}(\sigma )_a^{n(r)\pm \rho (\sigma ),\mu }=U^{}(\sigma )_a^{n(r)\mu },U(\sigma )_{n(r)\pm \rho (\sigma ),\mu }^a=U(\sigma )_{n(r)\mu }^a$$
(2.12b)
$$\underset{r,\mu }{}U^{}(\sigma )_a^{n(r)\mu }U(\sigma )_{n(r)\mu }^b=\delta _a^b,U(\sigma )_{n(r)\mu }^aU^{}(\sigma )_a^{n(s)\nu }=\delta _\mu ^\nu \delta _{n(r)}^{n(s)}$$
(2.12c)
$$\delta _{n(r)}^{n(s)}\delta _{n(r)n(s),0\text{ mod }\rho (\sigma )},\underset{s}{}f(n(s))\delta _{n(s)}^{n(r)}=f(n(r))$$
(2.12d)
where the last identity holds for all $`f`$ with period $`\rho (\sigma )`$. In what follows, relations such as the first relation in (2.12c) will generally be written as
$$U^{}(\sigma )_a^{n(r)\mu }U(\sigma )_{n(r)\mu }^b=\delta _a^b$$
(2.13)
with an implied summation on repeated indices.
Using the $`H`$-covariance of $`\rho `$ in (2.6d), the eigenvalue problem can be recast in the form
$$\omega (\rho U^{})^{}=(\rho U^{})^{}E^{}.$$
(2.14)
This tells us first that the eigenvalues $`E^{}`$ are in the spectrum $`E`$
$$E_{n(r)}(\sigma )^{}=E_{\rho (\sigma )n(r)}(\sigma )=E_{n(r)}(\sigma )$$
(2.15)
and, in particular, that $`n(r)\{n(r)\}`$ is a spectral index when $`n(r)`$ is a spectral index. Moreover, (2.15) tells us that the degeneracy of $`E_{n(r)}(\sigma )`$ is the same as that of $`E_{n(r)}(\sigma )`$.
Since $`n(r)`$ is only determined mod $`\rho (\sigma )`$, it is useful to define the twist class $`\overline{n}(r)`$
$$\overline{n}(r)\overline{n(r)}n(r)\rho (\sigma )\frac{n(r)}{\rho (\sigma )},\overline{n}(r)\{0,\mathrm{},\rho (\sigma )1\}$$
(2.16)
where $`x`$ is the largest integer less than or equal to $`x`$. The twist class is an integer which evaluates $`n(r)`$ in the fundamental range shown, and twist classes will control the monodromies of twisted fields in the orbifolds below. Note that $`\overline{n}(r)=n(r)`$ when $`n(r)`$ is in the fundamental range. We also compute the twist class corresponding to $`n(r)`$
$`\overline{n(r)}=\overline{\overline{n}(r)}=\{\begin{array}{cc}\rho (\sigma )\overline{n}(r)\text{ when }\overline{n}(r)0& \\ 0\text{ when }\overline{n}(r)=0& \end{array}`$ (2.19)
where we have used $`x=x+1`$ for $`x`$ not an integer.
In what follows, the description of each sector $`\sigma `$ of the general current-algebraic orbifold $`A(H)/H`$ is given in terms of the solution $`\{U^{}(\sigma )`$, $`E(\sigma )\}`$ of the $`H`$-eigenvalue problem (2.3). It is unlikely that a solution of the $`H`$-eigenvalue problem can be found across all $`H`$, but solution of the $`H`$-eigenvalue problem is straightforward for any particular choice of $`H`$: The solutions for the cyclic permutation groups $`H=_\lambda `$
$$g=_{I=0}^{\lambda 1}𝔤^I,𝔤^I𝔤$$
(2.20a)
$$aa,I,a=1,\mathrm{},\text{dim}𝔤,I=0,\mathrm{},\lambda 1$$
(2.20b)
$$\omega (h_\sigma )_a^b\omega (h_\sigma )_{aI}^{bJ}=\delta _a^b\delta _{I+\sigma ,J\text{ mod }\lambda },\sigma =0,\mathrm{},\lambda 1$$
(2.20c)
$$n(r),\mu r,aj,\overline{r}=0,\mathrm{},\rho (\sigma )1,j=0,\mathrm{},\frac{\lambda }{\rho (\sigma )}1$$
(2.20d)
$$U^{}(\sigma )_a^{n(r)\mu }U^{}(\sigma )_{aI}^{rbj}=\frac{\delta _a^b}{\sqrt{\rho (\sigma )}}e^{\frac{2\pi iN(\sigma )r(jI)}{\lambda }}\delta _{j,I\text{ mod }\frac{\lambda }{\rho (\sigma )}},E_r(\sigma )=e^{\frac{2\pi ir}{\rho (\sigma )}}$$
(2.20e)
were given in Ref. References, which also defines the integers $`N(\sigma )`$. In this case, the labels $`(a,j)`$ are the degeneracy indices of the spectral indices $`n(r)=r`$, with $`\overline{r}=r\rho (\sigma )r/\rho (\sigma )`$. The $`U^{}`$ periodicity $`rr\pm \rho (\sigma )`$ is a consequence of the support $`((jI)\rho /\lambda )`$ of the Kronecker factor. More generally, the eigenvectors of each $`H`$-eigenvalue problem are the basis elements of discrete Fourier transforms with period $`\rho (\sigma )`$ in the spectral indices $`\{n(r)\}`$.
The solutions for the permutation group $`H=S_N`$ and the general group of inner automorphisms of simple $`g`$ will be given in Secs. 8 and 9 respectively. The setups for the permutation group $`H=𝔻_\lambda `$ and the general group of outer automorphisms of simple $`g`$ are given in Apps. Appendix I. and Appendix J. respectively. Many other large examples remain to be worked out in further detail, including the other permutation subgroups of $`S_N`$ (see also Subsec. 3.1).
### 2.4 The twisted currents of the orbifold $`A(H)/H`$
We turn next to the twisted currents $`\widehat{J}`$ of the current-algebraic orbifold
$$\frac{A(H)}{H},HAut(g)$$
(2.21)
where $`A(H)`$ is any $`H`$-invariant CFT.
For each conjugacy class $`\sigma `$, the eigencurrent$`^{\text{References}}`$ $`𝒥(\sigma )`$ is defined as the linear combination of the untwisted currents $`J`$
$$𝒥_{n(r)\mu }(z)\chi \left(\sigma \right)_{n\left(r\right)\mu }U(\sigma )_{n(r)\mu }^aJ_a(z),J_a(z)=\chi \left(\sigma \right)_{n\left(r\right)\mu }^1U^{}(\sigma )_a^{n(r)\mu }𝒥_{n(r)\mu }(z)$$
(2.22a)
$$𝒥_{n(r)\mu }(z)^{}=E_{n(r)}(\sigma )𝒥_{n(r)\mu }(z)$$
(2.22b)
which (according to (2.6a) and (2.10)) has a diagonal response $`E_{n(r)}(\sigma )`$ to the automorphism group. The numbers $`\chi \left(\sigma \right)_{n\left(r\right)\pm \rho \left(\sigma \right),\mu }=\chi \left(\sigma \right)_{n\left(r\right)\mu }`$ with $`\chi \left(0\right)=1`$ comprise an otherwise arbitrary set of normalization constants.
In general orbifold theory$`^{\text{References},\text{References},\text{References},\text{References},\text{References},\text{References}}`$ the sectors of each orbifold $`A(H)/H`$ are labeled by the conjugacy classes $`\sigma =0,\mathrm{},N_c1`$ of $`H`$. The twisted currents $`\widehat{J}(\sigma )`$ of sector $`\sigma `$ are obtained from the eigencurrents of sector $`\sigma `$ by the OPE isomorphism$`^{\text{References},\text{References}}`$
$$𝒥_{n(r)\mu }(z)\begin{array}{c}\hfill \\ ^\sigma \hfill \end{array}\widehat{J}_{n(r)\mu }(z)$$
(2.23a)
$$\text{automorphisms }E_{n(r)}(\sigma )\begin{array}{c}\hfill \\ ^\sigma \hfill \end{array}\text{monodromies }E_{n(r)}(\sigma )$$
(2.23b)
which includes the automorphism/monodromy exchange in (2.23b).
The OPE isomorphism (2.4) gives the twisted current system$`^{\text{References}}`$ of sector $`\sigma `$
$$\widehat{J}_{n(r)\mu }(z)\widehat{J}_{n(s)\nu }(w)=\frac{𝒢_{n(r)\mu ;n(s)\nu }(\sigma )}{(zw)^2}+\frac{i_{n(r)\mu ;n(s)\nu }^{n(t)\delta }(\sigma )\widehat{J}_{n(t)\delta }(w)}{zw}+O((zw)^0)$$
(2.24a)
$$\widehat{J}_{n(r)\mu }(ze^{2\pi i})=E_{n(r)}(\sigma )\widehat{J}_{n(r)\mu }(z),E_{n(r)}(\sigma )=e^{\frac{2\pi in(r)}{\rho (\sigma )}}$$
(2.24b)
$$\widehat{J}_{n(r)_\pm \rho (\sigma ),\mu }(z)=\widehat{J}_{n(r)\mu }(z),\widehat{J}_{\rho (\sigma )n(r),\mu }(z)=\widehat{J}_{n(r),\mu }(z)$$
(2.24c)
$$𝒢_{n(r)\mu ;n(s)\nu }(\sigma )=\chi \left(\sigma \right)_{n\left(r\right)\mu }\chi \left(\sigma \right)_{n\left(s\right)\nu }U(\sigma )_{n(r)\mu }^aU(\sigma )_{n(s)\nu }^bG_{ab}$$
(2.24d)
$$_{n(r)\mu ;n(s)\nu }^{n(t)\delta }(\sigma )=\chi \left(\sigma \right)_{n\left(r\right)\mu }\chi \left(\sigma \right)_{n\left(s\right)\nu }\chi \left(\sigma \right)_{n\left(t\right)\delta }^1U(\sigma )_{n(r)\mu }^aU(\sigma )_{n(s)\nu }^bf_{ab}^cU^{}(\sigma )_c^{n(t)\delta }$$
(2.24e)
$$\mathrm{\#}\{\widehat{J}(\sigma )\}=\mathrm{\#}\{J\}=\text{dim}g$$
(2.24f)
and (2.24b) tells us that the twisted current $`\widehat{J}_{n(r)\mu }`$ has twist class $`\overline{n}(r)`$ in (2.16). The relations (2.24d) and (2.24e) are the orbifold duality transformations
$$G\begin{array}{c}\hfill \\ ^\sigma \hfill \end{array}𝒢(G;\sigma ),f\begin{array}{c}\hfill \\ ^\sigma \hfill \end{array}(f;\sigma )$$
(2.25)
from the untwisted metric $`G_{ab}`$ and structure constants $`f_{ab}^c`$ to the twisted metric $`𝒢(\sigma )`$ and the twisted structure constants $`(\sigma )`$ of sector $`\sigma `$.
The twisted tensors $`𝒢(\sigma )`$ and $`(\sigma )`$ inherit the spectral index periodicity
$$𝒢_{n(r)\pm \rho (\sigma ),\mu ;n(s)\nu }(\sigma )=𝒢_{n(r)\mu ;n(s)\nu }(\sigma ),_{n(r)\pm \rho (\sigma ),\mu ;n(s)\nu }^{n(t)\delta }(\sigma )=_{n(r)\mu ;n(s)\nu }^{n(t)\delta }(\sigma )$$
(2.26)
from the periodicity (2.12b) of the eigenvectors $`U(\sigma )`$ and $`U^{}(\sigma )`$, and similarly for $`n(s)`$ and $`n(t)`$. For the same reason, the periodicity $`n(r)n(r)\pm \rho (\sigma )`$ holds for each spectral index of each of the twisted tensors introduced below.
For the untwisted sector $`\sigma =0`$ one has $`U^{}(0)=1`$, $`\overline{n}(r)=0`$, $`\mu =a`$ and
$$𝒢_{ab}(0)=G_{ab},_{ab}^c(0)=f_{ab}^c,\widehat{J}_a(\sigma =0)=𝒥_a(\sigma =0)=J_a$$
(2.27)
so the system (2.4) reduces in this case to the OPE’s (2.1a) of the affine Lie algebra of the $`H`$-covariant algebra $`g`$. The number of twisted currents in (2.24f) is independent of $`\sigma `$ because $`U^{}(\sigma )`$ is an invertible square matrix. The twisted current system (2.4) was called $`𝔤(HAut(g);\sigma )`$ in Ref. References.
The twisted metric of sector $`\sigma `$ also satisfies
$$𝒢_{n(r)\mu ;n(s)\nu }(\sigma )=𝒢_{n(s)\nu ;n(r)\mu }(\sigma )$$
(2.28a)
$$𝒢_{n(r)\mu ;n(s)\nu }(\sigma )(1E_{n(r)}(\sigma )E_{n(s)}(\sigma ))=0$$
(2.28b)
$$𝒢_{n(r)\mu ;n(s)\nu }(\sigma )=\delta _{n(r)+n(s),0\text{ mod }\rho (\sigma )}𝒢_{n(r)\mu ;n(r),\nu }(\sigma )$$
(2.28c)
where the selection rule for $`𝒢`$ in (2.28b) is the dual<sup>e</sup><sup>e</sup>eTo prove the selection rule (2.28b), use (2.6c) to replace $`G`$ by $`\omega ^2G`$ in the orbifold duality transformation (2.24d), followed by the appropriate form of the $`H`$-eigenvalue problem in (2.10). All the selection rules below follow similarly from the orbifold duality transformations, the $`H`$-eigenvalue problem and the corresponding $`H`$-invariances of the untwisted tensors. in sector $`\sigma `$ of the $`H`$-invariance of $`G`$ in (2.6c), and (2.28c) is the solution<sup>f</sup><sup>f</sup>fWhen $`H`$ is a direct product group, further Kronecker factors can occur in the reduced twisted tensor $`𝒢_{n(r)\mu ;n(r),\nu }(\sigma )`$. Similarly, further Kronecker factors can occur in the reduced forms of each twisted tensor below. of the selection rule. Similarly, the twisted structure constants of sector $`\sigma `$ satisfy
$$_{n(r)\mu ;n(s)\nu }^{n(t)\delta }(\sigma )=_{n(s)\nu ;n(r)\mu }^{n(t)\delta }(\sigma )$$
(2.29a)
$$_{n(r)\mu ;n(s)\nu }^{n(u)\mu ^{}}(\sigma )_{n(t)\delta ;n(u)\mu ^{}}^{n(v)\gamma }(\sigma )+_{n(s)\nu ;n(t)\delta }^{n(u)\mu ^{}}(\sigma )_{n(r)\mu ;n(u)\mu ^{}}^{n(v)\gamma }(\sigma )$$
$$+_{n(t)\delta ;n(r)\mu }^{n(u)\mu ^{}}(\sigma )_{n(s)\nu ;n(u)\mu ^{}}^{n(v)\gamma }(\sigma )=0$$
(2.29b)
$$_{n(r)\mu ;n(s)\nu ;n(t)\delta }(\sigma )_{n(r)\mu ;n(s)\nu }^{n(u)ϵ}(\sigma )𝒢_{n(u)ϵ;n(t)\delta }(\sigma )=_{n(r)\mu ;n(t)\delta ;n(s)\nu }(\sigma )$$
(2.29c)
$$_{n(r)\mu ;n(s)\nu }^{n(t)\delta }(\sigma )(1E_{n(r)}(\sigma )E_{n(s)}(\sigma )E_{n(t)}(\sigma )^{})=0$$
(2.29d)
$$_{n(r)\mu ;n(s)\nu }^{n(t)\delta }(\sigma )=\delta _{n(r)+n(s)n(t),0\text{ mod }\rho (\sigma )}_{n(r)\mu ,n(s)\nu }^{n(r)+n(s),\delta }(\sigma )$$
(2.29e)
$$_{n(r)\mu ,n(s)\nu }^{n(r)+n(s),\delta }(\sigma )=0\text{unless }n(r)+n(s)\{n(r)\}.$$
(2.29f)
The relation (2.29b) is the Jacobi identity of the twisted structure constants, and the twisted structure constants in (2.29c) with all indices down are totally antisymmetric. The $``$-selection rule (2.29d) (and its solution in (2.29e,2.29e)) is the dual in sector $`\sigma `$ of the $`H`$-invariance (2.6c) of the untwisted structure constants.
The $`H`$-invariance conditions (2.6c) also imply<sup>g</sup><sup>g</sup>gTo prove that $`𝒢`$ is a class function, use the duality transformation (2.24d), the change of $`U`$ in (2.30b) and the $`H`$-invariance of $`G`$ in (2.6c). With the relevant duality transformation and $`H`$-invariance, each twisted tensor introduced below is similarly proved to be a class function. that the twisted metric and twisted structure constants are class functions$`^{\text{References}}`$ under conjugation in $`H`$
$$\omega (h_\sigma )v^{}(\sigma )\omega (h_\sigma )v(\sigma ),v(\sigma )v^{}(\sigma )=1,v(\sigma )H$$
(2.30a)
$$U(\sigma )U(\sigma )v(\sigma ),U^{}(\sigma )v^{}(\sigma )U^{}(\sigma ),$$
(2.30b)
$$𝒢(U(\sigma )v(\sigma );\sigma )=𝒢(U(\sigma );\sigma ),(U(\sigma )v(\sigma );\sigma )=(U(\sigma );\sigma )$$
(2.30c)
where the $`H`$-eigenvalue problem (2.10) was used to obtain the change of $`U(\sigma )`$, $`U^{}(\sigma )`$ under conjugation.
A more informative presentation of the twisted current system (2.4) is
$$\widehat{J}_{n(r)\mu }(z)\widehat{J}_{n(s)\nu }(w)=\frac{\delta _{n(r)+n(s),0\text{ mod }\rho (\sigma )}𝒢_{n(r)\mu ;n(r),\nu }(\sigma )}{(zw)^2}$$
$$+\frac{i_{n(r)\mu ;n(s)\nu }^{n(r)+n(s),\delta }(\sigma )\widehat{J}_{n(r)+n(s),\delta }(w)}{zw}+O((zw)^0)$$
(2.31a)
$$\widehat{J}_{n(r)\mu }(ze^{2\pi i})=e^{\frac{2\pi in(r)}{\rho (\sigma )}}\widehat{J}_{n(r)\mu }(z),\sigma =0,\mathrm{},N_c1$$
(2.31b)
where the $``$ term is summed only on the repeated degeneracy index $`\delta `$. In this form, we have incorporated the solutions of the $`𝒢`$\- and $``$-selection rules to exhibit the grading of the orbifold system. Looking back, one sees that this grading is a consequence of the $`H`$-covariance of the untwisted affine algebra $`g`$.
### 2.5 The general orbifold affine-Virasoro construction
The stress tensor $`\widehat{T}_\sigma `$ of sector $`\sigma `$ of $`A(H)/H`$ follows by first rewriting the stress tensor $`T`$ of the untwisted sector in terms of the eigencurrents and then using the OPE isomorphism (2.23a) to obtain<sup>h</sup><sup>h</sup>hSchematically, $`T=L:JJ:=:𝒥𝒥:\begin{array}{c}\hfill \\ ^\sigma \hfill \end{array}\widehat{T}_\sigma =:\widehat{J}\widehat{J}:`$. the derived OPE isomorphism$`^{\text{References}}`$
$$T(z)\begin{array}{c}\hfill \\ ^\sigma \hfill \end{array}\widehat{T}_\sigma (z).$$
(2.32)
The result is the general orbifold affine-Virasoro construction$`^{\text{References}}`$
$$\widehat{T}_\sigma (z)\widehat{T}_\sigma (w)=\frac{\widehat{c}(\sigma )/2}{(zw)^4}+(\frac{2}{(zw)^2}+\frac{_w}{zw})\widehat{T}_\sigma (w)+O((zw)^0)$$
(2.33a)
$$\widehat{T}_\sigma (z)=^{n(r)\mu ;n(s)\nu }(\sigma ):\widehat{J}_{n(r)\mu }(z)\widehat{J}_{n(s)\nu }(z):,\widehat{c}(\sigma )=c=2G_{ab}L^{ab}$$
(2.33b)
$$^{n(r)\mu ;n(s)\nu }(\sigma )=\chi \left(\sigma \right)_{n\left(r\right)\mu }^1\chi \left(\sigma \right)_{n\left(s\right)\nu }^1L^{ab}U^{}(\sigma )_a^{n(r)\mu }U^{}(\sigma )_b^{n(s)\nu }$$
(2.33c)
$$^{n(r)\mu ;n(s)\nu }(\sigma )=^{n(s)\nu ;n(r)\mu }(\sigma )$$
(2.33d)
which reduces in the untwisted sector of each orbifold to the general $`H`$-invariant affine-Virasoro construction $`\widehat{T}_{\sigma =0}=T`$ in (2.7a). Here (2.33c) is the orbifold duality transformation
$$L\begin{array}{c}\hfill \\ ^\sigma \hfill \end{array}(L;\sigma )$$
(2.34)
from the $`H`$-invariant inverse inertia tensor $`(0)=L`$ of the untwisted sector to the twisted inverse inertia tensor $`(\sigma )`$ of sector $`\sigma `$.
The symbol $`:():`$ in (2.33b) denotes the orbifold extension of OPE normal ordering$`^{\text{References},\text{References},\text{References}}`$
$$:\widehat{J}_{n(r)\mu }(z)\widehat{J}_{n(s)\nu }(w):\widehat{J}_{n(r)\mu }(z)\widehat{J}_{n(s)\nu }(w)\frac{𝒢_{n(r)\mu ;n(s)\nu }(\sigma )}{(zw)^2}\frac{i_{n(r)\mu ;n(s)\nu }^{n(t)\delta }(\sigma )\widehat{J}_{n(t)\delta }(w)}{zw}$$
(2.35a)
$$:\widehat{J}_{n(r)\mu }(w)\widehat{J}_{n(s)\nu }(w):=_w\frac{dz}{2\pi i}\frac{\widehat{J}_{n(r)\mu }(z)\widehat{J}_{n(s)\nu }(w)}{zw}$$
(2.35b)
where the contour in (2.35b) does not encircle the origin. We emphasize that the simplicity of the result (2.33b) would be compromised should other ordering prescriptions be employed (see Eqs. (3.11) and (H.2b)). Following the usage in Refs. References, References, and References, we will refer to the prescription (2.5) as OPE normal ordering.<sup>i</sup><sup>i</sup>iThe OPE normal ordering in (2.5) is the natural finite part of the orbifold operator product, but, except in the untwisted sectors,$`^{\text{References}}`$ this ordering is generally not a true normal ordering (with zero vev). Indeed, this fact accounts for the non-zero ground state conformal weights of most orbifold sectors.
The general orbifold affine-Virasoro construction (2.32) is the center of the orbifold program. Although the OPE isomorphisms (2.4) and (2.32) gave a simple derivation of this result, the Virasoro property (2.33a) of $`\widehat{T}_\sigma `$ was checked against the earlier general Virasoro construction$`^{\text{References}}`$ for cyclic permutation orbifolds. The Virasoro property (2.33a) was also checked by direct OPE computation in Ref. References (see also Sec. 5 and App. Appendix D.). Because of the OPE normal ordering, one finds that every step of this computation follows by the duality algorithm$`^{\text{References}}`$
$$a\begin{array}{c}\hfill \\ ^\sigma \hfill \end{array}n(r)\mu ,G\begin{array}{c}\hfill \\ ^\sigma \hfill \end{array}𝒢,f\begin{array}{c}\hfill \\ ^\sigma \hfill \end{array},L\begin{array}{c}\hfill \\ ^\sigma \hfill \end{array}$$
(2.36)
from the corresponding step in the direct OPE verification of the Virasoro property of the general affine-Virasoro construction $`T`$ in (2.3b). The final result of this computation is that the twisted inverse inertia tensor $`(\sigma )`$ must satisfy the general orbifold Virasoro master equation$`^{\text{References}}`$
$$^{n(r)\mu ;n(s)\nu }(\sigma )=2^{n(r)\mu ;n(r^{})\mu ^{}}(\sigma )𝒢_{n(r^{})\mu ^{};n(s^{})\nu ^{}}(\sigma )^{n(s^{})\nu ^{};n(s)\nu }(\sigma )$$
(2.37)
$$^{n(r^{})\mu ^{};n(s^{})\nu ^{}}(\sigma )^{n(t)\delta ;n(t^{})\delta ^{}}(\sigma )_{n(r^{})\mu ^{};n(t)\delta }^{n(r)\mu }(\sigma )_{n(s^{})\nu ^{};n(t^{})\delta ^{}}^{n(s)\nu }(\sigma )$$
$$^{n(r^{})\mu ^{};n(s^{})\nu ^{}}(\sigma )_{n(r^{})\mu ^{};n(t)\delta }^{n(t^{})\delta ^{}}(\sigma )_{n(s^{})\nu ^{};n(t^{})\delta ^{}}^{(n(r)\mu }(\sigma )^{n(s)\nu );n(t)\delta }(\sigma )$$
which can indeed be obtained by the duality algorithm (2.36) from the Virasoro master equation in (2.3c). Moreover, the general orbifold Virasoro master equation (2.37) is dual (by the orbifold duality transformation (2.33c)) to the Virasoro master equation. A final connection is that (2.37) reduces to the Virasoro master equation (2.3c) when the automorphism group $`H`$ is trivial. The result (2.37) generalizes the earlier orbifold Virasoro master equation$`^{\text{References}}`$ for cyclic permutation orbifolds.
The twisted inverse inertia tensors $`(\sigma )`$ of the various sectors of each orbifold are related by the orbifold duality transformation (2.33c) and its inverse $`L()`$ according to
$$^{n(r)^{}\mu ^{};n(s)^{}\nu ^{}}(\sigma ^{})$$
$$=\frac{\chi \left(\sigma \right)_{n\left(r\right)\mu }}{\chi \left(\sigma ^{}\right)_{n\left(r\right)^{}\mu ^{}}}\frac{\chi \left(\sigma \right)_{n\left(s\right)\nu }}{\chi \left(\sigma ^{}\right)_{n\left(r\right)^{}\mu ^{}}}^{n(r)\mu ;n(s)\nu }(\sigma )(U(\sigma )U^{}(\sigma ^{}))_{n(r)\mu }^{n(r)^{}\mu ^{}}(U(\sigma )U^{}(\sigma ^{}))_{n(s)\nu }^{n(s)^{}\nu ^{}}.$$
(2.38)
Similar relations are easily found among the sectors $`\sigma `$ of the twisted tensors $`𝒢(\sigma )`$, $`(\sigma )`$ introduced above, as well as among the sectors of the other twisted tensors introduced below.
The twisted inverse inertia tensors $`(\sigma )`$ also satisfy$`^{\text{References}}`$
$$^{n(r)\mu ;n(s)\nu }(\sigma )(1E_{n(r)}(\sigma )E_{n(s)}(\sigma ))=0$$
(2.39a)
$$^{n(r)\mu ;n(s)\nu }(\sigma )=\delta _{n(r)+n(s),0\text{ mod }\rho (\sigma )}^{n(r)\mu ;n(r),\nu }(\sigma )$$
(2.39b)
where the $``$-selection rule in (2.39a) is the dual in sector $`\sigma `$ of the $`H`$-invariance (2.7a) of the inverse inertia tensor $`L`$ in the untwisted sector. Another consequence of the $`H`$-invariance (2.7a) is that each $`(\sigma )`$ is a class function$`^{\text{References}}`$
$$(U(\sigma )v(\sigma );\sigma )=(U(\sigma );\sigma )$$
(2.40)
under the $`H`$-conjugation in (g).
Incorporation of the $``$-selection rule (2.5) gives a more informative presentation of the general orbifold affine-Virasoro construction
$$\widehat{T}_\sigma (z)=^{n(r)\mu ;n(r),\nu }(\sigma ):\widehat{J}_{n(r)\mu }(z)\widehat{J}_{n(r),\nu }(z):,\widehat{T}_\sigma (ze^{2\pi i})=\widehat{T}_\sigma (z)$$
(2.41a)
$$\widehat{c}(\sigma )=2𝒢_{n(r)\mu ;n(r),\nu }(\sigma )^{n(r)\mu ;n(r),\nu }(\sigma )=2G_{ab}L^{ab}=c,\sigma =0,\mathrm{},N_c1$$
(2.41b)
which shows that the orbifold stress tensors have trivial monodromy. Moreover, the $``$-selection rule (2.5) is consistent with the general orbifold Virasoro master equation (2.37), and the two can be combined to obtain a reduced master equation for $`^{n(r)\mu ;n(r),\nu }(\sigma )`$. Looking back, one sees that this consistency is a consequence of the underlying $`H`$-symmetry of the CFT $`A(H)`$.
### 2.6 The cyclic permutation orbifolds $`A(_\lambda )/_\lambda `$
As an example, we recall the seminal case$`^{\text{References}}`$ of the general cyclic permutation orbifold $`A(_\lambda )/_\lambda `$,
$$aaI,L^{ab}L^{aI,bJ}=L_{IJ}^{ab},I=0,\mathrm{},\lambda 1$$
(2.42a)
$$n(r),\mu r,aj,\widehat{J}_{n(r)\mu }\widehat{J}_{aj}^{(r)},\chi \left(\sigma \right)_{n\left(r\right)\mu }\chi \left(\sigma \right)_{raj}=\sqrt{\rho (\sigma )}$$
(2.42b)
$$𝒢_{raj;sbl}(\sigma )=\rho (\sigma )k\eta _{ab}\delta _{jl}\delta _{r+s,0\text{ mod }\rho (\sigma )},_{raj;sbl}^{tcm}(\sigma )=f_{ab}^c\delta _{jl}\delta _l^m\delta _{r+st,0\text{ mod }\rho (\sigma )}$$
(2.42c)
$$\widehat{J}_{aj}^{(r)}(z)\widehat{J}_{bl}^{(s)}(w)=\delta _{jl}\{\frac{\rho (\sigma )k\eta _{ab}\delta _{r+s,0\text{ mod }\rho (\sigma )}}{(zw)^2}+\frac{if_{ab}^c\widehat{J}_{cj}^{(r+s)}(w)}{(zw)}\}+O((zw)^0)$$
(2.42d)
$$\widehat{J}_{aj}^{(r)}(ze^{2\pi i})=e^{\frac{2\pi ir}{\rho (\sigma )}}\widehat{J}_{aj}^{(r)}(z),\widehat{J}_{aj}^{(r\pm \rho (\sigma ))}(z)=\widehat{J}_{aj}^{(r)}(z)$$
(2.42e)
$$\widehat{T}_\sigma (z)=\underset{r=0}{\overset{\rho (\sigma )1}{}}\underset{j,l=0}{\overset{\frac{\lambda }{\rho (\sigma )}1}{}}^{raj;r,bl}(\sigma ):\widehat{J}_{aj}^{(r)}(z)\widehat{J}_{bl}^{(r)}(z):$$
(2.42f)
$$^{raj;r,bl}(\sigma )=\frac{1}{\rho (\sigma )}\underset{s=0}{\overset{\rho (\sigma )1}{}}e^{\frac{2\pi iN(\sigma )rs}{\rho (\sigma )}}L_{\frac{\lambda }{\rho (\sigma )}s+jl}^{ab}$$
(2.42g)
$$\widehat{\mathrm{\Delta }}_0(\sigma )=\frac{\lambda k\eta _{ab}}{4\rho ^2(\sigma )}\{\frac{\rho ^2(\sigma )1}{3}L_0^{ab}\underset{r=1}{\overset{\rho (\sigma )1}{}}csc^2(\frac{\pi N(\sigma )r}{\rho (\sigma )})L_{\frac{\lambda }{\rho (\sigma )}r}^{ab}\}$$
(2.42h)
$$a,b=1,\mathrm{},\text{dim}𝔤,\overline{r},\overline{s}=0,\mathrm{},\rho (\sigma )1,j,l=0,\mathrm{},\frac{\lambda }{\rho (\sigma )}1,\sigma =0,\mathrm{},\lambda 1$$
(2.42i)
whose relation to the orbifold Virasoro master equation$`^{\text{References}}`$ provided a nontrivial check of the orbifold program.
The special case (2.6) follows from the discussion above and the $`_\lambda `$-eigendata in (2.3). In this case the twisted current system (2.42d,2.42d) describes a sector-dependent set of semisimple orbifold affine algebras.$`^{\text{References}}`$ The orbifold duality transformation $`(L;\sigma )`$ in (2.42g) is a set of discrete Fourier transforms of $`L`$, with spectral index periodicity $`rr+\rho (\sigma )`$, where $`L`$ is the inverse inertia tensor of any $`_\lambda `$(permutation)-invariant CFT $`A(_\lambda )`$. The twisted inverse inertia tensors $`^{raj;r,bl}(\sigma )`$ were called $`_r^{a(j)b(l)}(\sigma )`$ in Ref. References and the result in (2.42h) is the ground state conformal weight of sector $`\sigma `$.
The orbifold duality transformations for $`𝒢`$ and $``$ in (2.42c) can also be written as discrete Fourier transforms and, indeed, because of the periodicity of the eigenvectors of each $`H`$-eigenvalue problem, all the orbifold duality transformations of this paper
$$𝒢(G;\sigma ),(f;\sigma ),(L;\sigma ),\mathrm{}$$
(2.43)
are discrete Fourier transforms with period $`\rho (\sigma )`$ in the spectral indices $`\{n(r)\}`$.
## 3 The Mode Formulation of Orbifold Theory
As an application of the local formulation above, we turn now to the mode formulation of the general current-algebraic orbifold.
### 3.1 The general twisted current algebra
The first step in the mode formulation is to find the twisted current algebra of each sector of each orbifold $`A(H)/H`$, beginning with the general twisted current system (2.4).
The mode expansion of the twisted currents
$$\widehat{J}_{n(r)\mu }(z)=\underset{m}{}\widehat{J}_{n(r)\mu }(m+\frac{n\left(r\right)}{\rho \left(\sigma \right)})z^{(m+\frac{n(r)}{\rho (\sigma )})1}$$
(3.1a)
$$\widehat{J}_{n(r)\mu }(z)=\widehat{J}_{n(r)\pm \rho (\sigma ),\mu }(z)=\underset{m}{}\widehat{J}_{n(r)\pm \rho (\sigma ),\mu }(m+\frac{n\left(r\right)\pm \rho \left(\sigma \right)}{\rho \left(\sigma \right)})z^{(m+\frac{n(r)\pm \rho (\sigma )}{\rho (\sigma )})1}$$
(3.1b)
follows from the monodromies (2.31b) and the periodicity (2.24c). From (3.1b) we find the periodicity of the modes
$$\widehat{J}_{n(r)_\pm \rho (\sigma ),\mu }(m1+\frac{n\left(r\right)\pm \rho \left(\sigma \right)}{\rho \left(\sigma \right)})=\widehat{J}_{n(r)_\pm \rho (\sigma ),\mu }(m+\frac{n\left(r\right)}{\rho \left(\sigma \right)})=\widehat{J}_{n(r)\mu }(m+\frac{n\left(r\right)}{\rho \left(\sigma \right)})$$
(3.2a)
$$\widehat{J}_{n(r),\mu }(m\frac{n\left(r\right)}{\rho \left(\sigma \right)})=\widehat{J}_{\rho (\sigma )n(r),\mu }(m1+\frac{\rho \left(\sigma \right)n\left(r\right)}{\rho \left(\sigma \right)})$$
(3.2b)
$$\widehat{J}_{n(r)_\pm \rho (\sigma ),\mu }(m+\frac{n\left(r\right)\pm \rho \left(\sigma \right)}{\rho \left(\sigma \right)})\widehat{J}_{n(r)\mu }(m+\frac{n\left(r\right)}{\rho \left(\sigma \right)})$$
(3.2c)
which includes and generalizes the mode periodicity of orbifold affine algebra.$`^{\text{References},\text{References},\text{References}}`$ Note that the modes do not satisfy the naive periodicity (3.2c).
We introduce next the operators $`\widehat{\widehat{J}}`$ with trivial monodromy
$$\widehat{\widehat{J}}_{n(r)\mu }(z)z^{\frac{n(r)}{\rho (\sigma )}}\widehat{J}_{n(r)\mu }(z)$$
(3.3a)
$$\widehat{\widehat{J}}_{n(r)\mu }(z)\widehat{\widehat{J}}_{n(s)\nu }(w)=\delta _{n(r)+n(s),0\text{ mod }\rho (\sigma )}\omega ^{\frac{n(r)+n(s)}{\rho (\sigma )}}(\frac{1}{(zw)^2}+\frac{n(r)/\rho (\sigma )}{w(zw)})𝒢_{n(r)\mu ;n(r),\nu }(\sigma )$$
$$+\frac{i_{n(r)\mu ;n(s)\nu }^{n(r)+n(s),\delta }(\sigma )\widehat{\widehat{J}}_{n(r)+n(s),\delta }(w)}{zw}+O((zw)^0).$$
(3.3b)
The Kronecker delta in the first term of (3.3b) guarantees that these OPE’s are free of branch cuts, which is a necessary condition for the monodromies (2.31b) to be consistent with the OPE’s (2.31a).
Then standard analysis of the OPE’s (3.3b) gives the general twisted current algebra $`\widehat{𝔤}(\sigma )\widehat{𝔤}(HAut(g);\sigma )`$
$$[\widehat{J}_{n(r)\mu }(m+\frac{n\left(r\right)}{\rho \left(\sigma \right)}),\widehat{J}_{n(s)\nu }(n+\frac{n\left(s\right)}{\rho \left(\sigma \right)})]=i_{n(r)\mu ;n(s)\nu }^{n(r)+n(s),\delta }(\sigma )\widehat{J}_{n(r)+n(s),\delta }(m+n+\frac{n\left(r\right)+n\left(s\right)}{\rho \left(\sigma \right)})$$
$$+(m+\frac{n\left(r\right)}{\rho \left(\sigma \right)})\delta _{m+n+\frac{n(r)+n(s)}{\rho (\sigma )},0}𝒢_{n(r)\mu ;n(r),\nu }(\sigma )$$
(3.4a)
$$m,n,\sigma =0,\mathrm{},N_c1$$
(3.4b)
in sector $`\sigma `$ of each orbifold $`A(H)/H`$. Here $`N_c`$ is the number of conjugacy classes of $`H`$, and the duality transformations for $`𝒢`$ and $``$ are given in Eq. (2.4). To obtain (3.1) we have used the identity
$$\delta _{n(r)+n(s),0\text{ mod }\rho (\sigma )}\delta _{m+n+\frac{n(r)+n(s)}{\rho (\sigma )},0}=\delta _{m+n+\frac{n(r)+n(s)}{\rho (\sigma )},0}$$
(3.5)
to simplify the $`𝒢`$ term.
The general twisted current algebra (3.1) shows a natural grading across all orbifolds. The general twisted current algebra also satisfies the Jacobi identity (see App. 9.5), which is the ultimate check of the consistency of the principle (2.4).
Looking back, one sees that the grading and the consistency of the general twisted current algebra are consequences of the $`H`$-covariance of the underlying affine algebra (2.1b) on $`g`$. The general twisted current algebra is also consistent with the mode periodicity relations (3.1b) and, in the untwisted sector $`\sigma =0`$, the general twisted current algebra reduces to the $`H`$-covariant affine algebra on $`g`$.
Another feature of the general twisted current algebra is the integral affine subalgebra
$$\widehat{𝔤}^{(0)}(\sigma )\widehat{𝔤}(\sigma ):[\widehat{J}_{0\mu }(m),\widehat{J}_{0\nu }(n)]=i_{0\mu ;0\nu }^{0\delta }(\sigma )\widehat{J}_{0\delta }(m+n)+m\delta _{m+n,0}𝒢_{0\mu ;0\nu }(\sigma )$$
(3.6)
which is an untwisted affine Lie algebra generated by the set of currents in twist class $`\overline{n}(r)=0`$. The integral affine subalgebra $`\widehat{𝔤}^{(0)}(\sigma )`$ is non-trivial for all known non-abelian $`\widehat{𝔤}(\sigma )`$.
As an example of $`\widehat{𝔤}(\sigma )`$, we mention the semisimple orbifold affine algebra at orbifold affine level $`\widehat{k}(\sigma )=\rho (\sigma )k`$
$$\widehat{𝔤}(\sigma )=\widehat{𝔤}(_\lambda \text{(permutation)}Aut(g);\sigma ),g=_{I=0}^{\lambda 1}𝔤^I,𝔤^I𝔤$$
(3.7a)
$$[\widehat{J}_{aj}^{(r)}(m+\frac{r}{\rho \left(\sigma \right)}),\widehat{J}_{bl}^{(s)}(n+\frac{s}{\rho \left(\sigma \right)})]=\delta _{jl}\{if_{ab}^c\widehat{J}_{cj}^{(r+s)}(m+n+\frac{r+s}{\rho \left(\sigma \right)})+(m+\frac{r}{\rho \left(\sigma \right)})\delta _{m+n+\frac{r+s}{\rho (\sigma )},0}\rho (\sigma )k\eta _{ab}\}$$
(3.7b)
$$\widehat{J}_{aj}^{(r\pm \rho (\sigma ))}(m+\frac{r}{\rho \left(\sigma \right)})=\widehat{J}_{aj}^{(r)}(m+\frac{r}{\rho \left(\sigma \right)})$$
(3.7c)
$$a,b=1,\mathrm{},\text{dim}𝔤,\overline{r},\overline{s}=0,\mathrm{},\rho (\sigma )1,j,l=0,\mathrm{},\frac{\lambda }{\rho (\sigma )}1,\sigma =0,\mathrm{},\lambda 1$$
(3.7d)
which holds<sup>j</sup><sup>j</sup>jFor prime $`\lambda `$ one has $`\rho (\sigma )=\lambda `$ and $`j=0`$ for all $`\lambda 1`$ twisted sectors, so each twisted sector has only the simple orbifold affine algebra at level $`\widehat{k}=\lambda k`$. This was first seen$`^{\text{References}}`$ in the characters of cyclic copy orbifolds. in sector $`\sigma `$ of each cyclic permutation orbifold $`A(_\lambda )/_\lambda `$. Using the data (2.42b,2.42b,2.42h), the form (3.1) follows from the general twisted current algebra (3.1). The algebra is also equivalent to the twisted current system in (2.6).
Here is a summary of what is known about (the $`H`$-eigenvalue problem and) the twisted current algebra $`\widehat{𝔤}(HAut(g);\sigma )`$ :
$``$ $`H`$ is any subgroup of $`S_N`$(permutation) $``$ sets of commuting orbifold affine
algebras$`^{\text{References}\text{References},\text{References},\text{References}\text{References}}`$ at various levels.
$``$ $`H`$ is a group of inner or outer automorphisms of simple $`g`$ $``$ inner-automorphically
twisted$`^{\text{References},\text{References},\text{References},\text{References},\text{References}}`$ and outer-automorphically twisted$`^{\text{References},\text{References},\text{References}}`$ affine Lie algebras.
$``$ $`H`$ is a group of inner automorphisms of each copy $`𝔤^I𝔤`$ times a permutation group
which acts among the copies $``$ doubly-twisted affine algebras.$`^{\text{References},\text{References},\text{References}}`$
The results$`^{\text{References}}`$ for $`H=_\lambda `$(permutation) are collected in Eqs. (2.3), (2.6) and (3.1). The case $`H=S_N`$(permutation) is discussed in detail in Sec. 8 and the case of the general group $`H=H(d)`$ of inner automorphisms of simple $`g`$ is discussed in Sec. 9. The setups for the case $`H=𝔻_\lambda `$(permutation) and the outer automorphism groups of simple $`g`$ are included in Apps. Appendix I. and Appendix J. respectively.
There are many other cases for which the twisted current algebras have not yet been worked out (e.g. $`H`$ is a group of outer automorphisms of each copy $`𝔤^I𝔤`$ times a permutation group which acts among the copies). Another basis for the general twisted current algebra $`\widehat{𝔤}(\sigma )`$ is given in App. Appendix B..
### 3.2 Operator products and mode ordering
In addition to OPE normal ordering (2.5), it is convenient to introduce a mode ordering<sup>k</sup><sup>k</sup>kMany mode orderings are possible, the most effective ordering being the one that is closest to “normal” (zero vev), given the mode characterization of the ground state of a given sector. The $`M`$-ordering in (3.8) is a normal ordering for the orbifold affine algebras of the general permutation orbifold (see Secs. 7 and 8).
$$:\widehat{J}_{n(r)\mu }(m+\frac{n\left(r\right)}{\rho \left(\sigma \right)})\widehat{J}_{n(s)\nu }(n+\frac{n\left(s\right)}{\rho \left(\sigma \right)}):_M\theta (m+\frac{n\left(r\right)}{\rho \left(\sigma \right)}0)\widehat{J}_{n(s)\nu }(n+\frac{n\left(s\right)}{\rho \left(\sigma \right)})\widehat{J}_{n(r)\mu }(m+\frac{n\left(r\right)}{\rho \left(\sigma \right)})$$
$$+\theta (m+\frac{n\left(r\right)}{\rho \left(\sigma \right)}<0)\widehat{J}_{n(r)\mu }(m+\frac{n\left(r\right)}{\rho \left(\sigma \right)})\widehat{J}_{n(s)\nu }(n+\frac{n\left(s\right)}{\rho \left(\sigma \right)})$$
(3.8)
denoted by the subscript $`M`$. Then one may use (3.1a), (3.1) and (3.8) to compute the exact operator product of the twisted currents
$$\widehat{J}_{n(r)\mu }(z)\widehat{J}_{n(s)\nu }(w)=(\frac{w}{z})^{\frac{\overline{n}(r)}{\rho (\sigma )}}\{[\frac{1}{(zw)^2}+\frac{\overline{n}(r)/\rho (\sigma )}{w(zw)}]𝒢_{n(r)\mu ;n(s)\nu }(\sigma )$$
$$+\frac{i_{n(r)\mu ;n(s)\nu }^{n(r)+n(s),\delta }(\sigma )\widehat{J}_{n(r)+n(s),\delta }(w)}{zw}\}+:\widehat{J}_{n(r)\mu }(z)\widehat{J}_{n(s)\nu }(w):_M$$
(3.9)
where $`\overline{n}(r)`$ is the twist class of $`\widehat{J}_{n(r)\mu }`$ (see Eq. (2.16)). The partial conversion $`n(r)\overline{n}(r)`$ seen here is detailed in App. Appendix C., and it is easily checked that the result (3.9) is consistent with the OPE’s (2.4). We also note that every term in (3.9) is periodic under $`n(r)n(r)+\rho (\sigma )`$, $`\overline{n}(r)\overline{n}(r)`$.
Comparing (3.9) with (2.5), we may express the OPE normal ordered product
$$:\widehat{J}_{n(r)\mu }(z)\widehat{J}_{n(s)\nu }(z):=:\widehat{J}_{n(r)\mu }(z)\widehat{J}_{n(s)\nu }(z):_M\frac{i\overline{n}(r)}{z\rho (\sigma )}_{n(r)\mu ;n(s)\nu }^{n(r)+n(s),\delta }(\sigma )\widehat{J}_{n(r)+n(s),\delta }(z)$$
$$+\frac{1}{z^2}\frac{\overline{n}(r)}{2\rho (\sigma )}(1\frac{\overline{n}(r)}{\rho (\sigma )})𝒢_{n(r)\mu ;n(s)\nu }(\sigma )$$
(3.10)
in terms of an $`M`$ ordered product.
With the relation (3.10), the general orbifold stress tensor (2.41a) can be expressed in terms of an M ordered product as
$`\widehat{T}_\sigma (z)`$ $`=`$ $`{\displaystyle \underset{r,\mu ,\nu }{}}^{n(r)\mu ;n(r),\nu }(\sigma )\{:\widehat{J}_{n(r)\mu }(z)\widehat{J}_{n(r),\nu }(z):_M{\displaystyle \frac{i\overline{n}(r)}{\rho (\sigma )}}_{n(r)\mu ;n(r),\nu }^{0\delta }(\sigma ){\displaystyle \frac{\widehat{J}_{0\delta }(z)}{z}}`$ (3.11)
$`+{\displaystyle \frac{1}{z^2}}{\displaystyle \frac{\overline{n}(r)}{2\rho (\sigma )}}(1{\displaystyle \frac{\overline{n}(r)}{\rho (\sigma )}})𝒢_{n(r)\mu ;n(r),\nu }(\sigma )\}.`$
In this form, we note the presence of the Freericks-Halpern$`^{\text{References}}`$ term $`\widehat{J}_{0\delta }(z)/z`$, where $`\widehat{J}_{0\delta }`$ are the generators of the integral affine subalgebra.
### 3.3 The Virasoro generators of sector $`\sigma `$
The OPE normal ordered product of the twisted current modes can also be expressed in terms of $`M`$ ordering
$$:\widehat{J}_{n(r)\mu }(m+\frac{n\left(r\right)}{\rho \left(\sigma \right)})\widehat{J}_{n(s)\nu }(n+\frac{n\left(s\right)}{\rho \left(\sigma \right)}):=:\widehat{J}_{n(r)\mu }(m+\frac{n\left(r\right)}{\rho \left(\sigma \right)})\widehat{J}_{n(s)\nu }(n+\frac{n\left(s\right)}{\rho \left(\sigma \right)}):_M$$
$$i\frac{\overline{n}(r)}{\rho (\sigma )}_{n(r)\mu ;n(s)\nu }^{n(r)+n(s),\delta }(\sigma )\widehat{J}_{n(r)+n(s),\delta }(m+n+\frac{n\left(r\right)+n\left(s\right)}{\rho \left(\sigma \right)})$$
$$+\frac{\overline{n}(r)}{2\rho (\sigma )}(1\frac{\overline{n}(r)}{\rho (\sigma )})\delta _{m+n+\frac{n(r)+n(s)}{\rho (\sigma )},0}𝒢_{n(r)\mu ;n(r),\nu }(\sigma ).$$
(3.12)
Then (using the mode expansion (3.1a) of the twisted currents) the general orbifold stress tensor in (2.5) or (3.11) gives the orbifold Virasoro generators
$$\widehat{T}_\sigma (z)=\underset{m}{}L_\sigma (m)z^{m2},\sigma =0,\mathrm{},N_c1$$
(3.13a)
$`L_\sigma (m)`$ $`=`$ $`{\displaystyle \underset{r,\mu ,\nu }{}}^{n(r)\mu ;n(r),\nu }(\sigma ){\displaystyle \underset{p}{}}:\widehat{J}_{n(r)\mu }(p+\frac{n\left(r\right)}{\rho \left(\sigma \right)})\widehat{J}_{n(r),\nu }(mp\frac{n\left(r\right)}{\rho \left(\sigma \right)}):`$ (3.13c)
$`=`$ $`{\displaystyle \underset{r,\mu ,\nu }{}}^{n(r)\mu ;n(r),\nu }(\sigma )\{{\displaystyle \underset{p}{}}:\widehat{J}_{n(r)\mu }(p+\frac{n\left(r\right)}{\rho \left(\sigma \right)})\widehat{J}_{n(r),\nu }(mp\frac{n\left(r\right)}{\rho \left(\sigma \right)}):_M`$
$`i{\displaystyle \frac{\overline{n}(r)}{\rho (\sigma )}}_{n(r)\mu ;n(r),\nu }^{0\delta }(\sigma )\widehat{J}_{0\delta }(m)`$
$`+\delta _{m,0}{\displaystyle \frac{\overline{n}(r)}{2\rho (\sigma )}}(1{\displaystyle \frac{\overline{n}(r)}{\rho (\sigma )}})𝒢_{n(r)\mu ;n(r),\nu }(\sigma )\}`$
for all sectors $`\sigma `$ of each orbifold $`A(H)/H`$. This result is the mode form of the general orbifold affine-Virasoro construction. For the untwisted sector $`\sigma =0`$, the general construction (3.12) reduces to the mode formulation of the general $`H`$-invariant affine-Virasoro construction$`^{\text{References}}`$ in (2.7a).
The simplest example of this result is the orbifold $`U(1)/_2`$:
$$J(z)J(w)=\frac{1}{(zw)^2}+O((zw)^0),T(z)=\frac{1}{2}:J^2(z):$$
(3.14a)
$$\sigma =1:J(z)^{}=J(w),\rho (1)=2,U(1)=\chi \left(1\right)=1$$
(3.14b)
$$\widehat{J}(z)\widehat{J}(w)=\frac{1}{(zw)^2}+O((zw)^0),\widehat{J}(ze^{2\pi i})=\widehat{J}(z)$$
(3.14c)
$$\widehat{J}(z)=\underset{m}{}\widehat{J}(m+\frac{1}{2})z^{(m+\frac{1}{2})1},[\widehat{J}(m+\frac{1}{2}),\widehat{J}(n+\frac{1}{2})]=(m+\frac{1}{2})\delta _{m+n+1,0}$$
(3.14d)
$$L_{\sigma =1}(m)=\frac{1}{2}\underset{p}{}:\widehat{J}(p+\frac{1}{2})\widehat{J}(mp\frac{1}{2}):_M+\delta _{m,0}\frac{1}{16}$$
(3.14e)
$$\widehat{J}((m+\frac{1}{2})0)|0=(L_{\sigma =1}(m0)\delta _{m,0}\frac{1}{16})|0=0.$$
(3.14f)
Here $`\widehat{J}\widehat{J}_{n(r)=1}=\widehat{J}_{n(r)=1}=i\widehat{X}`$ is the current of the standard $`\frac{1}{2}`$-integrally moded scalar field $`\widehat{X}`$. This twisted field was invented$`^{\text{References}}`$ by Halpern and Thorn, who used it to construct the first twisted sector of an orbifold.
For the seminal non-abelian example,$`^{\text{References}}`$ the Virasoro generators of the permutation orbifolds $`A(_\lambda )/_\lambda `$
$`L_\sigma (m)`$ $`=`$ $`{\displaystyle \underset{r=0}{\overset{\rho (\sigma )1}{}}}[{\displaystyle \underset{j,l=0}{\overset{\frac{\lambda }{\rho (\sigma )}1}{}}}^{raj;r,bl}(\sigma ){\displaystyle \underset{p}{}}:\widehat{J}_{aj}^{(r)}(p+\frac{r}{\rho \left(\sigma \right)})\widehat{J}_{bl}^{(r)}(mp\frac{r}{\rho \left(\sigma \right)}):_M`$ (3.15)
$`+{\displaystyle \underset{j=0}{\overset{\frac{\lambda }{\rho (\sigma )}1}{}}}^{raj;r,bj}(\sigma )\{{\displaystyle \frac{ir}{\rho (\sigma )}}f_{ab}^c\widehat{J}_{cj}^{(0)}(m)+\delta _{m,0}{\displaystyle \frac{r}{2\rho (\sigma )}}(1{\displaystyle \frac{r}{\rho (\sigma )}})k\rho (\sigma )\eta _{ab}\}]`$
follow from (2.6) and (3.12).
## 4 The Orbifold Adjoint Operation
### 4.1 The orbifold conjugation matrix
To study the adjoint operation in the twisted sectors, we first introduce the orbifold conjugation matrix $``$,
$$\rho \begin{array}{c}\hfill \\ ^\sigma \hfill \end{array}(\rho ;\sigma )$$
(4.1a)
$$_{n(r)\mu }^{n(s)\nu }(\sigma )\chi \left(\sigma \right)_{n\left(r\right)\mu }^{}\chi \left(\sigma \right)_{n\left(s\right)\nu }^1U(\sigma )_{n(r)\mu }^a\rho _a^bU^{}(\sigma )_b^{n(s)\nu }$$
(4.1b)
$$(\sigma )^{}(\sigma )=(\sigma )(\sigma )^{}=1$$
(4.1c)
$$_{n(r)\mu }^{n(s)\nu }(\sigma )(1E_{n(r)}(\sigma )E_{n(s)}(\sigma ))=0$$
(4.1d)
$$_{n(r)\mu }^{n(s)\nu }(\sigma )=\delta _{n(r)+n(s),0\text{ mod }\rho (\sigma )}_{n(r)\mu }^{n(r),\nu }(\sigma )$$
(4.1e)
which is the dual in sector $`\sigma `$ (see the following subsection) of the conjugation matrix $`\rho `$ in Eq. (b). Eq. (4.1a) and its realization in (4.1b) is another orbifold duality transformation on the same footing as the earlier ones for $`𝒢,`$ and $``$. From Eq. (4.1b) and the $`H`$-eigenvalue problem (2.10), one sees that the $``$-selection rule in (4.1d) (and its solution in (4.1e)) is the dual in sector $`\sigma `$ of the $`H`$-covariance (2.6d) of the conjugation matrix $`\rho `$. The orbifold conjugation matrix $``$ also controls the complex conjugation of the other twisted tensors<sup>l</sup><sup>l</sup>lTo prove (4.2c), for example, use the duality transformation $`(L)`$ in (2.33c) and the $`L^{}`$ relation in (2.5), followed by the inverse $`L()`$ of the duality transformation.
$$𝒢_{n(r)\mu ;n(s)\nu }(\sigma )^{}=_{n(r)\mu }^{n(t)\mu ^{}}(\sigma )_{n(s)\nu }^{n(u)\nu ^{}}(\sigma )𝒢_{n(t)\mu ^{};n(u)\nu ^{}}(\sigma )$$
(4.2a)
$$_{n(r)\mu ;n(s)\nu }^{n(t)\delta }(\sigma )^{}=_{n(r)\mu }^{n(u)\mu ^{}}(\sigma )_{n(s)\nu }^{n(v)\nu ^{}}(\sigma )_{n(u)\mu ^{};n(v)\nu ^{}}^{n(w)\delta ^{}}(\sigma )_{n(w)\delta ^{}}^{n(t)\delta }(\sigma )^{}$$
(4.2b)
$$^{n(r)\mu ;n(s)\nu }(\sigma )^{}=^{n(t)\mu ^{};n(u)\nu ^{}}(\sigma )_{n(t)\mu ^{}}^{n(r)\mu }(\sigma )^{}_{n(u)\nu ^{}}^{n(s)\nu }(\sigma )^{}$$
(4.2c)
in each sector $`\sigma `$ of each orbifold $`A(H)/H`$. These results are dual to Eqs. (2.2b) and (2.5), so that, in particular, the relation (4.2c) holds when the CFT $`A(H)`$ is unitary.
Another consequence of the $`H`$-covariance of $`\rho `$ in (2.6d) is that each $`(\sigma )`$ is a class function
$$(U(\sigma )v(\sigma );\sigma )=(U(\sigma );\sigma )$$
(4.3)
under the $`H`$-conjugation in (g).
### 4.2 The adjoint of the twisted currents
We can now give the adjoint of the twisted currents of sector $`\sigma `$
$$\widehat{J}_{n(r)\mu }(m+\frac{n\left(r\right)}{\rho \left(\sigma \right)})^{}=\underset{\nu }{}_{n(r)\mu }^{n(r),\nu }(\sigma )\widehat{J}_{n(r),\nu }(m\frac{n\left(r\right)}{\rho \left(\sigma \right)})$$
(4.4a)
$$\widehat{J}_{n(r)\mu }(m+\frac{n\left(r\right)}{\rho \left(\sigma \right)})^{}=\widehat{J}_{n(r)\mu }(m+\frac{n\left(r\right)}{\rho \left(\sigma \right)})$$
(4.4b)
where $`(\sigma )`$ is the orbifold conjugation matrix defined above. The orbifold adjoint operation (4.4a) is the dual<sup>m</sup><sup>m</sup>mThe local form $`\widehat{J}(z)^{}=\widehat{J}(z^1)z^2`$ of the adjoint operation (4.4a), including the form (4.1b) of $``$, follows by duality transformation from the local form $`J(z)^{}=\rho J(z^1)z^2`$ of the adjoint operation (b) in the untwisted sector. in sector $`\sigma `$ of the adjoint operation (2.2a) on the untwisted currents. The consistency relation (4.4b) is easily checked with Eqs. (4.4a) and (A.3).
The orbifold adjoint operation (4.4a) defines a real form of the general twisted current algebra $`\widehat{𝔤}(\sigma )`$ in (3.1) because the adjoint of the algebra
$$[\widehat{J}_{n(r)\mu }(m+\frac{n\left(r\right)}{\rho \left(\sigma \right)})^{},\widehat{J}_{n(s)\nu }(n+\frac{n\left(s\right)}{\rho \left(\sigma \right)})^{}]=(m+\frac{n\left(r\right)}{\rho \left(\sigma \right)})𝒢_{n(r)\mu ;n(r),\nu }(\sigma )^{}\delta _{m+n+\frac{n(r)+n(s)}{\rho (\sigma )},0}$$
$$+i_{n(r)\mu ;n(s)\nu }^{n(r)+n(s),\delta }(\sigma )^{}\widehat{J}_{n(r)+n(s),\delta }(m+n+\frac{n\left(r\right)+n\left(s\right)}{\rho \left(\sigma \right)})^{}$$
(4.5)
is consistent with $`\widehat{𝔤}(\sigma )`$. To verify this statement, use (A.4a) and (A.4b).
In the case$`^{\text{References}}`$ of the permutation orbifolds $`A(_\lambda )/_\lambda `$, the twisted current algebra is the semisimple orbifold affine algebra (3.1) and we obtain the orbifold adjoint operation
$$\rho _a^b\rho _{aI}^{bJ}=\rho _a^b\delta _I^J,I,J=0,\mathrm{},\lambda 1$$
(4.6a)
$$(\sigma )_{raj}^{sbl}(\sigma )=\rho _a^b\delta _j^l\delta _{r+s,0\text{ mod }\rho (\sigma )},\widehat{J}_{aj}^{(r)}(m+\frac{r}{\rho \left(\sigma \right)})^{}=\rho _a^b\widehat{J}_{bj}^{(r)}(m\frac{r}{\rho \left(\sigma \right)})$$
(4.6b)
from (2.3), (2.6) and (4.4a). This adjoint operation is in agreement with the standard adjoint operation for orbifold affine algebra$`^{\text{References},\text{References}}`$ and it is known$`^{\text{References}}`$ that this adjoint operation guarantees unitarity of the twisted affine Hilbert space when the untwisted affine Hilbert space is unitary.
More generally, as we will note in Secs. 7, 8 and 9, the orbifold adjoint operation defines a unitary twisted affine Hilbert space for each sector of the orbifolds<sup>n</sup><sup>n</sup>nIn all these cases, unitarity of the twisted affine Hilbert spaces follows easily from known orbifold induction procedures.$`^{\text{References},\text{References},\text{References},\text{References},\text{References},\text{References}}`$
$``$ $`A(H)/H`$, where $`H`$ is any subgroup of $`S_N`$(permutation)
$``$ $`A(H(d))/H(d)`$, where $`H(d)`$ is any group of inner automorphisms of simple $`g`$
$``$ $`A(H)/H`$, where $`H`$ is a product of any subgroup of $`S_N`$(permutation) times a group of
inner automorphisms
when the untwisted affine Hilbert space is unitary. On the basis of this evidence, we conjecture that the orbifold adjoint operation implies unitarity of all the twisted affine Hilbert spaces of $`A(H)/H`$, given the unitarity of the untwisted affine Hilbert space.
### 4.3 The adjoint of the orbifold Virasoro generators
The adjoint (4.4a) of the twisted currents gives the desired adjoint operation on the Virasoro generators
$$L_\sigma (m)^{}=L_\sigma (m)$$
(4.7)
in every sector $`\sigma `$ of each orbifold $`A(H)/H`$. To see this, apply the orbifold adjoint operation (4.4a) to the mode form (3.13c) of the Virasoro generators and use the $`^{}`$ relation (A.4c). In further detail, one finds that the $`(\widehat{J})^2`$, $`(\widehat{J})^1`$ and $`(\widehat{J})^0`$ terms in (3.13c) satisfy the adjoint relation (4.7) separately.
It follows that all the twisted sectors of the orbifolds bulleted in Subsec. 4.2 are unitary when the CFT’s $`A(H)`$ are unitary,<sup>o</sup><sup>o</sup>oThis result was established for the permutation orbifolds $`A(_\lambda )/_\lambda `$ in Ref. References. and similarly for all $`A(H)/H`$ if our conjecture in that subsection can be proven.
## 5 The $`\widehat{T}\widehat{J}`$ OPE’s of $`A(H)/H`$
The OPE isomorphisms also allow us to compute the $`\widehat{T}_\sigma \widehat{J}(\sigma )`$ OPE’s of orbifold theory.
To begin this discussion, we remind the reader of the $`TJ`$ OPE of the general affine-Virasoro construction$`^{\text{References},\text{References},\text{References},\text{References},\text{References}}`$
$$T(z)J_a(w)=M(L)_a^b(\frac{1}{(zw)^2}+\frac{_w}{zw})J_b(w)+\frac{N(L)_a^{bc}:J_b(w)J_c(w):}{zw}+O((zw)^0)$$
(5.1a)
$$[L(m),J_a(n)]=nM(L)_a^bJ_b(m+n)+N(L)_a^{bc}\underset{p}{}:J_b(p)J_c(m+np):$$
(5.1b)
$$M(L)_a^b=2G_{ac}L^{cb}+f_{ad}^eL^{dc}f_{ce}^b,N(L)_a^{bc}=if_{ad}^{(b}L^{c)d}$$
(5.1c)
$$M(L)_{ab}=M(L)_a^cG_{cd}=M(L)_{ba},N(L)_a^{bc}=N(L)_a^{cb}$$
(5.1d)
$$M(L)^{}=\rho M(L)\rho ^{}=M(L)^{},N(L)_a^{bc}=\rho _a^dN(L)_d^{ef}\rho _c^b\rho _f^c.$$
(5.1e)
In the special case of the $`H`$-invariant CFT’s $`A(H)`$, one finds that the $`H`$-invariances of the tensors $`M(L)`$ and $`N(L)`$
$$[\omega ,M(L)]=0,\omega _a^dN(L)_d^{ef}(\omega ^{})_e^b(\omega ^{})_f^c=N(L)_a^{bc},\omega HAut(g)$$
(5.2)
follow from (5.1c), (2.6c) and (2.7a).
For the general orbifold $`A(H)/H`$, the OPE isomorphisms for $`\widehat{J}(\sigma )`$ and $`\widehat{T}_\sigma `$ in (2.23a) and (2.32) combine to give the derived OPE isomorphism
$$T(z)𝒥_{n(r)\mu }(w)\begin{array}{c}\hfill \\ ^\sigma \hfill \end{array}\widehat{T}_\sigma (z)\widehat{J}_{n(r)\mu }(w)$$
(5.3)
and this isomorphism gives the general $`\widehat{T}_\sigma \widehat{J}(\sigma )`$ OPE in sector $`\sigma `$ of $`A(H)/H`$
$`\widehat{T}_\sigma (z)\widehat{J}_{n(r)\mu }(w)`$ $`=`$ $`_{n(r)\mu }^{n(s)\nu }(\sigma )[{\displaystyle \frac{1}{(zw)^2}}+{\displaystyle \frac{_w}{(zw)}}]\widehat{J}_{n(s)\nu }(w)`$
$`+{\displaystyle \frac{𝒩_{n(r)\mu }^{n(s)\nu ;n(t)\delta }(\sigma ):\widehat{J}_{n(s)\nu }(w)\widehat{J}_{n(t)\delta }(w):}{zw}}+O((zw)^0)`$
$$_{n(r)\mu }^{n(s)\nu }(\sigma )=\chi \left(\sigma \right)_{n\left(r\right)\mu }\chi \left(\sigma \right)_{n\left(s\right)\nu }^1U(\sigma )_{n(r)\mu }^aM(L)_a^bU^{}(\sigma )_b^{n(s)\nu }$$
(5.4b)
$$𝒩_{n(r)\mu }^{n(s)\nu ;n(t)\delta }(\sigma )=\frac{\chi \left(\sigma \right)_{n\left(r\right)\mu }}{\chi \left(\sigma \right)_{n\left(s\right)\nu }\chi \left(\sigma \right)_{n\left(t\right)\delta }}U(\sigma )_{n(r)\mu }^aN(L)_a^{bc}U^{}(\sigma )_b^{n(s)\nu }U^{}(\sigma )_c^{n(t)\delta }.$$
(5.4c)
The twisted tensors $`(\sigma )`$ and $`𝒩(\sigma )`$ in (5.4b,5.4b) are the dual in sector $`\sigma `$ of the untwisted tensors $`M(L)`$ and $`N(L)`$
$$M(L)\begin{array}{c}\hfill \\ ^\sigma \hfill \end{array}(M(L);\sigma ),N(L)\begin{array}{c}\hfill \\ ^\sigma \hfill \end{array}𝒩(N(L);\sigma )$$
(5.5)
and the same result (5) can be obtained by direct OPE calculation (see App. Appendix D.) in the twisted sectors of the orbifolds. The orbifold duality transformations (5.5) (and their realization in (5.4b,5.4b)) are on the same footing as the earlier ones for $`𝒢,,`$ and $``$.
One also finds the $``$\- and $`𝒩`$-selection rules and their solutions
$$_{n(r)\mu }^{n(s)\nu }(\sigma )(1E_{n(r)}(\sigma )E_{n(s)}(\sigma )^{})=0$$
(5.6a)
$$_{n(r)\mu }^{n(s)\nu }(\sigma )=\delta _{n(r)n(s),0\text{ mod }\rho (\sigma )}_{n(r)\mu }^{n(r)\nu }(\sigma )$$
(5.6b)
$$𝒩_{n(r)\mu }^{n(s)\nu ;n(t)\delta }(\sigma )(1E_{n(r)}(\sigma )E_{n(s)}(\sigma )^{}E_{n(t)}(\sigma )^{})=0$$
(5.6c)
$$𝒩_{n(r)\mu }^{n(s)\nu ;n(t)\delta }(\sigma )=\delta _{n(r)n(s)n(t),0\text{ mod }\rho (\sigma )}𝒩_{n(r)\mu }^{n(s)\nu ;n(r)n(s),\delta }(\sigma )$$
(5.6d)
$$𝒩_{n(r)\mu }^{n(s)\nu ;n(r)n(s),\delta }(\sigma )=0\text{unless }n(r)n(s)\{n(r)\}$$
(5.6e)
which are the dual in sector $`\sigma `$ of the $`H`$-invariances (5.2) of $`M(L)`$ and $`N(L)`$ in the untwisted sector.
Another consequence of the $`H`$-invariances in (5.2) is that the twisted tensors $`(\sigma )`$ and $`𝒩(\sigma )`$ are class functions
$$(U(\sigma )v(\sigma );\sigma )=(U(\sigma );\sigma ),𝒩(U(\sigma )v(\sigma );\sigma )=𝒩(U(\sigma );\sigma )$$
(5.7)
under the $`H`$-conjugation in (g).
Using the $``$\- and $`𝒩`$-selection rules in (5), one obtains the $`\widehat{T}_\sigma \widehat{J}(\sigma )`$ OPE’s in the more informative presentation
$`\widehat{T}_\sigma (z)\widehat{J}_{n(r)\mu }(w)`$ $`=`$ $`_{n(r)\mu }^{n(r)\nu }(\sigma )[{\displaystyle \frac{1}{(zw)^2}}+{\displaystyle \frac{_w}{(zw)}}]\widehat{J}_{n(r)\nu }(w)`$
$`+{\displaystyle \frac{𝒩_{n(r)\mu }^{n(s)\nu ;n(r)n(s),\delta }(\sigma ):\widehat{J}_{n(s)\nu }(w)\widehat{J}_{n(r)n(s),\delta }(w):}{zw}}+O((zw)^0)`$
which shows the consistency of the monodromies on both sides of the relation. The associated commutator<sup>p</sup><sup>p</sup>pThe complex conjugates $`^{}`$ and $`𝒩^{}`$ can be computed with (5.1e), and the results follow the pattern seen in Eq. (l): a factor of $`^{}`$ or $``$ for an up or down index respectively. Then it is easily checked that (5.9) is consistent with the adjoint operations in (4.4a) and (4.7).
$$[L_\sigma (m),\widehat{J}_{n(r)\mu }(n+\frac{n\left(r\right)}{\rho \left(\sigma \right)})]=(n+\frac{n\left(r\right)}{\rho \left(\sigma \right)})_{n(r)\mu }^{n(r)\nu }(\sigma )\widehat{J}_{n(r)\nu }(m+n+\frac{n\left(r\right)}{\rho \left(\sigma \right)})$$
(5.9)
$$+𝒩_{n(r)\mu }^{n(s)\nu ;n(r)n(s),\delta }(\sigma )\underset{p}{}:\widehat{J}_{n(s)\nu }(p+\frac{n\left(s\right)}{\rho \left(\sigma \right)})\widehat{J}_{n(r)n(s),\delta }(m+np+\frac{n\left(r\right)n\left(s\right)}{\rho \left(\sigma \right)}):$$
follows from (5).
## 6 The Orbifolds of the ($`H`$ and Lie $`h`$)-invariant CFT’s
Our presentation of these “doubly-invariant” CFT’s and their orbifolds will follow the progression
$$A(\text{Lie}h)A(\text{Lie}h(H))\frac{A(\text{Lie}h(H))}{H}$$
(6.1)
starting with a review of Lie symmetry in current-algebraic conformal field theory.
### 6.1 The Lie $`h`$-invariant CFT’s
The Lie $`h`$-invariant CFT’s$`^{\text{References}\text{References},\text{References},\text{References}}`$ on $`g`$, called collectively $`A(\text{Lie}h)`$, are those CFT’s with a Lie symmetry<sup>q</sup><sup>q</sup>qThe Lie $`h`$-invariant CFT’s $`A(\text{Lie}h)`$ are also invariant (at least) under the connected part $`(\text{Lie}H)_c`$ of the corresponding Lie group Lie $`H`$. Moreover $`(\text{Lie}H)_c\text{Lie}GAut(g)`$ and $`(\text{Lie}H)_cAut(h)`$, so one may consider the orbifolds by a Lie group $`A(\text{Lie}h)/(\text{Lie}H)_c`$ for any Lie $`h`$-invariant CFT. We shall not do so here. $`hg`$, which may be realized globally or locally. Large numbers of Lie $`h`$-invariant CFT’s are known,$`^{\text{References}}`$ and simple examples of $`A(\text{Lie}h)`$ include the WZW model $`A_g(\text{Lie}g)`$, described by the affine-Sugawara construction$`^{\text{References},\text{References}\text{References},\text{References}}`$ on $`g`$, and the $`g/h`$ coset constructions.$`^{\text{References},\text{References},\text{References},\text{References}}`$
To describe this class of CFT’s more precisely, we begin with the index decomposition
$$a=(A,I),A\text{ in }h,I\text{ in }g/h:G_{AI}=f_{AB}^I=f_{AI}^B=0$$
(6.2a)
$$J_A(z)J_B(w)=\frac{G_{AB}}{(zw)^2}+\frac{if_{AB}^CJ_C(w)}{zw}+O((zw)^0)$$
(6.2b)
where $`h`$ is any reductive subalgebra of the ambient algebra $`g`$ and $`J_h=\{J_A\}`$ are the $`h`$ currents. The Lie $`h`$-invariant CFT’s $`A(\text{Lie}h)`$ are those CFT’s on $`g`$ whose inverse inertia tensor is invariant under infinitesimal transformations generated by Lie $`h`$
$$T=L^{ab}:J_aJ_b:,\delta _AL^{ab}=L^{c(a}f_{cA}^{b)}=iN(L)_A^{ab}=0,A=1,\mathrm{},\text{dim}h$$
(6.3)
where $`N(L)`$ is defined in (5). The inverse inertia tensors of $`A(\text{Lie}h)`$ satisfy a consistent$`^{\text{References}}`$ reduced Virasoro master equation.
At least for all unitary CFT’s in $`A(\text{Lie}h)`$, it follows that$`^{\text{References},\text{References}}`$
$$M(L)_A^I=0,M(L)_A^CM(L)_C^B=M(L)_A^B$$
(6.4a)
$$T(z)J_A(w)=M(L)_A^B(\frac{1}{(zw)^2}+\frac{_w}{zw})J_B(w)+O((zw)^0)$$
(6.4b)
where $`M(L)_A^B`$ is the $`h`$ block of $`M(L)`$ (see Eq. (5)). As seen in (6.1), the essential simplifications of the Lie $`h`$-invariant CFT’s are that a) $`M(L)_A^B`$ is a projector and b) there is no two-current term on the right side of (6.4b). It then follows$`^{\text{References},\text{References}}`$ in a left eigenbasis of $`M(L)_A^B`$ that the $`h`$ currents are either $`(1,0)`$ or $`(0,0)`$ operators of $`T(z)`$. For the affine-Sugawara constructions, one has $`h=g`$ and all the $`h`$ currents are $`(1,0)`$ operators, so that Lie $`h`$ is realized globally. For the $`g/h`$ coset constructions the $`h`$ currents are $`(0,0)`$ operators, so that Lie $`h`$ is realized locally. For general Lie $`h`$-invariant CFT’s, the affine $`h`$ subalgebras (6.2b) decompose as$`^{\text{References}}`$
$$h=h_0h_1$$
(6.5)
where $`h_0`$ and $`h_1`$ are the closed affine subalgebras of the $`(0,0)`$ and the $`(1,0)`$ operators respectively.<sup>r</sup><sup>r</sup>rIn what follows, we will generally assume unitarity of the Lie $`h`$-invariant CFT’s and hence the Lie $`h`$ relations (6.1). For the affine-Sugawara and coset constructions, however, it is clear that (6.1) holds without assumption of unitarity. Section 9 gives further evidence that the relations (6.1) may follow directly from the Virasoro master equation.
### 6.2 The ($`H`$ and Lie $`h`$)-invariant CFT’s
We consider next the ($`H`$ and Lie $`h`$)-invariant CFT’s, called collectively $`A(\text{Lie}h(H))`$,
$$A(\text{Lie}h(H))A(\text{Lie}h),A(\text{Lie}h(H))A(H)$$
(6.6a)
$$hg,HAut(g),HAut(h).$$
(6.6b)
These are the CFT’s which are simultaneously invariant under some Lie $`h`$ and also a finite group $`H`$, and the CFT’s of $`A(\text{Lie}h(H))`$ can be used to from the orbifolds $`A(\text{Lie}h(H))/H`$. Because $`g`$, $`h`$ and $`H`$ are not fixed, these doubly-invariant CFT’s provide a different slice of orbifold theory than the large examples in Subsec. 2.6 and Secs. 8 and 9.
We will not attempt to classify the ($`H`$ and Lie $`h`$)-invariant CFT’s, but the set is large, including the general WZW model
$$A_g(H)=A_g(\text{Lie}g)=A(\text{Lie}g(H)),HAut(g)$$
(6.7)
which is described by the affine-Sugawara construction$`^{\text{References},\text{References}\text{References},\text{References}}`$ on $`g`$, and the general $`H`$-invariant coset construction$`^{\text{References},\text{References}}`$
$$\frac{g}{h}(H)A(\text{Lie}h(H)).$$
(6.8)
Many other examples in $`A(\text{Lie}h(H))`$ are known including $`a)`$ the level families that live on the Lie $`h`$-invariant graphs with a graph symmetry$`^{\text{References},\text{References},\text{References}}`$ and $`b)`$ the Lie $`(h=g_{diag})`$-invariant CFT’s$`^{\text{References},\text{References}}`$ in $`A(_\lambda (\text{permutation}))`$, whose stress tensors describe the untwisted sectors of the $`G_{diag(\sigma )}`$-invariant cyclic orbifolds.
Another large class of doubly-invariant CFT’s $`A(\text{Cartan}g(H(d)))`$ is discussed in Sec. 9, where $`H=H(d)`$ is any group of inner automorphisms. These doubly-invariant CFT’s are particularly important because, as we shall see, they underlie the connection between inner-automorphic orbifolds and stress-tensor spectral flow.
For CFT’s in $`A(\text{Lie}h(H))`$, the subalgebra $`hg`$ is an $`H`$-covariant subalgebra$`^{\text{References}}`$
$$\omega (h_\sigma )_A^I=0,J_A^{}=\omega (h_\sigma )_A^BJ_B,\omega (h_\sigma )H,A=1,\mathrm{},\text{dim}h$$
(6.9)
of the $`H`$-covariant ambient algebra $`g`$. Here $`\omega (h_\sigma )_A^B`$ is the $`h`$ block of $`\omega (h_\sigma )`$ and $`J_A^{}`$ satisfies the same OPE as $`J_A`$, given in (6.2b). Then, using (2.6b,2.6b) and (6.2b), the properties
$$\omega (h_\sigma )_A^C\omega ^{}(h_\sigma )_C^B=\delta _A^B,\omega (h_\sigma )_I^A=0$$
(6.10a)
$$\omega (h_\sigma )_A^C\omega (h_\sigma )_B^DG_{CD}=G_{AB},\omega (h_\sigma )_A^D\omega (h_\sigma )_B^Ef_{DE}^F\omega ^{}(h_\sigma )_F^C=f_{AB}^C$$
(6.10b)
are obtained for all $`h_\sigma `$. The relations in (6.10b) express the $`H`$-invariance of the $`h`$ tensors $`G_{AB}`$ and $`f_{AB}^C`$. It follows that each matrix action $`\omega (h_\sigma )`$ is block diagonal
$$\omega (h_\sigma )=\left(\begin{array}{cc}\omega (h_\sigma )_A^B& 0\\ 0& \omega (h_\sigma )_I^J\end{array}\right)$$
(6.11)
and that the $`h`$ and $`g/h`$ blocks $`\omega (h_\sigma )_A^B`$ and $`\omega (h_\sigma )_I^J`$ of $`\omega (h_\sigma )`$ are separately unitary. Then the unitary eigenvector matrices $`U(\sigma )`$ and $`U^{}(\sigma )`$ may also be taken block diagonal with each block separately unitary. Finally, Eqs. (5.2) and (6.11) imply the vanishing commutator
$$\omega (h_\sigma )_A^CM(L)_C^B=M(L)_A^C\omega (h_\sigma )_C^B$$
(6.12)
among the $`h`$ blocks of $`\omega (h_\sigma )`$ and $`M(L)`$.
### 6.3 The orbifolds $`A(\text{Lie}h(H))/H`$
We turn now to the orbifolds by $`H`$ of the doubly-invariant CFT’s $`A(\text{Lie}h(H))`$
$$\frac{A(\text{Lie}h(H))}{H}\frac{A(H)}{H},hg,HAut(g),HAut(h).$$
(6.13)
The twisted $`g`$ currents $`\widehat{J}_{n(r)\mu }`$ of these orbifolds satisfy the twisted current algebra
$$\widehat{𝔤}(HAut(g),HAut(h);\sigma )$$
(6.14)
whose form is included in (3.1). In this case, the ambient algebra (6.14) has a twisted $`h`$ subalgebra$`^{\text{References}}`$ generated by the twisted $`h`$ currents $`\widehat{J}_{\widehat{𝔥}(\sigma )}`$, which we will discuss below. When Lie $`h`$ is realized locally, one must learn to gauge the twisted $`h`$ currents in an action formulation of these orbifolds.
To study the twisted $`h`$ currents in a “conformal weight” basis (see Eq. (6.3)), we begin with the simultaneous left-eigenvalue problem
$$U(\sigma )_{\widehat{n}(r)\theta _i\mu }^B\omega (h_\sigma )_B^A=E_{\widehat{n}(r)}(\sigma )U(\sigma )_{\widehat{n}(r)\theta _i\mu }^A$$
(6.15a)
$$U(\sigma )_{\widehat{n}(r)\theta _i\mu }^BM(L)_B^A=\theta _iU(\sigma )_{\widehat{n}(r)\theta _i\mu }^A$$
(6.15b)
$$E_{\widehat{n}(r)}(\sigma )=e^{\frac{2\pi i\widehat{n}(r)}{\rho (\sigma )}},\widehat{n}(r)\{n(r)\},\theta _i\{0,1\},\sigma =0,\mathrm{},N_c1$$
(6.15c)
defined on the (unitary) $`h`$ block of $`U`$. Eq. (6.15a) is the induced $`H`$-eigenvalue problem$`^{\text{References}}`$ for the $`H`$-covariant subalgebra $`h`$, and the consistency of the simultaneous eigenvalue problem (6.3) is guaranteed by the vanishing commutator in Eq. (6.12). The allowed values of $`\theta _i`$ in (6.15c) follow from $`M^2=M`$ in (6.4a), and we emphasize that the matrix $`M(L)`$ and its eigenvalues $`\theta _i(\sigma )=\theta _i(0)=\theta _i`$ are independent of $`\sigma `$.
The analysis$`^{\text{References}}`$ reviewed in Sec. 2 for the ambient algebra $`g`$ can now be applied, mutatis mutandis, to the $`H`$-covariant subalgebra $`h`$. With (6.1) and (6.3), we find that the $`h`$-eigencurrents $`𝒥_h(\sigma )=\{𝒥_{\widehat{n}(r)\theta _i\mu }\}`$ of sector $`\sigma `$
$$𝒥_{\widehat{n}(r)\theta _i\mu }(z)\chi \left(\sigma \right)_{\widehat{n}\left(r\right)\theta _i\mu }U(\sigma )_{\widehat{n}(r)\theta _i\mu }^AJ_A(z),𝒥_{\widehat{n}(r)\theta _i\mu }(z)^{}=E_{\widehat{n}(r)}(\sigma )𝒥_{\widehat{n}(r)\theta _i\mu }(z)$$
(6.16a)
$$T(z)𝒥_{\widehat{n}(r)\theta _i\mu }(w)=\theta _i(\frac{1}{(zw)^2}+\frac{_w}{zw})𝒥_{\widehat{n}(r)\theta _i\mu }(w)+O((zw)^0)$$
(6.16b)
have conformal weight $`\theta _i\{0,1\}`$ and simultaneously a diagonal response $`E_{\widehat{n}(r)}(\sigma )`$ to the automorphism group $`H`$.
Then using the OPE’s (6.2b) of the $`h`$ current algebra, the $`h`$-eigencurrents (6.16a) and the OPE isomorphism
$$𝒥_h(\sigma )\begin{array}{c}\hfill \\ ^\sigma \hfill \end{array}\widehat{J}_{\widehat{𝔥}(\sigma )},\widehat{J}_{\widehat{𝔥}(\sigma )}=\{\widehat{J}_{\widehat{n}(r)\theta _i\mu }\}$$
(6.17a)
$$\text{automorphisms }E_{\widehat{n}(r)}(\sigma )\begin{array}{c}\hfill \\ ^\sigma \hfill \end{array}\text{monodromies }E_{\widehat{n}(r)}(\sigma )$$
(6.17b)
one finds the twisted $`h`$ current system (read Eq. (2.24a) with $`n(r)\mu \widehat{n}(r)\theta _i\mu `$). The twisted $`h`$ tensors of this system are
$$𝒢_{\widehat{n}(r)\theta _i\mu ;\widehat{n}(s)\theta _j\nu }(\sigma )=\chi \left(\sigma \right)_{\widehat{n}\left(r\right)\theta _i\mu }\chi \left(\sigma \right)_{\widehat{n}\left(s\right)\theta _j\nu }U(\sigma )_{\widehat{n}(r)\theta _i\mu }^AU(\sigma )_{\widehat{n}(s)\theta _j\nu }^BG_{AB}$$
(6.18a)
$$_{\widehat{n}(r)\theta _i\mu ;\widehat{n}(s)\theta _j\nu }^{\widehat{n}(t)\theta _k\delta }(\sigma )=\frac{\chi \left(\sigma \right)_{\widehat{n}\left(r\right)\theta _i\mu }\chi \left(\sigma \right)_{\widehat{n}\left(s\right)\theta _j\nu }}{\chi \left(\sigma \right)_{\widehat{n}\left(t\right)\theta _k\delta }}U(\sigma )_{\widehat{n}(r)\theta _i\mu }^AU(\sigma )_{\widehat{n}(s)\theta _j\nu }^Bf_{AB}^CU^{}(\sigma )_C^{\widehat{n}(t)\theta _k\delta }$$
(6.18b)
$`_{\widehat{n}\left(r\right)\theta _i\mu ;\widehat{n}\left(s\right)\theta _j\nu }^{\widehat{n}\left(u\right)\theta _lϵ}\left(\sigma \right)_{\widehat{n}\left(t\right)\theta _k\delta ;\widehat{n}\left(u\right)\theta _lϵ}^{\widehat{n}\left(v\right)\theta _m\gamma }\left(\sigma \right)+_{\widehat{n}\left(s\right)\theta _j\nu ;\widehat{n}\left(t\right)\theta _k\delta }^{\widehat{n}\left(u\right)\theta _lϵ}\left(\sigma \right)_{\widehat{n}\left(r\right)\theta _i\mu ;\widehat{n}\left(u\right)\theta _lϵ}^{\widehat{n}\left(v\right)\theta _m\gamma }\left(\sigma \right)`$
$$+_{\widehat{n}(t)\theta _k\delta ;\widehat{n}(r)\theta _i\mu }^{\widehat{n}(u)\theta _lϵ}(\sigma )_{\widehat{n}(s)\theta _j\nu ;\widehat{n}(u)\theta _lϵ}^{\widehat{n}(v)\theta _m\gamma }(\sigma )=0$$
(6.18c)
$$_{\widehat{n}(r)\theta _i\mu ;\widehat{n}(s)\theta _j\nu ;\widehat{n}(t)\theta _k\delta }(\sigma )_{\widehat{n}(r)\theta _i\mu ;\widehat{n}(s)\theta _j\nu }^{\widehat{n}(u)\theta _lϵ}(\sigma )𝒢_{\widehat{n}(u)\theta _lϵ;\widehat{n}(t)\theta _k\delta }(\sigma )=_{\widehat{n}(r)\theta _i\mu ;\widehat{n}(t)\theta _k\delta ;\widehat{n}(s)\theta _j\nu }(\sigma )$$
(6.18d)
$$𝒢_{\widehat{n}(r)\theta _i\mu ;\widehat{n}(s)\theta _j\nu }(\sigma )(1E_{\widehat{n}(r)}(\sigma )E_{\widehat{n}(s)}(\sigma ))=0$$
(6.18e)
$$_{\widehat{n}(r)\theta _i\mu ;\widehat{n}(s)\theta _j\nu }^{\widehat{n}(t)\theta _k\delta }(\sigma )(1E_{\widehat{n}(r)}(\sigma )E_{\widehat{n}(s)}(\sigma )E_{\widehat{n}(t)}(\sigma )^{})=0$$
(6.18f)
and the corresponding twisted $`h`$ subalgebra is
$$\widehat{𝔥}(\sigma )\widehat{𝔥}(HAut(h);\sigma )\widehat{𝔤}(HAut(g),HAut(h);\sigma )$$
(6.19a)
$$[\widehat{J}_{\widehat{n}(r)\theta _i\mu }(m+\frac{\widehat{n}\left(r\right)}{\rho \left(\sigma \right)}),\widehat{J}_{\widehat{n}(s)\theta _j\nu }(n+\frac{\widehat{n}\left(s\right)}{\rho \left(\sigma \right)})]=(m+\frac{\widehat{n}\left(r\right)}{\rho \left(\sigma \right)})\delta _{m+n+\frac{\widehat{n}(r)+\widehat{n}(s)}{\rho (\sigma )},0}𝒢_{\widehat{n}(r)\theta _i\mu ;\widehat{n}(r),\theta _j\nu }(\sigma )$$
$$+i_{\widehat{n}(r)\theta _i\mu ;\widehat{n}(s)\theta _j\nu }^{\widehat{n}(r)+\widehat{n}(s),\theta _k\delta }(\sigma )\widehat{J}_{\widehat{n}(r)+\widehat{n}(s),\theta _k\delta }(m+n+\frac{\widehat{n}\left(r\right)+\widehat{n}\left(s\right)}{\rho \left(\sigma \right)})$$
(6.19b)
$$\mathrm{\#}\{\widehat{J}_{\widehat{𝔥}(\sigma )}\}=\mathrm{\#}\{J_h\}=\text{dim}h.$$
(6.19c)
The subalgebra $`\widehat{𝔥}(\sigma )`$ is the dual in sector $`\sigma `$ of the untwisted affine $`h`$ subalgebra in (6.2b), and the Jacobi identity of $`\widehat{𝔥}(\sigma )`$ follows with (6.18c,6.18c).
We turn next to the twisted inverse inertia tensor $`(\sigma )`$ in the stress tensor $`\widehat{T}_\sigma `$ of each orbifold in $`A(\text{Lie}h(H))/H`$. In addition to the selection rules (2.5) – which are dual to the $`H`$-symmetry of $`A(\text{Lie}h(H))`$ – we find that $`(\sigma )`$ also satisfies the condition<sup>s</sup><sup>s</sup>sA special case of this condition was first seen for cyclic permutation orbifolds in Ref. References.
$$^{n(u)\gamma ;(n(s)\nu }(\sigma )_{n(u)\gamma ;\widehat{n}(r)\theta _i\mu ;}^{n(t)\delta )}(\sigma )=i𝒩_{\widehat{n}(r)\theta _i\mu }^{n(s)\nu ;n(t)\delta }(\sigma )=0$$
(6.20)
which is dual to the Lie $`h`$ condition (6.3). Moreover, the $`\widehat{T}_\sigma \widehat{J}_{\widehat{𝔥}(\sigma )}`$ relations
$$_{\widehat{n}(r)\theta _i\mu }^{\widehat{n}(s)\theta _j\nu }(\sigma )=\theta _i\delta _{\widehat{n}(r)\theta _i\mu }^{\widehat{n}(s)\theta _j\nu }$$
(6.21a)
$$\widehat{T}_\sigma (z)\widehat{J}_{\widehat{n}(r)\theta _i\mu }(w)=\theta _i(\frac{1}{(zw)^2}+\frac{_w}{zw})\widehat{J}_{\widehat{n}(r)\theta _i\mu }(w)+O((zw)^0)$$
(6.21b)
$$[L_\sigma (m),\widehat{J}_{\widehat{n}(r)\theta _i\mu }(n+\frac{\widehat{n}\left(r\right)}{\rho \left(\sigma \right)})]=\theta _i(n+\frac{\widehat{n}\left(r\right)}{\rho \left(\sigma \right)})\widehat{J}_{\widehat{n}(r)\theta _i\mu }(m+n+\frac{\widehat{n}\left(r\right)}{\rho \left(\sigma \right)})$$
(6.21c)
can be obtained from (6.16b), (6.3) and the derived isomorphism (2.32), or as a special case of (5) and (5.9). These results show that each twisted $`h`$ current of any orbifold $`A(\text{Lie}h(H))/H`$ is either$`^{\text{References}}`$ a twisted $`(1,0)`$ operator (defined by $`\theta _i=1`$) or a twisted $`(0,0)`$ operator (defined by $`\theta _i=0`$). Looking back, we see that the properties (6.20) are orbifold reflections of the Lie $`h`$ properties in (6.1) and (6.16b).
### 6.4 WZW orbifolds, coset orbifolds and orbifold $`K`$-conjugation
We work out here the general WZW orbifold$`^{\text{References}}`$ and the general coset orbifold,$`^{\text{References}}`$ which provide simple sets of examples in $`A(\text{Lie}h(H))/H`$. In this discussion, we emphasize the role of the WZW orbifolds in orbifold $`K`$-conjugation,$`^{\text{References}}`$ and the role of orbifold $`K`$-conjugation in the construction of the coset orbifolds.
The affine-Sugawara construction$`^{\text{References},\text{References}\text{References},\text{References}}`$ on $`g`$
$$A_g(H)=A(\text{Lie}g(H)),HAut(h=g)$$
(6.22a)
$$T_g=L_g^{ab}:J_aJ_b:,L_g^{ab}=_I\frac{\eta _I^{ab}}{2k_I+Q_I}$$
(6.22b)
$$L_g^{cd}\omega _c^a\omega _c^b=L_g^{ab},\omega HAut(h=g)$$
(6.22c)
$$N(L_g)_a^{bc}=0,M(L_g)_a^b=\delta _a^b$$
(6.22d)
$$[L_g(m),J_a(n)]=nJ_a(m+n)$$
(6.22e)
is always a (global Lie$`h=g`$)-invariant CFT which is also $`H`$-invariant under any $`HAut(g)`$. Then the description of the general WZW orbifold$`^{\text{References}}`$
$$\frac{A_g(H)}{H}=\frac{A(\text{Lie}g(H))}{H},\widehat{𝔥}(\sigma )=\widehat{𝔤}(\sigma )$$
(6.23a)
$$\widehat{T}_{\widehat{𝔤}(\sigma )}=_{\widehat{𝔤}(\sigma )}^{n(r)\mu ;n(r),\nu }(\sigma ):\widehat{J}_{n(r)\mu }\widehat{J}_{n(r),\nu }:,\sigma =0,\mathrm{},N_c1$$
(6.23b)
$$_{\widehat{𝔤}(\sigma )}^{n(r)\mu ;n(r),\nu }(\sigma )=\chi \left(\sigma \right)_{n\left(r\right)\mu }^1\chi \left(\sigma \right)_{n\left(r\right),\nu }^1L_g^{ab}U^{}(\sigma )_a^{n(r)\mu }U^{}(\sigma )_b^{n(r),\nu }$$
(6.23c)
$$𝒩_{n(r)\mu }^{n(s)\nu ;n(t)\delta }(\sigma )=0,_{n(r)\mu }^{n(s)\nu }(\sigma )=\delta _{n(r)\mu }^{n(s)\nu }$$
(6.23d)
$$[L_\sigma ^{\widehat{𝔤}(\sigma )}(m),\widehat{J}_{n(r)\mu }(n+\frac{n\left(r\right)}{\rho \left(\sigma \right)})]=(n+\frac{n\left(r\right)}{\rho \left(\sigma \right)})\widehat{J}_{n(r)\mu }(m+n+\frac{n\left(r\right)}{\rho \left(\sigma \right)})$$
(6.23e)
is obtained from the results of the previous subsection. In the WZW orbifolds, the twisted $`g`$ currents $`\{\widehat{J}(\sigma )\}`$ satisfy the general twisted current algebra $`\widehat{𝔤}(\sigma )`$ in (3.1) and each of the twisted $`g`$ currents is a twisted (1,0) operator (with $`\theta _i(\sigma )=\theta _i=1`$) under the orbifold affine-Sugawara construction $`\widehat{T}_{\widehat{𝔤}(\sigma )}`$. The case of the WZW cyclic permutation orbifolds was discussed in Ref. References and further discussion of WZW orbifolds is found in Secs. 7 and 9.
Returning for a moment to the general affine-Virasoro construction (2.1), we remind the reader that the affine-Sugawara construction $`T_g`$ also plays a central role in the operation known as $`K`$-conjugation$`^{\text{References},\text{References},\text{References},\text{References},\text{References},\text{References}}`$
$$A\begin{array}{c}\hfill \\ ^K\hfill \end{array}\stackrel{~}{A}$$
(6.24a)
$$T_g(z)=T(z)+\stackrel{~}{T}(z),\stackrel{~}{T}(z)T(w)=O((zw)^0)$$
(6.24b)
$$L_g^{ab}=L^{ab}+\stackrel{~}{L}^{ab},c_g=c+\stackrel{~}{c}$$
(6.24c)
which relates $`K`$-conjugate pairs $`A`$, $`\stackrel{~}{A}`$ of current-algebraic CFT’s on $`g`$. The simplest example of $`K`$-conjugation is the $`g/h`$ coset construction$`^{\text{References},\text{References},\text{References},\text{References}}`$
$$T_{g/h}=T_gT_h,c_{g/h}=c_gc_h,[L_{g/h}(m),J_A(n)]=0,A=1,\mathrm{},\text{dim}h$$
(6.25)
which is a (local Lie $`h`$)-invariant CFT (with $`\theta _i=0`$ for each $`h`$ current).
We emphasize that $`K`$-conjugation is closed on the space of $`H`$-invariant CFT’s
$$A(H)\begin{array}{c}\hfill \\ ^K\hfill \end{array}\stackrel{~}{A}(H),HAut(g)$$
(6.26)
because the affine-Sugawara construction $`T_g`$ describes the $`H`$-invariant CFT $`A_g(H)`$.
For the corresponding orbifolds by $`H`$, one finds that $`K`$-conjugation (6.26) and the duality transformation (2.33c) combine to give orbifold $`K`$-conjugation$`^{\text{References}}`$
$$\frac{A(H)}{H}\begin{array}{c}\hfill \\ ^K\hfill \end{array}\frac{\stackrel{~}{A}(H)}{H}$$
(6.27a)
$$\widehat{T}_{\widehat{𝔤}(\sigma )}(z)=\widehat{T}_\sigma (z)+\widehat{\stackrel{~}{T}}_\sigma (z),\widehat{\stackrel{~}{T}}_\sigma (z)\widehat{T}_\sigma (w)=O((zw)^0)$$
(6.27b)
$$_{\widehat{𝔤}(\sigma )}^{n(r)\mu ;n(s)\nu }(\sigma )=^{n(r)\mu ;n(s)\nu }(\sigma )+\stackrel{~}{}^{n(r)\mu ;n(s)\nu }(\sigma )$$
(6.27c)
$$\widehat{c}_{\widehat{𝔤}(\sigma )}=c_g,\widehat{c}(\sigma )=c,\widehat{\stackrel{~}{c}}(\sigma )=\stackrel{~}{c}$$
(6.27d)
which relates $`K`$-conjugate pairs of stress tensors in $`A(H)/H`$ through the orbifold affine-Sugawara construction $`\widehat{T}_{\widehat{𝔤}(\sigma )}`$.
The simplest example of orbifold $`K`$-conjugation, namely the general coset orbifold
$$\frac{\frac{g}{h}(H)}{H}\frac{A(\text{Lie}h(H))}{H},hg,HAut(h),HAut(g)$$
(6.28a)
$$T_{g/h}(z)=T_g(z)T_h(z)\begin{array}{c}\hfill \\ ^\sigma \hfill \end{array}\widehat{T}_{\frac{g/h}{H}}(z)_\sigma =\widehat{T}_{\widehat{𝔤}(\sigma )}\widehat{T}_{\widehat{𝔥}(\sigma )}$$
(6.28b)
$$T_h=L_h^{AB}:J_AJ_B:\begin{array}{c}\hfill \\ ^\sigma \hfill \end{array}\widehat{T}_{\widehat{𝔥}(\sigma )}=_{\widehat{𝔥}(\sigma )}^{\widehat{n}(r)0\mu ;\widehat{n}(r),0\nu }(\sigma ):\widehat{J}_{\widehat{n}(r)0\mu }\widehat{J}_{\widehat{n}(r),0\nu }:$$
(6.28c)
$$_{\widehat{𝔥}(\sigma )}^{\widehat{n}(r)0\mu ;\widehat{n}(s)0\nu }(\sigma )=\chi \left(\sigma \right)_{\widehat{n}\left(r\right)0\mu }^1\chi \left(\sigma \right)_{\widehat{n}\left(s\right)0\nu }^1L_h^{AB}U^{}(\sigma )_A^{\widehat{n}(r)0\mu }U^{}(\sigma )_B^{\widehat{n}(s)0\nu }$$
(6.28d)
is also found in $`A(\text{Lie}h(H))/H`$. The untwisted sectors of these orbifolds are formed from the $`H`$-invariant coset constructions $`\frac{g}{h}(H)`$, and the twisted $`h`$ currents $`\widehat{J}_{\widehat{𝔥}(\sigma )}=\{\widehat{J}_{\widehat{n}(r)0\mu }\}`$ have $`\theta _i(\sigma )=\theta _i=0`$ in this case because the untwisted $`h`$ currents $`J_A`$ are $`(0,0)`$ operators of $`T_{g/h}`$.
In further detail, the twisted $`h`$ currents are twisted $`(1,0)`$ operators under $`\widehat{T}_{\widehat{𝔤}(\sigma )}`$ and $`\widehat{T}_{\widehat{𝔥}(\sigma )}`$ and hence twisted $`(0,0)`$ operators under $`\widehat{T}_{(g/h)/H}`$:
$$[L_\sigma ^{\widehat{𝔤}(\sigma )}(m),\widehat{J}_{\widehat{n}(r)0\mu }(n+\frac{\widehat{n}\left(r\right)}{\rho \left(\sigma \right)})]=(n+\frac{\widehat{n}\left(r\right)}{\rho \left(\sigma \right)})\widehat{J}_{\widehat{n}(r)0\mu }(m+n+\frac{\widehat{n}\left(r\right)}{\rho \left(\sigma \right)})$$
(6.29a)
$$[L_\sigma ^{\widehat{𝔥}(\sigma )}(m),\widehat{J}_{\widehat{n}(r)0\mu }(n+\frac{\widehat{n}\left(r\right)}{\rho \left(\sigma \right)})]=(n+\frac{\widehat{n}\left(r\right)}{\rho \left(\sigma \right)})\widehat{J}_{\widehat{n}(r)0\mu }(m+n+\frac{\widehat{n}\left(r\right)}{\rho \left(\sigma \right)})$$
(6.29b)
$$[L_\sigma ^{\frac{g/h}{H}}(m),\widehat{J}_{\widehat{n}(r)0\mu }(n+\frac{\widehat{n}\left(r\right)}{\rho \left(\sigma \right)})]=0.$$
(6.29c)
The eigenvalue $`\theta _i=0`$ of the twisted $`h`$ currents controls the vanishing final commutator with the coset orbifold Virasoro generators. The results (6.28a,6.28a) and (6.29c) are equivalent to those given in Ref. References, which also discusses the orbifolds of the $`_\lambda `$-invariant coset constructions in further detail. Finally, the twisted $`h`$ currents satisfy the twisted $`h`$ subalgebra
$$[\widehat{J}_{\widehat{n}(r)0\mu }(m+\frac{\widehat{n}\left(r\right)}{\rho \left(\sigma \right)}),\widehat{J}_{\widehat{n}(s)0\nu }(n+\frac{\widehat{n}\left(s\right)}{\rho \left(\sigma \right)})]=(m+\frac{\widehat{n}\left(r\right)}{\rho \left(\sigma \right)})\delta _{m+n+\frac{\widehat{n}(r)+\widehat{n}(s)}{\rho (\sigma )},0}𝒢_{\widehat{n}(r)0\mu ;\widehat{n}(r),0\nu }(\sigma )$$
$$+i_{\widehat{n}(r)0\mu ;\widehat{n}(s)0\nu }^{\widehat{n}(r)+\widehat{n}(s),0\delta }(\sigma )\widehat{J}_{\widehat{n}(r)+\widehat{n}(s),0\delta }(m+n+\frac{\widehat{n}\left(r\right)+\widehat{n}\left(s\right)}{\rho \left(\sigma \right)})$$
(6.30)
which is dual to the untwisted $`h`$ algebra whose OPE’s are given in (6.2b).
Other applications of orbifold $`K`$-conjugation are found in the following subsection and Subsec. 9.5.
### 6.5 Decomposition of the twisted $`h`$ subalgebras
With the help of orbifold $`K`$-conjugation, we will show in this subsection that the twisted affine $`h`$ subalgebra $`\widehat{𝔥}(\sigma )`$ of each sector of each orbifold $`A(\text{Lie}h(H))/H`$ decomposes as
$$\widehat{𝔥}(\sigma )=\widehat{𝔥}_0(\sigma )\widehat{𝔥}_1(\sigma )$$
(6.31)
where $`\widehat{𝔥}_0(\sigma )`$ and $`\widehat{𝔥}_1(\sigma )`$ are the closed affine subalgebras of the twisted (0,0) and (1,0) operators respectively. The algebra $`\widehat{𝔥}(\sigma )\widehat{𝔤}(\sigma )`$ is given in (6.18f) and the orbifold argument given below parallels the argument in Ref. References for the decomposition (6.5) of the $`h`$ subalgebra of any Lie $`h`$-invariant CFT.
For the orbifolds $`A(\text{Lie}h(H))/H`$, we may combine orbifold $`K`$-conjugation (6.4) with the $`L_\sigma `$, $`\widehat{J}_{\widehat{𝔥}(\sigma )}`$ commutator in (6.20) to find the relation
$$\frac{A(\text{Lie}h(H))}{H}\begin{array}{c}\hfill \\ ^K\hfill \end{array}\frac{\stackrel{~}{A}(\text{Lie}h(H))}{H}:\theta _i^{\widehat{𝔤}(\sigma )}=1=\theta _i+\stackrel{~}{\theta }_i$$
(6.32)
among the “conformal weights” $`\theta `$ of the twisted $`h`$ currents of each $`K`$-conjugate orbifold pair in $`A(\text{Lie}h(H))/H`$.
In what follows, we consider various pairs of generators $`\widehat{J}_{\widehat{𝔥}(\sigma )}`$ of the twisted $`h`$ subalgebra $`\widehat{𝔥}(\sigma )`$ of sector $`\sigma `$ in any particular orbifold $`A(\text{Lie}h(H))/H`$, starting with an arbitrary pair of twisted $`(0,0)`$ operators. According to (6.20), the commutator of this pair of twisted currents commutes with the stress tensor $`\widehat{T}_\sigma `$ of this orbifold
$$[L_\sigma (m),[\widehat{J}_{\widehat{n}(r)0\mu }(n+\frac{\widehat{n}\left(r\right)}{\rho \left(\sigma \right)}),\widehat{J}_{\widehat{n}(s)0\nu }(p+\frac{\widehat{n}\left(s\right)}{\rho \left(\sigma \right)})]]=0$$
(6.33)
so the set of twisted $`(0,0)`$ operators of this orbifold is closed under commutation. Similarly the set of twisted $`(1,0)`$ operators of this orbifold is closed under commutation because, according to (6.32), these currents and their commutators are twisted $`(0,0)`$ operators of the $`K`$-conjugate stress tensors $`\widehat{\stackrel{~}{T}}_\sigma `$. To prove that the twisted $`(0,0)`$ and twisted $`(1,0)`$ operators of the orbifold commute, use the chain rule and the commutator (6.20) to see that
$$[L_\sigma (m),[\widehat{J}_{\widehat{n}(r)0\mu }(n+\frac{\widehat{n}\left(r\right)}{\rho \left(\sigma \right)}),\widehat{J}_{\widehat{n}(s)1\nu }(p+\frac{\widehat{n}\left(s\right)}{\rho \left(\sigma \right)})]]$$
$$=(p+\frac{\widehat{n}\left(s\right)}{\rho \left(\sigma \right)})[\widehat{J}_{\widehat{n}(r)0\mu }(n+\frac{\widehat{n}\left(r\right)}{\rho \left(\sigma \right)}),\widehat{J}_{\widehat{n}(s)1\nu }(p+m+\frac{\widehat{n}\left(s\right)}{\rho \left(\sigma \right)})].$$
(6.34)
Then use the twisted $`h`$ algebra (6.18f) on both sides, followed by the commutator (6.20) again on the left side. The result of this computation
$$𝒢_{n(r)0\mu ,n(s)1\nu }=_{n(r)0\mu ,n(s)1\nu }^{n(r)+n(s),\theta _k\delta }=0$$
(6.35)
establishes the decomposition (6.31) of the twisted $`h`$ subalgebra $`\widehat{𝔥}(\sigma )`$ of sector $`\sigma `$. As discussed in Subsec. 6.4, the general WZW orbifold and the general coset orbifold provide simple examples of this decomposition.
## 7 About Permutation Orbifolds
In this section, we discuss some features common to all permutation orbifolds
$$\frac{A(H)}{H},HS_N(\text{permutation})Aut(g),g=_{I=0}^{K1}𝔤^I,KN$$
(7.1)
where the copies $`𝔤^I𝔤`$ in the ($`HS_N`$) permutation-invariant CFT’s are taken at level $`k`$. In this case, the general twisted current algebra $`\widehat{𝔤}(\sigma )`$ of sector $`\sigma `$ consists of a set of commuting orbifold affine algebras$`^{\text{References},\text{References},\text{References},\text{References}}`$ at various orders (see also Sec. 8). We have also checked that the orbifold adjoint operation (4.4a) gives the standard$`^{\text{References},\text{References}}`$ adjoint operation (see Eq. (4.6b) and Sec. 8) for each commuting orbifold affine algebra in each sector $`\sigma `$, so that unitarity of the twisted sectors of any permutation orbifold $`A(H)/H`$ follows from unitarity of the permutation-invariant CFT $`A(H)`$.
From the orbifold induction procedure for orbifold affine algebras$`^{\text{References}}`$ we know that the ground state $`|0_\sigma `$ of sector $`\sigma `$ (which corresponds to the twist field of sector $`\sigma `$) is a twisted affine primary state
$$\widehat{J}_{n(r)\mu }((m+\frac{n\left(r\right)}{\rho \left(\sigma \right)})0)|0_\sigma =0,_\sigma 0|\widehat{J}_{n(r)\mu }((m+\frac{n\left(r\right)}{\rho \left(\sigma \right)})0)=0$$
(7.2)
where the second relation follows from the first by the orbifold adjoint operation. It follows from (7.2) that $`M`$ ordering, defined in (3.8), is a true normal ordering
$$\widehat{J}_{n(r)\mu }(z)_\sigma =:\widehat{J}_{n(r)\mu }(z)\widehat{J}_{n(s)\nu }(w):_M_\sigma =0$$
(7.3a)
$$\underset{p}{}:\widehat{J}_{n(r)\mu }(p+\frac{n\left(r\right)}{\rho \left(\sigma \right)})\widehat{J}_{n(r),\nu }(mp\frac{n\left(r\right)}{\rho \left(\sigma \right)}):_M|0_\sigma =0,\text{for}m0,n(r),\mu ,\nu $$
(7.3b)
for all permutation orbifolds. The relations (7.2) and (7.2) hold for all $`n(r)`$ and $`n(s)`$, not necessarily in the fundamental range. As a consequence of (7.3a), we obtain the twisted current-current correlators in each sector $`\sigma `$ of the general permutation orbifold
$$\widehat{J}_{n(r)\mu }(z)\widehat{J}_{n(s)\nu }(w)_\sigma =\delta _{n(r)+n(s),0\text{ mod }\rho (\sigma )}(\frac{w}{z})^{\frac{\overline{n}(r)}{\rho (\sigma )}}[\frac{1}{(zw)^2}+\frac{\overline{n}(r)/\rho (\sigma )}{w(zw)}]𝒢_{n(r)\mu ;n(r),\nu }(\sigma )$$
(7.4)
from the exact operator product (3.9). Correlators of more than two twisted currents can be computed from (3.1a), (3.1) and (7.2).
Using (3.13c), (7.2) and (7.3b), we may also compute the ground state conformal weight $`\widehat{\mathrm{\Delta }}_0(\sigma )`$ of sector $`\sigma `$
$$L_\sigma (m0)|0_\sigma =\delta _{m,0}\widehat{\mathrm{\Delta }}_0(\sigma )|0_\sigma $$
(7.5a)
$`\widehat{\mathrm{\Delta }}_0(\sigma )`$ $`=`$ $`{\displaystyle \underset{r,s,\mu ,\nu }{}}{\displaystyle \frac{\overline{n}(r)}{2\rho (\sigma )}}(1{\displaystyle \frac{\overline{n}(r)}{\rho (\sigma )}})^{n(r)\mu ;n(s)\nu }(\sigma )𝒢_{n(r)\mu ;n(s)\nu }(\sigma )`$ (7.5b)
$`=`$ $`{\displaystyle \underset{r,\mu ,\nu }{}}{\displaystyle \frac{\overline{n}(r)}{2\rho (\sigma )}}(1{\displaystyle \frac{\overline{n}(r)}{\rho (\sigma )}})^{n(r)\mu ;n(r),\nu }(\sigma )𝒢_{n(r)\mu ;n(r),\nu }(\sigma ).`$ (7.5c)
The ground state conformal weights are class functions and the factor $`\overline{n}(r)`$ in (7.5c) tells us that twist class $`\overline{n}(r)=0`$ (the integral affine subalgebra) does not contribute to $`\widehat{\mathrm{\Delta }}_0(\sigma )`$. When the CFT $`A(H)`$ is unitary, the ground state conformal weights $`\widehat{\mathrm{\Delta }}_0(\sigma )`$ must be real
$$\widehat{\mathrm{\Delta }}_0(\sigma )^{}=\widehat{\mathrm{\Delta }}_0(\sigma )$$
(7.6)
because the twisted sectors of these orbifolds are also unitary in this case. To see this explicitly, use the $`𝒢^{}`$ relation (A.4a), the $`^{}`$ relation (A.4c), the relation (A.3) for $``$ and the symmetry of $`𝒢`$ and $``$.
Using the duality transformations for $`𝒢`$ and $``$ in Eqs. (2.24d) and (2.33c), the ground state conformal weights (7.5b) can also be expressed in terms of the inverse inertia tensor $`L^{ab}`$ of the untwisted sector<sup>t</sup><sup>t</sup>tThe structure $`\frac{\overline{n}}{\rho }(1\frac{\overline{n}}{\rho })`$ was first seen in the ground state conformal weights of the free-field examples discussed by Dixon, Harvey, Vafa and Witten in Ref. References.
$$\widehat{\mathrm{\Delta }}_0(\sigma )=L^{ac}G_{cb}\underset{r}{}\frac{\overline{n}(r)}{2\rho (\sigma )}(1\frac{\overline{n}(r)}{\rho (\sigma )})𝒫(\overline{n}(r);\sigma )_a^b$$
(7.7a)
$$𝒫(\overline{n}(r);\sigma )_a^b=_\mu U^{}(\sigma )_a^{n(r)\mu }U(\sigma )_{n(r)\mu }^b,\omega (h_\sigma )_a^b=_rE_{n(r)}(\sigma )𝒫(\overline{n}(r);\sigma )_a^b$$
(7.7b)
$$𝒫(\overline{n}(r);\sigma )_a^c𝒫(\overline{n}(s);\sigma )_c^b=\delta _{n(r)}^{n(s)}𝒫(\overline{n}(r);\sigma )_a^b$$
(7.7c)
where $`𝒫(\overline{n}(r);\sigma )`$ is the projector onto the $`\overline{n}(r)`$ subspace of sector $`\sigma `$. Other properties of these projectors are collected in App. Appendix A..
To go further for $`HS_N`$, we need a more explicit notation for the semisimplicity of the Lie algebra $`g`$ in (7.1) and the degeneracy indices of the $`H`$-eigenvalue problem:
$$aa,I,n(r),\mu n(r),aj,a=1,\mathrm{},\text{dim}𝔤,I=0,\mathrm{},K1$$
(7.8a)
$$G_{ab}G_{aI,bJ}=k\eta _{ab}\delta _{IJ},f_{ab}^cf_{aI,bJ}^{cM}=f_{ab}^c\delta _{IJ}\delta _J^M$$
(7.8b)
$$J_{aI}(z)J_{bJ}(w)=\delta _{IJ}\{\frac{k\eta _{ab}}{(zw)^2}+\frac{f_{ab}^cJ_{cJ}(w)}{zw}\}+O((zw)^0)$$
(7.8c)
$$L^{ab}L^{aI,bJ},\rho _a^b\rho _{aI}^{bJ}=\rho _a^b\delta _I^J,\omega (h_\sigma )_a^b\omega (h_\sigma )_{aI}^{bJ}=\delta _a^b\omega (h_\sigma )_I^J$$
(7.8d)
$$U^{}(\sigma )_a^{n(r)\mu }U^{}(\sigma )_{aJ}^{n(r)bj}=\delta _a^bU^{}(\sigma )_J^{n(r)j},\omega (h_\sigma )_I^JU^{}(\sigma )_J^{n(r)j}=U^{}(\sigma )_I^{n(r)j}E_{n(r)}(\sigma )$$
(7.8e)
$$𝒫(\overline{n}(r);\sigma )_a^b𝒫(\overline{n}(r);\sigma )_{aI}^{bJ}=\delta _a^b𝒫(\overline{n}(r);\sigma )_I^J$$
(7.8f)
$$𝒫(\overline{n}(r);\sigma )_I^J=\underset{j}{}U^{}(\sigma )_I^{n(r)j}U(\sigma )_{n(r)j}^J$$
(7.8g)
$$𝒢_{n(r)aj;n(s)bl}(\sigma )k\eta _{ab},_{n(r)aj;n(s)bl}^{n(t)cm}(\sigma )f_{ab}^c,_{n(r)aj}^{n(s)bl}(\sigma )\rho _a^b$$
(7.8h)
$$T=L^{aI,bJ}:J_{aI}J_{bJ}:\begin{array}{c}\hfill \\ ^\sigma \hfill \end{array}\widehat{T}_\sigma =^{n(r)aj;n(s)bl}(\sigma ):\widehat{J}_{n(r)aj}\widehat{J}_{n(s)bl}:.$$
(7.8i)
Here the $`a`$ indices to the right of the arrows are the Lie algebra indices of the copies $`𝔤^I`$, with Killing metric $`\eta _{ab}`$, and $`\mu =(a,j)`$ labels the $`\overline{n}(r)`$ subspace. The reduced matrix $`U^{}(\sigma )_J^{n(r)j}`$ (which solves the reduced eigenvalue problem in (7.8e)) is also unitary, and the reduced matrix $`𝒫(\overline{n}(r);\sigma )_I^J`$ is also a projector onto the $`\overline{n}(r)`$ subspace.
Then the ground state conformal weight (7.7a) takes the form
$$\widehat{\mathrm{\Delta }}_0(\sigma )=kL^{aI,bJ}\eta _{ab}\underset{r}{}\frac{\overline{n}(r)}{2\rho (\sigma )}(1\frac{\overline{n}(r)}{\rho (\sigma )})𝒫(\overline{n}(r);\sigma )_{IJ}$$
(7.9)
for all $`HS_N`$, where $`P_{IJ}=P_I^K\delta _{KJ}`$. For the case of the cyclic permutation orbifolds $`A(_\lambda )/_\lambda `$, a more explicit form of the ground state conformal weights can be obtained
$$L^{aI,bJ}=L_{IJ}^{ab},𝒫(r;\sigma )_{IJ}=\frac{1}{\rho (\sigma )}e^{\frac{2\pi iN(\sigma )(JI)r}{\lambda }}\delta _{I,J\text{ mod }\frac{\lambda }{\rho (\sigma )}}$$
(7.10a)
$$\widehat{\mathrm{\Delta }}_0(\sigma )=\frac{\lambda k\eta _{ab}}{4\rho ^2(\sigma )}\{\frac{\rho ^2(\sigma )1}{3}L_0^{ab}\underset{r=1}{\overset{\rho (\sigma )1}{}}csc^2(\frac{\pi N(\sigma )r}{\rho (\sigma )})L_{\frac{\lambda }{\rho (\sigma )}r}^{ab}\},\sigma =1,\mathrm{},\lambda 1$$
(7.10b)
from (2.3) and (7.9), in agreement with the result given in Ref. References.
Returning to the general permutation orbifold, we find that the reduced matrices $`\omega `$ and $`U`$ satisfy
$$\underset{I}{}\omega (h_\sigma )_I^J=1$$
(7.11a)
$$(\underset{I}{}U^{}(\sigma )_I^{n(r)j})(1E_{n(r)}(\sigma ))=(\underset{I}{}U(\sigma )_{n(r)j}^I)(1E_{n(r)}(\sigma ))=0$$
(7.11b)
$$(\underset{I}{}U^{}(\sigma )_I^{n(r)j}),(\underset{I}{}U(\sigma )_{n(r)j}^I)\text{ and }\underset{I}{}𝒫(\overline{n}(r);\sigma )_I^J\delta _{n(r),0\text{ mod }\rho (\sigma )}$$
(7.11c)
for all $`HS_N`$. Here (7.11a) follows because $`\omega (h_\sigma )_{IJ}`$ is a permutation matrix, while the selection rules (7.11b) follow from (7.11a) and the reduced eigenvalue problem.
The relations (7.9) and (7.11c) give a simple form for the ground state conformal weights of the permutation orbifolds $`A(S_N)/S_N`$
$$L^{aI,bJ}=\lambda ^{ab}\delta ^{IJ}+l^{ab},\underset{I}{}𝒫(\overline{n}(r);\sigma )_I^I=\text{dim}[\overline{n}(r)]$$
(7.12a)
$$\widehat{\mathrm{\Delta }}_0(\sigma )=k\eta _{ab}\lambda ^{ab}\underset{r}{}\frac{\overline{n}(r)}{2\rho (\sigma )}(1\frac{\overline{n}(r)}{\rho (\sigma )})\text{dim}[\overline{n}(r)]$$
(7.12b)
where $`\text{dim}[\overline{n}(r)]`$ is the dimension of the $`\overline{n}(r)`$ subspace. This simplification depends on the form of the inverse inertia tensor in (7.12a), which describes the general $`S_N`$-invariant CFT. A more explicit form of these ground state conformal weights will be obtained in Sec. 8.
The relations (7) also give a simple form for the ground state conformal weights of the general permutation copy orbifold<sup>u</sup><sup>u</sup>uFor the case $`H=_\lambda `$(permutation), results equivalent to (u) were given in Refs. References and References.
$$\frac{\times _{I=0}^{K1}A_I}{H}\frac{A(H)}{H},HS_N(\text{permutation})$$
(7.13a)
$$L^{aI,bJ}=\lambda ^{ab}\delta ^{IJ}$$
(7.13b)
$$\widehat{\mathrm{\Delta }}_0(\sigma )=\frac{c}{2}\underset{r}{}\frac{\overline{n}(r)}{2\rho (\sigma )}(1\frac{\overline{n}(r)}{\rho (\sigma )})\text{dim}[\overline{n}(r)],c=2k\eta _{ab}\lambda ^{ab}$$
(7.13c)
where $`A_IA`$ are $`K`$ copies (permuted by $`HS_N`$) of any affine-Virasoro construction $`A`$ with central charge $`c`$. Included in this set of copy orbifolds is the general WZW permutation orbifold (see also Subsec. 6.4), whose twisted inverse inertia tensors are given by
$$_{\widehat{𝔤}(\sigma )}^{n(r)aj;n(s)bl}(\sigma )=\frac{x_𝔤}{2(x_𝔤+\stackrel{~}{h}_𝔤)}𝒢^{n(r)aj;n(s)bl}(\sigma ),𝒢^{n(r)aj;n(s)bl}(\sigma )\frac{\eta ^{ab}}{k},\sigma =0,\mathrm{},N_c1.$$
(7.14)
Here $`N_c`$ is the number of conjugacy classes of $`HS_N`$, $`\stackrel{~}{h}_𝔤`$ is the dual Coxeter number of each copy $`𝔤`$ and $`x_𝔤`$ is the invariant level of affine $`𝔤`$. The twisted tensor $`𝒢`$ with all indices up is the inverse of the twisted metric $`𝒢_{n(r)aj;n(s)bl}(\sigma )`$ in (7.8h).
For both $`A(S_N)/S_N`$ and the general permutation copy orbifold (u) we find an $`ab`$ symmetry of $``$
$$^{n(r)aj;n(s)bl}(\sigma )=^{n(r)bj;n(s)al}(\sigma )$$
(7.15a)
$$\widehat{T}_\sigma (z)=^{n(r)aj;n(r),bl}(\sigma ):\widehat{J}_{n(r)aj}(z)\widehat{J}_{n(r),bl}(z):_M+\frac{1}{z^2}\widehat{\mathrm{\Delta }}_0(\sigma )$$
(7.15b)
so that no linear terms in the currents appear in the $`M`$-ordered form of the stress tensors. For $`_\lambda `$-(permutation) copy orbifolds, many examples of this phenomenon were seen in Ref. References.
## 8 The Permutation Orbifolds $`A(S_N)/S_N`$
### 8.1 The permutation group $`S_N`$
As a large example beyond $`A(_\lambda )/_\lambda `$, we will work out the permutation orbifolds
$$\frac{A(S_N)}{S_N},S_N(\text{permutation})Aut(g)$$
(8.1)
in further detail, where $`A(S_N)`$ is any $`S_N`$(permutation)-invariant affine-Virasoro construction. The case of $`A(S_3)/S_3`$ was worked out previously in Ref. References. In the $`S_N`$-invariant CFT’s, we have $`N`$ copies $`𝔤^I`$ of an affine Lie algebra on simple $`𝔤`$ with currents $`J_{aI}`$
$$g=_{I=0}^{N1}𝔤^I,𝔤^I𝔤$$
(8.2a)
$$[J_{aI}(m),J_{bJ}(n)]=\delta _{IJ}\{if_{ab}^cJ_{cI}(m+n)+mk\eta _{ab}\delta _{m+n,0}\}$$
(8.2b)
$$J_{aI}^{}=\omega _I^JJ_{aJ},\omega S_N\text{(permutation)}$$
(8.2c)
$$a,b=1,\mathrm{},\text{dim}𝔤,I,J=0,\mathrm{},N1,m,n$$
(8.2d)
and the action of $`S_N`$(permutation) on the currents is given in (8.2c).
In the cyclic notation for the elements of $`S_N`$, the action of a cycle of length $`l`$
$$(I_{\widehat{j}=0}\mathrm{}I_{\widehat{j}=l1}):I_{\widehat{j}}=0,\mathrm{},N1,\widehat{j}=0,\mathrm{},l1,I_{\widehat{j}}I_{\widehat{i}}\text{ when }\widehat{j}\widehat{i}$$
(8.3)
is a cyclic permutation of the copies of $`𝔤`$ according to
$$𝔤^{I_{\widehat{j}}}𝔤^{I_{\widehat{j}+1}},\widehat{j}=0,\mathrm{},l2;𝔤^{I_{l1}}𝔤^{I_0}.$$
(8.4)
All elements of $`S_N`$ are products of disjoint cycles, and the general element can be written as a product of elements of disjoint cyclic groups $`_{\sigma _j}`$
$$\underset{j=0}{\overset{n1}{}}(I_{\widehat{j}=0}^j\mathrm{}I_{\widehat{j}=\sigma _j1}^j),I_{\widehat{j}}^jI_{\widehat{i}}^i\text{ when }ji$$
(8.5)
where $`\sigma _j`$ is the length of the $`j`$th disjoint cycle, $`\widehat{j}`$ counts within a cycle and $`n`$ (the number of disjoint cycles) is some positive integer.
A property of each element of $`S_N`$ is its cycle type $`\stackrel{}{\sigma }`$, which is the collection of lengths of cycles written in order of decreasing length
$$\stackrel{}{\sigma }\{\sigma _0,\mathrm{},\sigma _{n1}\},\sigma _{j+1}\sigma _j,\underset{j=0}{\overset{n(\stackrel{}{\sigma })1}{}}\sigma _j=N.$$
(8.6)
The cycle types $`\stackrel{}{\sigma }`$ are the $`\sigma _{j+1}\sigma _j`$ partitions of $`N`$ for any $`n=n(\stackrel{}{\sigma })`$, and $`\stackrel{}{\sigma }`$ also labels the conjugacy classes of $`S_N`$. Conversely, the elements of conjugacy class $`\stackrel{}{\sigma }`$ are obtained by relaxing the restriction on the ordering of the lengths in $`\stackrel{}{\sigma }`$. As an example, the table below
| Conj.Class $`\stackrel{}{\sigma }`$ | Elements |
| --- | --- |
| $`\{\sigma _0,\sigma _1,\sigma _2\}`$={1,1,1} | (0)(1)(2) |
| $`\{\sigma _0,\sigma _1\}`$={2,1} | (01)(2),(02)(1),(12)(0) |
| $`\{\sigma _0\}`$={3} | (012),(021) |
(8.7)
shows the conjugacy classes of $`S_3`$.
Modeling our choice on the cycle types, we choose one representative from each conjugacy class $`\stackrel{}{\sigma }`$ by determining $`\{I_{\widehat{j}}^j\}`$ as follows
$$j,\widehat{j}\{I_{\widehat{j}}^j\}:I_{\widehat{j}}^j\underset{k=0}{\overset{j1}{}}\sigma _k+\widehat{j},\sigma _{j+1}\sigma _j,\underset{j=0}{\overset{n(\stackrel{}{\sigma })1}{}}\sigma _j=N$$
(8.8a)
$$j=0,\mathrm{},n(\stackrel{}{\sigma })1,\widehat{j}=0,\mathrm{},\sigma _j1$$
(8.8b)
given $`\stackrel{}{\sigma }`$ and the ranges of $`j,\widehat{j}`$ in (8.8b). Here $`j`$ labels the disjoint cycles and $`\widehat{j}`$ labels the integers in each disjoint cycle. For example, the chosen representatives for the elements of $`S_3`$ and $`S_4`$ are
| $`S_3`$: Conj. Class $`\stackrel{}{\sigma }`$ | Representative |
| --- | --- |
| {1,1,1} | (0)(1)(2)=$`(I_0^0)(I_0^1)(I_0^2)`$ |
| {2,1} | (01)(2)=$`(I_0^0I_1^0)(I_0^1)`$ |
| {3} | (012)=$`(I_0^0I_1^0I_2^0)`$ |
| $`S_4`$: Conj. Class $`\stackrel{}{\sigma }`$ | Representative |
| --- | --- |
| {1,1,1,1} | (0)(1)(2)(3) |
| {2,1,1} | (01)(2)(3) |
| {2,2} | (01)(23) |
| {3,1} | (012)(3) |
| {4} | (0123). |
(8.9)
We will refer to the representative of conjugacy class $`\stackrel{}{\sigma }`$ as $`h_\stackrel{}{\sigma }S_N`$.
### 8.2 Automorphisms and the twisted currents
To construct the action $`\omega `$ of the automorphism $`h_\stackrel{}{\sigma }`$ on the currents, we begin by relabeling the copies of $`𝔤`$
$$𝔤^I𝔤^{j\widehat{j}},a,Ia,j\widehat{j},J_{aI}J_{aj\widehat{j}}$$
(8.10a)
$$\underset{k=0}{\overset{j1}{}}\sigma _k+\widehat{j}=I,j\{0,\mathrm{},n(\stackrel{}{\sigma })1\},\widehat{j}\{0,\mathrm{},\sigma _j1\}$$
(8.10b)
where $`I`$ is the semisimplicity index in (8.1), $`j`$ and $`\widehat{j}`$ are the unique solutions to (8.10b) given $`I`$, and $`\stackrel{}{\sigma }`$ solves Eq. (8.1). This relabeling (which is the inverse of (8.1)) corresponds to the pictorial representation in Fig. 2 of the action of $`h_\stackrel{}{\sigma }`$ on the copies
where $`\widehat{j}`$ labels the integers which are cyclically permuted $`\widehat{j}\widehat{j}+1`$ inside each box (disjoint cycle) $`j`$. It follows that the action $`\omega (\stackrel{}{\sigma })\omega (h_\stackrel{}{\sigma })`$ on the currents is
$$\omega \omega _{aI}^{bJ}(\stackrel{}{\sigma })=\omega _{aj\widehat{j}}^{bl\widehat{l}}(\stackrel{}{\sigma })=\delta _a^b\omega _{j\widehat{j}}^{l\widehat{l}}(\stackrel{}{\sigma })$$
(8.11a)
$$J_{aj\widehat{j}}^{}=\underset{l,\widehat{l}}{}\omega _{j\widehat{j}}^{l\widehat{l}}(\stackrel{}{\sigma })J_{al\widehat{l}},\omega _{j\widehat{j}}^{l\widehat{l}}(\stackrel{}{\sigma })=\delta _{jl}\delta _{\widehat{j}+1,\widehat{l}\text{ mod }\sigma _j}$$
(8.11b)
where the block-diagonal matrix $`\omega _{j\widehat{j}}^{l\widehat{l}}(\stackrel{}{\sigma })`$ is orthogonal. In what follows, all quantities have the spectral index periodicity $`\widehat{j}\widehat{j}+\sigma _j`$.
The next step is to solve the $`S_N`$-eigenvalue problem
$$n(r),\mu \widehat{j},aj,U^{}U^{}(\stackrel{}{\sigma })_{\widehat{j}aj}^{\widehat{l}bl}=\delta _a^bU^{}(\stackrel{}{\sigma })_{\widehat{j}j}^{\widehat{l}l}$$
(8.12a)
$$\underset{\widehat{l},l}{}\omega _{\widehat{m}m}^{\widehat{l}l}(\stackrel{}{\sigma })U^{}(\stackrel{}{\sigma })_{\widehat{l}l}^{\widehat{j}j}=U^{}(\stackrel{}{\sigma })_{\widehat{m}m}^{\widehat{j}j}E_{\widehat{j}}^j(\stackrel{}{\sigma })$$
(8.12b)
for the block-diagonal and unitary matrix $`U^{}(\stackrel{}{\sigma })`$ and the eigenvalues $`E(\stackrel{}{\sigma })`$:
$$U^{}(\stackrel{}{\sigma })_{\widehat{j}j}^{\widehat{l}l}=\frac{\delta _{jl}}{\sqrt{\sigma _j}}e^{\frac{2\pi i\widehat{j}\widehat{l}}{\sigma _j}},U(\stackrel{}{\sigma })_{\widehat{j}j}^{\widehat{l}l}=\frac{\delta _{jl}}{\sqrt{\sigma _j}}e^{\frac{2\pi i\widehat{j}\widehat{l}}{\sigma _j}},E_{\widehat{j}}^j(\stackrel{}{\sigma })=e^{\frac{2\pi i\widehat{j}}{\sigma _j}}$$
(8.13a)
$$\underset{\widehat{m},m}{}U(\stackrel{}{\sigma })_{\widehat{j}j}^{\widehat{m}m}U^{}(\stackrel{}{\sigma })_{\widehat{m}m}^{\widehat{l}l}=\delta _{jl}\delta _{\widehat{j}+\widehat{l},0\text{ mod }\sigma _j},\underset{\widehat{j},j}{}U^{}(\stackrel{}{\sigma })_{\widehat{j}j}^{\widehat{l}l}=\sqrt{\sigma _l}\delta _{\widehat{l},0\text{ mod }\sigma _l}.$$
(8.13b)
The second identity in Eq. (8.13b) is the form taken by Eq. (7.11c) in this case.
With the choice $`\chi _{\widehat{j}aj}\left(\stackrel{}{\sigma }\right)=\sqrt{\sigma _j}`$ we find the twisted current system of sector $`\stackrel{}{\sigma }`$ of the orbifold $`A(S_N)/S_N`$
$$\widehat{J}_{n(r)\mu }\widehat{J}_{aj}^{(\widehat{j})},𝒢𝒢_{\widehat{j}aj;\widehat{l}bl}(\stackrel{}{\sigma })=\delta _{jl}\sigma _jk\eta _{ab}\delta _{\widehat{j}+\widehat{l},0\text{ mod }\sigma _j}$$
(8.14a)
$$_{\widehat{j}aj;\widehat{l}bl}^{\widehat{m}cm}(\stackrel{}{\sigma })=f_{ab}^c\delta _{jl}\delta _l^m\delta _{\widehat{j}+\widehat{l},\widehat{m}\text{ mod }\sigma _j}$$
(8.14b)
$$\widehat{J}_{aj}^{(\widehat{j})}(z)\widehat{J}_{bl}^{(\widehat{l})}(w)=\delta _{jl}\{\frac{(\sigma _jk)\eta _{ab}\delta _{\widehat{j}+\widehat{l},0\text{ mod }\sigma _j}}{(zw)^2}+\frac{if_{ab}^c\widehat{J}_{cj}^{(\widehat{j}+\widehat{l})}(w)}{zw}\}+O((zw)^0)$$
(8.14c)
$$\widehat{J}_{aj}^{(\widehat{j})}(ze^{2\pi i})=e^{\frac{2\pi i\widehat{j}}{\sigma _j}}\widehat{J}_{aj}^{(\widehat{j})}(z),\widehat{J}_{aj}^{(\widehat{j}\pm \sigma _j)}(z)=\widehat{J}_{aj}^{(\widehat{j})}(z)$$
(8.14d)
$$\mathrm{\#}\{\widehat{J}\}=\text{dim}𝔤(\underset{j=0}{\overset{n(\stackrel{}{\sigma })1}{}}\sigma _j)=N\text{dim}𝔤=\text{dim}g=\mathrm{\#}\{J\}$$
(8.14e)
$$a,b=1,\mathrm{},\text{dim}𝔤,j,l=0,\mathrm{},n(\stackrel{}{\sigma })1,\overline{\widehat{j}}=0,\mathrm{},\sigma _j1,\overline{\widehat{l}}=0,\mathrm{},\sigma _l1$$
(8.14f)
where $`\overline{\widehat{j}}=\widehat{j}\sigma _j\widehat{j}/\sigma _j`$ evaluates $`\widehat{j}`$ in its fundamental range. According to (8.14d), the fraction
$$\frac{n(r)}{\rho (\sigma )}=\frac{\widehat{j}}{\sigma _j}$$
(8.15)
controls the monodromies of the twisted current $`\widehat{J}_{aj}^{(\widehat{j})}`$.
The modes of the twisted currents follow from the monodromies
$$\widehat{J}_{aj}^{(\widehat{j})}(z)=\underset{m}{}\widehat{J}_{aj}^{(\widehat{j})}(m+\frac{\widehat{j}}{\sigma _j})z^{(m+\frac{\widehat{j}}{\sigma _j})1}$$
(8.16a)
$$\widehat{J}_{aj}^{(\widehat{j}+\sigma _j)}(m1+\frac{\widehat{j}+\sigma _j}{\sigma _j})=\widehat{J}_{aj}^{(\widehat{j})}(m+\frac{\widehat{j}}{\sigma _j})$$
(8.16b)
$$\widehat{J}_{aj}^{(\widehat{j})}(m+\frac{\widehat{j}}{\sigma _j})|0_\stackrel{}{\sigma }=0\text{ when }(m+\frac{\widehat{j}}{\sigma _j})0$$
(8.16c)
where the periodicity and ground state condition are special cases of Eqs. (3.2a) and (7.2) respectively.
This leads to the twisted current algebra $`\widehat{𝔤}(\stackrel{}{\sigma })`$ of sector $`\stackrel{}{\sigma }`$
$$\widehat{𝔤}(\stackrel{}{\sigma })=\widehat{𝔤}(S_N(\text{permutation})Aut(g);\stackrel{}{\sigma })=_{j=0}^{n(\stackrel{}{\sigma })1}\widehat{𝔤}_{\sigma _j}$$
(8.17a)
$$[\widehat{J}_{aj}^{(\widehat{j})}(m+\frac{\widehat{j}}{\sigma _j}),\widehat{J}_{bl}^{(\widehat{l})}(n+\frac{\widehat{l}}{\sigma _l})]=\delta _{jl}\{if_{ab}^c\widehat{J}_{cj}^{(\widehat{j}+\widehat{l})}(m+n+\frac{\widehat{j}+\widehat{l}}{\sigma _j})+(\sigma _jk)\eta _{ab}(m+\frac{\widehat{j}}{\sigma _j})\delta _{m+n+\frac{\widehat{j}+\widehat{l}}{\sigma _j},0}\}$$
(8.17b)
$$a,b=1,\mathrm{},\text{dim}𝔤,j,l=0,\mathrm{},n(\stackrel{}{\sigma })1,\overline{\widehat{j}}=0,\mathrm{},\sigma _j1,\overline{\widehat{l}}=0,\mathrm{},\sigma _l1$$
(8.17c)
which is a special case of the general twisted current algebra in (3.1). This algebra consists of a product of $`n(\stackrel{}{\sigma })`$ commuting sets of orbifold affine algebras$`^{\text{References}}`$ $`\widehat{𝔤}_{\sigma _j}`$ of order $`\sigma _j`$, where the factor $`\widehat{𝔤}_{\sigma _j}`$ has orbifold affine level $`\widehat{k}_j=\sigma _jk`$ and $`k`$ is the level of the untwisted algebra. According to the identification (8.15), the quantity $`\overline{\widehat{j}}=\widehat{j}\sigma _j\widehat{j}/\sigma _j`$ is the effective twist class of $`\widehat{J}_{aj}^{(\widehat{j})}`$, relative to the order $`\sigma _j`$ of $`\widehat{𝔤}_{\sigma _j}`$.
For this case, the orbifold adjoint operation (4.4a) takes the form
$$_{\widehat{j}aj}^{\widehat{l}bl}(\stackrel{}{\sigma })=\delta _{jl}\rho _a^b\delta _{\widehat{j}+\widehat{l},0\text{ mod }\sigma _j}\widehat{J}_{aj}^{(\widehat{j})}(m+\frac{\widehat{j}}{\sigma _j})^{}=\rho _a^b\widehat{J}_{bj}^{(\widehat{j})}(m\frac{\widehat{j}}{\sigma _j}).$$
(8.18)
This is the standard adjoint operation of orbifold affine algebra,$`^{\text{References},\text{References}}`$ which guarantees unitarity of the twisted affine Hilbert space given unitarity of the untwisted affine Hilbert space. As noted above, unitarity of the permutation orbifolds $`A(S_N)/S_N`$ then follows from unitarity of the CFT $`A(S_N)`$.
### 8.3 The stress tensor of sector $`\stackrel{}{\sigma }`$
The stress tensors of the $`S_N`$-invariant CFT’s $`A(S_N)`$ are
$$T=L^{aJ,bL}\underset{J,L=0}{\overset{N1}{}}:J_{aJ}J_{bL}:,L^{aJ,bL}=\lambda ^{ab}\delta ^{JL}+l^{ab}$$
(8.19a)
$$\lambda ^{ab}=\lambda ^{ba},l^{ab}=l^{ba},c=2Nk\eta _{ab}(\lambda ^{ab}+l^{ab})$$
(8.19b)
where the inverse inertia tensors $`\lambda ^{ab}`$ and $`l^{ab}`$ satisfy the reduced Virasoro master equation
$$\lambda ^{ab}=2N\lambda ^{ac}\eta _{cd}\lambda ^{db}\lambda ^{cd}(\lambda ^{ef}+2l^{ef})f_{ce}^af_{df}^b(\lambda ^{cd}+l^{cd})f_{ce}^ff_{df}^{(a}\lambda ^{b)e}$$
(8.20a)
$$l^{ab}=4kl^{ac}\eta _{cd}(\lambda ^{db}+l^{db})+2k(N2)l^{ac}\eta _{cd}\lambda ^{db}l^{cd}l^{ef}f_{ce}^af_{df}^b(\lambda ^{cd}+l^{cd})f_{ce}^ff_{df}^{(a}l^{b)e}.$$
(8.20b)
This is a consistent set of $`\text{dim}𝔤(\text{dim}𝔤+1)`$ quadratic equations for the same number of unknowns and so the generically-expected number of inequivalent solutions at each level is
$$M(𝔤,N)=2^{(\text{dim}𝔤)^2}.$$
(8.21)
Remarkably, $`M(𝔤,N)`$ is independent of $`N`$.
Then the stress tensor of sector $`\stackrel{}{\sigma }`$ of $`A(S_N)/S_N`$
$$\widehat{T}_\stackrel{}{\sigma }=\underset{j,l=0}{\overset{n(\stackrel{}{\sigma })1}{}}\underset{\widehat{j}=0}{\overset{\text{ gcd }(\sigma _j,\sigma _l)1}{}}^{\widehat{j}aj;\widehat{j},bl}(\stackrel{}{\sigma }):\widehat{J}_{aj}^{(\frac{\widehat{j}\sigma _j}{\text{ gcd }(\sigma _j,\sigma _l)})}\widehat{J}_{bl}^{(\frac{\widehat{j}\sigma _l}{\text{ gcd }(\sigma _j,\sigma _l)})}:$$
(8.22a)
$$^{\widehat{j}aj;\widehat{j},bl}(\stackrel{}{\sigma })=\lambda ^{ab}\frac{\delta ^{jl}}{\sigma _j}+l^{ab}\delta _{\widehat{j},0}$$
(8.22b)
follows as a special case of the general Virasoro construction (2.32). Here $`z=`$gcd($`x,y`$) is the greatest integer such that $`x/z`$ and $`y/z`$ are also integers. Combining (8.22a) and (8.22b), we find a more transparent form of these stress tensors
$$\widehat{T}_\stackrel{}{\sigma }=\lambda ^{ab}\underset{j=0}{\overset{n(\stackrel{}{\sigma })1}{}}\frac{1}{\sigma _j}\underset{\widehat{j}=0}{\overset{\sigma _j1}{}}:\widehat{J}_{aj}^{(\widehat{j})}\widehat{J}_{bj}^{(\widehat{j})}:+l^{ab}\underset{j,l=0}{\overset{n(\stackrel{}{\sigma })1}{}}:\widehat{J}_{aj}^{(0)}\widehat{J}_{bl}^{(0)}:$$
(8.23a)
$$L_\stackrel{}{\sigma }(m)=\underset{p}{}\{\lambda ^{ab}\underset{j=0}{\overset{n(\stackrel{}{\sigma })1}{}}\frac{1}{\sigma _j}\underset{\widehat{j}=0}{\overset{\sigma _j1}{}}:\widehat{J}_{aj}^{(\widehat{j})}(p+\frac{\widehat{j}}{\sigma _j})\widehat{J}_{bj}^{(\widehat{j})}(mp\frac{\widehat{j}}{\sigma _j}):$$
$$+l^{ab}\underset{j,l=0}{\overset{n(\stackrel{}{\sigma })1}{}}:\widehat{J}_{aj}^{(0)}(p)\widehat{J}_{bl}^{(0)}(mp):\}$$
(8.23b)
$$\widehat{c}(\stackrel{}{\sigma })=c=2Nk\eta _{ab}(\lambda ^{ab}+l^{ab})$$
(8.23c)
$$\widehat{\mathrm{\Delta }}_0(\stackrel{}{\sigma })=\frac{k\eta _{ab}\lambda ^{ab}}{12}\underset{j=0}{\overset{n(\stackrel{}{\sigma })1}{}}(\sigma _j\frac{1}{\sigma _j})=\frac{k\eta _{ab}\lambda ^{ab}}{12}(N\underset{j=0}{\overset{n(\stackrel{}{\sigma })1}{}}\frac{1}{\sigma _j}).$$
(8.23d)
The ground state conformal weight of sector $`\stackrel{}{\sigma }`$ is given in (8.23d).
The result (8.23d) for the ground state conformal weights of $`A(S_N)/S_N`$ is in agreement with our earlier result (7.12b). To see this, start with (7.12b) and follow the steps
$`\widehat{\mathrm{\Delta }}_0(\sigma )`$ $`=`$ $`k\eta _{ab}\lambda ^{ab}{\displaystyle \underset{r}{}}{\displaystyle \frac{\overline{n}(r)}{2\rho (\sigma )}}(1{\displaystyle \frac{\overline{n}(r)}{\rho (\sigma )}})dim[\overline{n}(r)]=k\eta _{ab}\lambda ^{ab}{\displaystyle \underset{r,\mu (r)}{}}{\displaystyle \frac{\overline{n}(r)}{2\rho (\sigma )}}(1{\displaystyle \frac{\overline{n}(r)}{\rho (\sigma )}})`$ (8.24)
$`=`$ $`k\eta _{ab}\lambda ^{ab}{\displaystyle \underset{j=0}{\overset{n(\stackrel{}{\sigma })1}{}}}{\displaystyle \underset{\widehat{j}=0}{\overset{\sigma _j1}{}}}{\displaystyle \frac{\widehat{j}}{2\sigma _j}}(1{\displaystyle \frac{\widehat{j}}{\sigma _j}})={\displaystyle \frac{k\eta _{ab}\lambda ^{ab}}{12}}{\displaystyle \underset{j=0}{\overset{n(\stackrel{}{\sigma })1}{}}}(\sigma _j{\displaystyle \frac{1}{\sigma _j}}).`$
Here we have used Eq. (8.15) in the fundamental range (where $`\overline{\widehat{j}}=\widehat{j}`$) and the identities
$$\underset{\mu (r)}{}=dim[\overline{n}(r)],\underset{r,\mu (r)}{}=\underset{j=0}{\overset{n(\stackrel{}{\sigma })1}{}}\underset{\widehat{j}=0}{\overset{\sigma _j1}{}}$$
(8.25)
to convert from the notation of Sec. 7 to the present notation. As examples, the cases $`A(S_3)/S_3`$ and $`A(S_4)/S_4`$ are further discussed in App. Appendix E..
## 9 The Inner-Automorphic Orbifolds $`A(H(d))/H(d)`$
Our next large example is the set of all<sup>v</sup><sup>v</sup>vThe special case of inner-automorphic WZW orbifolds has been considered in Refs. References, References, References, References, References and References, and inner-automorphic coset orbifolds were also considered in Ref. References. inner-automorphic orbifolds $`A(H(d))/H(d)`$, where $`A(H(d))`$ is any inner automorphic invariant CFT. The story of these orbifolds is particularly interesting, not least because of their overlap with the orbifolds $`A(\text{Lie}h(H))/H`$. Indeed, we will argue that this overlap contains almost all the inner-automorphic orbifolds which can be equivalently described by stress-tensor spectral flow$`^{\text{References},\text{References},\text{References},\text{References}}`$ whereas the generic inner-automorphic orbifold apparently can not be described in this way.
### 9.1 Inner Automorphisms of simple $`g`$
We begin in the Cartan-Weyl basis of the general affine algebra on simple $`g`$
$$G_{ab}=k\eta _{ab},\eta _{ab}=\rho _a^b=\left(\begin{array}{cc}\delta _{AB}& 0\\ 0& \delta _{\alpha +\beta ,0}\end{array}\right)$$
(9.1a)
$$H_A(z)H_B(w)=\frac{k\delta _{AB}}{(zw)^2}+O((zw)^0),H_A(z)E_\alpha (w)=\frac{\alpha _AE_\alpha (w)}{zw}+O((zw)^0)$$
(9.1b)
$$E_\alpha (z)E_\beta (w)=\{\begin{array}{ccc}\frac{N_\gamma (\alpha ,\beta )E_\gamma (w)}{zw}+O((zw)^0)\hfill & & \text{if}\alpha +\beta =\gamma \hfill \\ \frac{k}{(zw)^2}+\frac{\alpha H(w)}{zw}+O((zw)^0)\hfill & & \text{if}\alpha +\beta =0\hfill \\ O((zw)^0)\hfill & & \text{otherwise}\hfill \end{array}$$
(9.1c)
$$a=(A,\alpha ),A,B=1,\mathrm{},\text{rank}g,\alpha ,\beta ,\gamma \mathrm{\Delta }(g).$$
(9.1d)
In this basis, the action of the general group $`H(d)`$ of inner automorphisms has the form
$$H(d)\text{Lie}GAut(g)$$
(9.2a)
$$H_A(z)^{}=H_A(z),E_\alpha (z)^{}=e^{2\pi i\sigma \alpha d}E_\alpha (z)$$
(9.2b)
$$\sigma =0,\mathrm{},\rho (1)1,N_c=\rho (1)$$
(9.2c)
where $`\text{Lie}G`$ is (the action in the adjoint of) the Lie group whose algebra is $`\text{Lie}g`$ and $`\rho (1)=\rho (\sigma =1)`$ is the order of the $`\sigma =1`$ element of $`H(d)`$. The general form<sup>w</sup><sup>w</sup>wFurther discussion of the allowed vectors $`d`$ is given in Ref. References. More general forms of the vector $`d`$ describe automorphism groups of infinite order. Our results below are well defined in this case, leading at least formally to orbifolds with an infinite number of sectors $`\sigma =0,1,\mathrm{}\mathrm{}`$. of the vector $`d`$
$$d=\frac{2}{N}\underset{i=1}{\overset{\text{rank}g}{}}\frac{q_i\lambda _i}{\alpha _i^2},N^+,q_i^+$$
(9.3a)
$$\alpha _i\lambda _j=\frac{\alpha _i^2}{2}\delta _{ij},\text{gcd}(N,q_1,\mathrm{},q_{\text{rank}g})=1$$
(9.3b)
gives $`\rho (1)=N`$. Here $`\{\alpha _i\}`$ and $`\{\lambda _i\}`$ are the simple roots and weights of $`g`$, and $`z=\text{gcd}(\{x_n\})`$ is the greatest integer such that $`(x_n/z)^+`$ for all n. Note that $`d`$ is inversely proportional to the length of the highest root. For $`\sigma 2`$ one finds that $`\rho (\sigma )=\rho (1)/\text{gcd}(\sigma ,\rho (1))`$, which holds for all cyclic groups. As an example, the grade automorphism is discussed in App. Appendix F..
We also need the $`H`$-invariance condition (2.7a) for the inverse inertia tensor in this case:
$$L^{A\alpha }(1e^{2\pi i\sigma \alpha d})=0,L^{\alpha \beta }(1e^{2\pi i\sigma (\alpha +\beta )d})=0,\text{ no restriction on }L^{AB}$$
(9.4a)
$$\sigma =0,\mathrm{},\rho (1)1.$$
(9.4b)
This gives the stress tensors of the inner automorphic invariant CFT’s $`A(H(d))`$
$$T=L^{ab}:J_aJ_b:,L^{ab}=L^{ba}$$
(9.5a)
$$L^{A\alpha }=0\text{ unless }\alpha d,L^{\alpha \beta }=0\text{ unless }(\alpha +\beta )d$$
(9.5b)
where (9.5b) is the solution of (9.4a).
As a simple example, the affine-Sugawara construction$`^{\text{References},\text{References}\text{References},\text{References}}`$ on $`g`$
$$T_g=L_g^{ab}:J_aJ_b:,L_g^{AB}=\frac{\delta ^{AB}}{2k+Q_\psi },L_g^{\alpha \beta }=\frac{\delta _{\alpha +\beta ,0}}{2k+Q_\psi },\underset{\alpha }{}\alpha _A\alpha _B=Q_\psi \delta _{AB}$$
(9.6)
describes an $`H(d)`$-invariant CFT for any $`H(d)`$. When we have in mind a particular $`H(d)`$, we say that the affine-Sugawara construction describes the $`H(d)`$-invariant CFT $`A_g(H(d))`$.
### 9.2 Inner-automorphically twisted currents
The action of the automorphism group $`H(d)`$ in (9.1d) is already diagonal, so the eigenvalues, spectral indices $`n(r)`$ and twist classes $`\overline{n}(r)`$ read
$$E_{n(r)}=e^{\frac{2\pi in(r)}{\rho (\sigma )}}=e^{2\pi i\sigma \alpha d},\sigma =0,\mathrm{},\rho (1)1$$
(9.7a)
$$n(r)n_\alpha =\rho (\sigma )\sigma \alpha d,n_A=0$$
(9.7b)
$$\overline{n}_\alpha =\rho (\sigma )(\sigma \alpha d+\sigma \alpha d),\overline{n}_A=0$$
(9.7c)
where (9.7c) is obtained from (2.16). By the same token, we may choose
$$U(\sigma )=U^{}(\sigma )=1,\chi \left(\sigma \right)=1𝒥_a(\sigma )=J_a,a=1,\mathrm{},\text{dim}g$$
(9.8)
so that the eigencurrents $`𝒥`$ are the untwisted currents $`J`$. It follows that the twisted tensors of the general inner-automorphic orbifold $`A(H(d))/H(d)`$ are identical to the tensors of the untwisted sector
$$_{ab}^c(\sigma )=f_{ab}^c,𝒢_{ab}(\sigma )=G_{ab},^{ab}(\sigma )=L^{ab},_a^b(\sigma )=\rho _a^b.$$
(9.9)
Then, the $``$-selection rule (2.39a) is nothing but the $`H`$-invariance of $`L`$ in (9.1), and the $`𝒢`$\- and $``$-selection rules in (2.28b) and (2.29d) are automatically satisfied by the metric and structure constants in the Cartan-Weyl basis.
Then Eq. (2.4) gives the inner-automorphically twisted current system of sector $`\sigma `$
$$\widehat{H}_A(z)\widehat{H}_B(w)=\frac{k\delta _{AB}}{(zw)^2}+O((zw)^0),\widehat{H}_A(z)\widehat{E}_\alpha (w)=\frac{\alpha _A\widehat{E}_\alpha (w)}{zw}+O((zw)^0)$$
(9.10a)
$$\widehat{E}_\alpha (z)\widehat{E}_\beta (w)=\{\begin{array}{ccc}\frac{N_\gamma (\alpha ,\beta )\widehat{E}_\gamma (w)}{zw}+O((zw)^0)\hfill & & \text{if}\alpha +\beta =\gamma \hfill \\ \frac{k}{(zw)^2}+\frac{\alpha \widehat{H}(w)}{zw}+O((zw)^0)\hfill & & \text{if}\alpha +\beta =0\hfill \\ O((zw)^0)\hfill & & \text{otherwise}\hfill \end{array}$$
(9.10b)
$$\widehat{E}_\alpha (ze^{2\pi i})=e^{2\pi i\sigma \alpha d}\widehat{E}_\alpha (z),\widehat{H}_A(ze^{2\pi i})=\widehat{H}_A(z)$$
(9.10c)
$$A,B=1,\mathrm{},\text{rank}g,\alpha ,\beta ,\gamma \mathrm{\Delta }(g)$$
(9.10d)
of each orbifold $`A(H(d))/H(d)`$. Here we have suppressed our usual labeling by the spectral indices of the twisted currents, but the monodromies (9.10c) are recorded in the modes
$$\widehat{H}_A(z)=\underset{m}{}\widehat{H}_A(m)z^{m1},\widehat{E}_\alpha (z)=\underset{m}{}\widehat{E}_\alpha (m\sigma \alpha d)z^{(m\sigma \alpha d)1}.$$
(9.11)
These expansions lead to the twisted current algebra $`\widehat{𝔤}(\sigma )=\widehat{𝔤}(H(d)Aut(g);\sigma )`$
$$[\widehat{H}_A(m),\widehat{H}_B(n)]=km\delta _{AB}\delta _{m+n,0},[\widehat{H}_A(m),\widehat{E}_\alpha (n\sigma \alpha d)]=\alpha _A\widehat{E}_\alpha (m+n\sigma \alpha d)$$
(9.12a)
$$[\widehat{E}_\alpha (m\sigma \alpha d),\widehat{E}_\beta (n\sigma \beta d)]=\{\begin{array}{ccc}N_\gamma (\alpha ,\beta )\widehat{E}_\gamma (m+n\sigma \gamma d)\hfill & & \text{if}\alpha +\beta =\gamma \hfill \\ \alpha \widehat{H}(m+n)+k(m\sigma \alpha d)\delta _{m+n,0}\hfill & & \text{if}\alpha +\beta =0\hfill \\ 0\hfill & & \text{otherwise}\hfill \end{array}$$
(9.12b)
$$A,B=1,\mathrm{},\text{rank}g,\alpha ,\beta ,\gamma \mathrm{\Delta }(g),\sigma =0,\mathrm{},\rho (1)1$$
(9.12c)
which is a special case of the general twisted current algebra (3.1). As expected, this algebra is a sector-dependent set of inner-automorphically twisted$`^{\text{References},\text{References},\text{References},\text{References}}`$ affine Lie algebras, and we note that, in this case, the integral affine subalgebra $`\widehat{𝔤}^{(0)}(\sigma )`$ of each sector is at least the affine Cartan subalgebra.
The adjoint of the twisted currents
$$\widehat{H}_A(m)^{}=\widehat{H}_A(m),\widehat{E}_\alpha (m\sigma \alpha d)^{}=\widehat{E}_\alpha (m+\sigma \alpha d)$$
(9.13)
follows from the orbifold adjoint operation (4.4a), using $`(\sigma )=\rho `$ in this case, and we know from the orbifold induction procedure for inner-automorphic twists$`^{\text{References},\text{References},\text{References},\text{References}}`$ that this adjoint guarantees unitarity of the orbifolds $`A(H(d))/H(d)`$ when the CFT $`A(H(d))`$ is unitary.
In this case, it is convenient for our discussion below to choose $`M^{}`$ ordering (mode ordering with respect to $`m`$) defined in App. Appendix G.. This gives the exact operator products
$$\widehat{H}_A(z)\widehat{H}_B(w)=\frac{k\delta _{AB}}{(zw)^2}+:\widehat{H}_A(z)\widehat{H}_B(w):_M^{}$$
(9.14a)
$$\widehat{H}_A(z)\widehat{E}_\alpha (w)=\frac{\alpha _A\widehat{E}_\alpha (w)}{zw}+:\widehat{H}_A(z)\widehat{E}_\alpha (w):_M^{}$$
(9.14b)
$$\widehat{E}_\alpha (z)\widehat{E}_\beta (w)=:\widehat{E}_\alpha (z)\widehat{E}_\beta (w):_M^{}+\{\begin{array}{ccc}(\frac{z}{w})^{\sigma \alpha d}\frac{N_\gamma (\alpha ,\beta )\widehat{E}_\gamma (w)}{zw}\hfill & & \text{if}\alpha +\beta =\gamma \hfill \\ (\frac{z}{w})^{\sigma \alpha d}[\frac{k}{(zw)^2}\frac{k\sigma \alpha d}{w(zw)}+\frac{\alpha \widehat{H}(w)}{zw}]\hfill & & \text{if}\alpha +\beta =0\hfill \\ 0\hfill & & \text{otherwise}\hfill \end{array}$$
(9.14c)
$$A,B=1,\mathrm{},\text{rank}g,\alpha ,\beta ,\gamma \mathrm{\Delta }(g)$$
(9.14d)
and then the relations
$$:\widehat{H}_A(z)\widehat{H}_B(z):=:\widehat{H}_A(z)\widehat{H}_B(z):_M^{},:\widehat{H}_A(z)\widehat{E}_\alpha (z):=:\widehat{H}_A(z)\widehat{E}_\alpha (w):_M^{}$$
(9.15a)
$$:\widehat{E}_\alpha (z)\widehat{E}_\beta (z):=:\widehat{E}_\alpha (z)\widehat{E}_\beta (z):_M^{}+\{\begin{array}{ccc}\frac{\sigma \alpha d}{z}N_\gamma (\alpha ,\beta )\widehat{E}_\gamma (z)\hfill & & \text{if}\alpha +\beta =\gamma \hfill \\ \frac{\sigma \alpha d}{z}\alpha \widehat{H}(z)\frac{k\sigma \alpha d(\sigma \alpha d+1)}{2z^2}\hfill & & \text{if}\alpha +\beta =0\hfill \\ 0\hfill & & \text{otherwise}\hfill \end{array}$$
(9.15b)
express the OPE normal ordered products in terms of $`M^{}`$ ordering.
### 9.3 The Virasoro generators of $`A(H(d))/H(d)`$
For each inner-automorphic orbifold $`A(H(d))/H(d)`$, the stress tensor of twisted sector $`\sigma `$
$$\widehat{T}_\sigma =L^{ab}:\widehat{J}_a\widehat{J}_b:=:(L^{AB}\widehat{H}_A\widehat{H}_B+L^{A\alpha }(\widehat{H}_A\widehat{E}_\alpha +\widehat{E}_\alpha \widehat{H}_A)+L^{\alpha \beta }\widehat{E}_\alpha \widehat{E}_\beta ):$$
(9.16)
is obtained from (2.41a), (9.4b) and (9.9). For these orbifolds, all the sector dependence of $`\widehat{T}_\sigma `$ resides in the twisted currents.
With $`M^{}`$ ordering, the corresponding Virasoro generators of sector $`\sigma `$ are
$$L_\sigma (m)=L_\sigma ^q(m)+L_\sigma ^l(m)\delta _{m,0}\frac{k}{2}\underset{\alpha }{}L^{\alpha ,\alpha }(\sigma \alpha d)^2$$
(9.17a)
$`L_\sigma ^q(m)`$ $``$ $`{\displaystyle \underset{p}{}}\{{\displaystyle \underset{A,B}{}}L^{AB}:\widehat{H}_A(p)\widehat{H}_B(mp):_M^{}`$
$`+{\displaystyle \underset{A,\alpha }{}}L^{A\alpha }:[\widehat{H}_A(p)\widehat{E}_\alpha ((mp+\sigma \alpha d)\sigma \alpha d)`$
$`+\widehat{E}_\alpha ((p+\sigma \alpha d)\sigma \alpha d)\widehat{H}_A(mp)]:_M^{}`$
$`+{\displaystyle \underset{\alpha ,\beta }{}}L^{\alpha \beta }:\widehat{E}_\alpha (p\sigma \alpha d)\widehat{E}_\beta ((mp+\sigma (\alpha +\beta )d)\sigma \beta d):_M^{}\}`$
$`L_\sigma ^l(m)`$ $``$ $`{\displaystyle \underset{\alpha +\beta =\gamma }{}}L^{\alpha \beta }N_\gamma (\alpha ,\beta )(\sigma \alpha d)\widehat{E}_\gamma ((m+\sigma \gamma d)\sigma \gamma d)`$
$`+{\displaystyle \underset{\alpha }{}}L^{\alpha ,\alpha }(\sigma \alpha d)\alpha \widehat{H}(m)`$
$$L^{A\alpha }=0\text{ unless }\alpha d,L^{\alpha \beta }=0\text{ unless }(\alpha +\beta )d,\sigma =0,\mathrm{},\rho (1)1$$
(9.17d)
where $`q`$ and $`l`$ label the terms quadratic and linear in the twisted current modes. The result (9.3) follows from (9.11), (9.14d) and (9.16), or as a special case of the general $`M^{}`$-ordered orbifold Virasoro generators in (H.2c). In this result, we have written the arguments of the twisted root operators in the form
$$\widehat{E}_\alpha ((\text{integer})\sigma \alpha d)$$
(9.18)
to exhibit their proper modeing. To see that the quantities in the inner parentheses are integers, one must use<sup>x</sup><sup>x</sup>xFor example, consider the term proportional to $`\widehat{E}_\alpha \widehat{E}_\beta `$ in Eq. (9.3). The quantity in the inner parentheses of $`\widehat{E}_\beta `$ is an integer because the solution of the $`L^{\alpha \beta }`$ selection rule allows the term to contribute only when $`\sigma (\alpha +\beta )d`$. This term can also be written simply as $`L^{\alpha \beta }\widehat{E}_\alpha (p\sigma \alpha d)\widehat{E}_\beta (mp+\sigma \alpha d)`$. Since this term can contribute only when $`\widehat{E}_\beta `$ is in the twist class of $`\widehat{E}_\alpha `$, the simple expression is actually in the form $`\widehat{J}_{n(r)}\widehat{J}_{n(r)}`$ of the general result (H.2c). Similarly, all the other terms in (9.3) can be put in the form of (H.2c). the solutions (9.17d) of the $`L`$-selection rules.
As a simple example of (9.3), we use the affine-Sugawara construction (9.6) to obtain the Virasoro generators of the general inner-automorphic WZW orbifold:
$$\frac{A_g(H(d))}{H(d)},\sigma =0,\mathrm{},\rho (1)1$$
(9.19a)
$`(2k+Q_\psi )L_\sigma ^{\widehat{𝔤}(\sigma )}(m)`$ $`=`$ $`{\displaystyle \underset{p}{}}\{:{\displaystyle \underset{\alpha }{}}\widehat{E}_\alpha (p\sigma \alpha d)\widehat{E}_\alpha (mp+\sigma \alpha d):_M^{}`$
$`+{\displaystyle \underset{A}{}}:\widehat{H}_A(p)\widehat{H}_A(mp):_M^{}\}+Q_\psi \{\sigma d\widehat{H}(m)\delta _{m,0}{\displaystyle \frac{k}{2}}\sigma ^2d^2\}`$
$$[L_\sigma ^{\widehat{𝔤}(\sigma )}(m),\widehat{H}_A(n)]=n\widehat{H}_A(m+n)$$
(9.19c)
$$[L_\sigma ^{\widehat{𝔤}(\sigma )}(m),\widehat{E}_\alpha (n\sigma \alpha d)]=(n\sigma \alpha d)\widehat{E}_\alpha (m+n\sigma \alpha d).$$
(9.19d)
This result (including the fact that $`\widehat{E}`$, $`\widehat{H}`$ are twisted $`(1,0)`$ operators) is a special case of the result given for the general WZW orbifold in Eq. (6.22e). In sector $`\sigma `$, the orbifold affine-Sugawara construction $`L_\sigma ^{\widehat{𝔤}(\sigma )}`$ is equivalent (with $`\sigma dd`$) to the inner-automorphically twisted affine-Sugawara construction of Refs. References and References.
Other special cases included in the result (9.3) are the Virasoro generators of the general inner-automorphic coset orbifold
$$\frac{\frac{g}{h}(H(d))}{H(d)}\frac{A(H(d))}{H(d)}.$$
(9.20)
These orbifolds were discussed at the level of stress-tensor spectral flow in Ref. References.
### 9.4 Action on the untwisted affine vacuum
In this and the following subsection, we study the action of the orbifold Virasoro generators (9.3) on a particular state $`|0`$ which satisfies
$$(\widehat{H}_A(m)\delta _{m,0}\sigma kd_A)|0=\widehat{E}_\alpha (m\sigma \alpha d)|0=0\text{ when }m0.$$
(9.21)
According to the orbifold induction procedure$`^{\text{References},\text{References},\text{References},\text{References}}`$ for inner-automorphically twisted affine Lie algebras, the state $`|0`$ is the untwisted affine vacuum (see Eq. (9.26b)). (Except for small $`\{\sigma \alpha d\}`$, the untwisted affine vacuum $`|0`$ is not the true ground state $`|0_\sigma `$ of twisted sector $`\sigma `$, but, so far as we know, identification of the true ground state is an unsolved problem.) As emphasized in Ref. References, the untwisted affine vacuum $`|0`$ is a twisted affine highest weight state only so long as $`\sigma \alpha d>1`$ for all $`\alpha \mathrm{\Delta }(g)`$. One may compute
$$0|L_\sigma (m)|0=\delta _{m,0}\{k^2\underset{A,B}{}L^{AB}(\sigma d_A)(\sigma d_B)+\frac{k}{2}\underset{\alpha }{}L^{\alpha ,\alpha }(\sigma \alpha d)^2\}$$
(9.22)
for all $`\{\sigma \alpha d\}`$, but, as also emphasized in Ref. References, $`|0`$ is not Virasoro primary in general unless $`\sigma \alpha d>1`$ for all $`\alpha \mathrm{\Delta }(g)`$.
This situation can presumably be avoided by computing on twisted affine primary states or, as we will study here, by choosing only $`L^{AB}`$ and $`L^{\alpha ,\alpha }`$ not equal to zero in (9.3). In this case, we find that $`|0`$ is Virasoro primary
$$L_\sigma (m)=\underset{p}{}\underset{A,B}{}L^{AB}:\widehat{H}_A(p)\widehat{H}_B(mp):_M^{}$$
(9.23a)
$$+\underset{\alpha }{}L^{\alpha ,\alpha }\{\underset{p}{}:\widehat{E}_\alpha (p\sigma \alpha d)\widehat{E}_\alpha (mp+\sigma \alpha d):_M^{}+\sigma \alpha d(\alpha \widehat{H}(m)\delta _{m,0}\sigma \alpha d)\}$$
$$L_\sigma (m0)|0=\delta _{m,0}\widehat{\mathrm{\Delta }}_0(L,d;\sigma )|0,\sigma =1,\mathrm{},\rho (1)1$$
(9.23b)
$$\widehat{\mathrm{\Delta }}_0(L,d;\sigma )=k^2\underset{A,B}{}L^{AB}(\sigma d_A)(\sigma d_B)+\frac{k}{2}\underset{\alpha }{}L^{\alpha ,\alpha }(\sigma \alpha d)^2$$
(9.23c)
without restriction on $`\{\sigma \alpha d\}`$. Recall from (9.5b) that there is no $`H`$-invariance restriction on $`L^{AB}`$ or $`L^{\alpha ,\alpha }`$.
### 9.5 Connection with spectral flow
When there is an orbifold induction procedure, one may rewrite orbifold Virasoro generators in terms of untwisted current modes<sup>y</sup><sup>y</sup>yThe inverse of this unconventional step$`^{\text{References}}`$ was introduced by Freericks and Halpern, who studied the twisted current formulation of inner-automorphic WZW orbifolds starting from the spectral flow discussed below. $`J_a(m)`$ which satisfy (2.1b). Such procedures are known for inner-automorphically twisted affine Lie algebra$`^{\text{References},\text{References},\text{References},\text{References}}`$ and orbifold affine algebra$`^{\text{References}}`$
$$\widehat{H}_A(m)=H_A(m)+\delta _{m,0}k\sigma d_A,\widehat{E}_\alpha (m\sigma \alpha d)=E_\alpha (m)$$
(9.24a)
$$\widehat{J}_{aj}^{(r)}(m+\frac{r}{\rho (\sigma )})=J_{aj}(\rho (\sigma )m+r),r=0,\mathrm{},\rho (\sigma )1$$
(9.24b)
as well as the doubly-twisted affine algebras$`^{\text{References},\text{References}}`$ which combine (9.24a) and (9.24b).
As discussed in App. Appendix H., rewriting orbifold stress tensors in terms of untwisted currents leads to generically exotic and unfamiliar forms of the Virasoro generators $`L_\sigma (m)`$. In these forms one sees a generic mode imbalance, in which the modes of the untwisted currents do not sum to the integer $`m`$.
We rewrite here only the special case $`L^{AB},L^{\alpha ,\alpha }0`$ of the inner-automorphic orbifolds in (9.4), which avoids this mode imbalance phenomenon. In this case one finds
$$L_\sigma (m)=L(m)+D(L,d;\sigma )H(m)+\delta _{m,0}\widehat{\mathrm{\Delta }}_0(L,d;\sigma ),\sigma =0,\mathrm{},\rho (1)1$$
(9.25a)
$$L(m)\underset{p}{}\{\underset{A,B}{}L^{AB}:H_A(p)H_B(mp):+\underset{\alpha }{}L^{\alpha ,\alpha }:E_\alpha (p)E_\alpha (mp):\}$$
(9.25b)
$$D(L,d;\sigma )^A\sigma \underset{B}{}d_BM(L)^{BA},M(L)^{AB}=2kL^{AB}+\underset{\alpha }{}L^{\alpha ,\alpha }\alpha ^A\alpha ^B$$
(9.25c)
$$(L_\sigma (m0)\delta _{m,0}\widehat{\mathrm{\Delta }}_0(L,d;\sigma ))|0=0,\widehat{\mathrm{\Delta }}_0(L,d;\sigma )=\frac{k\sigma ^2}{2}\underset{AB}{}d_Ad_BM(L)^{AB}$$
(9.25d)
where the matrix $`M(L)^{AB}`$ is $`\delta ^{AC}M(L)_C^B`$ and $`M(L)_A^B`$ is the Cartan block of $`M(L)_a^b`$ in (5.1c). In (9.25d), the conformal weight $`\widehat{\mathrm{\Delta }}_0(L,d;\sigma )`$ of the untwisted affine vacuum $`|0`$ is the same as that given in (9.23c). All the current modes in (9.5) are untwisted and we have written the operator $`L(m)`$ in (9.25b) as an OPE normal ordered product
$$:J_a(m)J_b(n):=:J_a(m)J_b(n):_M^{}=\theta (m0)J_b(n)J_a(m)+\theta (m0)J_a(m)J_b(n)$$
(9.26a)
$$H_A(m0)|0=E_\alpha (m0)|0=L(m1)|0=0$$
(9.26b)
although OPE normal ordering is the same as $`M^{}`$ ordering for untwisted currents.
The form (9.5) (but not the more general case in (LABEL:unbalancedgroup)) provides an opportunity to check some results of the orbifold program against known results in affine-Virasoro theory. In particular, we want to compare the result (9.5) to the general $`c`$-fixed conformal deformation$`^{\text{References},\text{References},\text{References},\text{References}}`$ of the general affine-Virasoro construction:
$$L(m;D)=L^{ab}\underset{p}{}:J_a(p)J_b(mp):+D^aJ_a(m)+\mathrm{\Delta }_0(D)\delta _{m,0}$$
(9.27a)
$$L^{ab}=2L^{ac}G_{cd}L^{db}L^{cd}L^{ef}f_{ce}^af_{df}^bL^{cd}f_{cd}^ff_{df}^{(a}L^{b)e},c(D)=c=2G_{ab}L^{ab}$$
(9.27b)
$$D^bM(L)_b^a=D^a,M(L)_a^b=2G_{ac}L^{cb}+f_{ad}^eL^{dc}f_{ce}^b$$
(9.27c)
$$(L(m0;D)\delta _{m,0}\mathrm{\Delta }_0(D))|0=0,\mathrm{\Delta }_0(D)=\frac{1}{2}G_{ab}D^aD^b.$$
(9.27d)
Here the first term in (9.27a) is the undeformed affine-Virasoro construction with central charge $`c`$, and $`\mathrm{\Delta }_0(D)`$ is the conformal weight of the untwisted affine vacuum $`|0`$ under the deformed Virasoro $`L(m;D)`$ as a function of the deformation $`D`$. The matrix $`M(L)`$ in (9.27c) is the same matrix defined in (5.1c), and the eigenvector condition $`DM(L)=D`$ constrains the allowed values of the deformation $`D`$. The eigenvector condition is equivalent to the requirement that $`DJ`$ is a $`(1,0)`$ operator of the undeformed Virasoro. Because the scale of $`D`$ is not fixed by the eigenvector condition, the $`c`$-fixed conformal deformations are also called stress-tensor spectral flow$`^{\text{References},\text{References},\text{References},\text{References}}`$ which is equivalent$`^{\text{References}}`$ to inner-automorphic twisting or spectral flow of the underlying currents.$`^{\text{References},\text{References},\text{References},\text{References}}`$ The first example of stress-tensor spectral flow was given by Bardakci and Halpern in Ref. References.
To begin this comparison, we note first that the orbifold result (9.5) must be in the spectral flow (9.5), with the identifications
$$L_{\sigma =0}(m)=L(m)=L(m;D=0),L^{ab}=\{L^{AB},L^{\alpha ,\alpha }\},\widehat{c}(\sigma )=c(D)=c$$
(9.28a)
$$D^A=D(L,d;\sigma )^A,D^\alpha =0,\mathrm{\Delta }_0(D)=\widehat{\mathrm{\Delta }}_0(L,d;\sigma )$$
(9.28b)
because the first term $`L(m)`$ in (9.25a) is precisely the Virasoro generator $`L_{\sigma =0}(m)`$ we started with in the untwisted sector of each orbifold. For the affine-Sugawara construction on $`g`$ in (9.6), the two systems (9.5) and (9.5) are indeed equivalent with
$$L_\sigma ^{\widehat{𝔤}(\sigma )}(m)=L_g(m;D_g)=\frac{\eta ^{ab}}{2k+Q_\psi }\underset{p}{}:J_a(p)J_b(mp):+D_g^AJ_A(m)+\mathrm{\Delta }_0(D_g)\delta _{m,0}$$
(9.29a)
$$M(L_g)_a^b=\delta _a^b,D_g^A=D(L_g,d;\sigma )=\sigma d^A,\mathrm{\Delta }_0(D_g)=\widehat{\mathrm{\Delta }}_0(L_g,d;\sigma )=\frac{k}{2}D_g^2$$
(9.29b)
$$\widehat{c}_g(\sigma )=c_g(D)=c_g,a=1,\mathrm{},\text{dim}g,A=1,\mathrm{},\text{rank}g,\sigma =0,\mathrm{},\rho (1)1$$
(9.29c)
and these Virasoro generators are also the same as those in (9.3). For arbitrary $`D_g`$ the result (9.28b) is the canonical example of stress-tensor spectral flow first studied by Freericks and Halpern in Ref. References, and applied more recently in Refs. References and References. As emphasized in Ref. References, the orbifold sectors in (9.28b) are special points of the spectral flow.
But the more general identification given in (9.5) is surprising because the “deformation” $`D(L,d;\sigma )`$ in (9.25c) is a function of the inverse inertia tensor $`L^{ab}`$ and it is not obvious that the conformal weight $`\widehat{\mathrm{\Delta }}_0(L,d;\sigma )`$ in (9.25d) of the untwisted affine vacuum agrees with the spectral flow form $`\mathrm{\Delta }_0(D)`$ in (9.27d).
The key to understanding this identification is an unsuspected<sup>z</sup><sup>z</sup>zMotivation for a Lie $`h`$ invariance of the orbifold Virasoro generators $`L_{\sigma =0}(m)=L(m;D=0)`$ comes from the spectral-flow side of the identification: Lie $`h`$-invariant CFT’s with a non-trivial global component $`h_1`$ (see (6.5)) are the only known CFT’s with (1,0) operators. Conversely, for any Lie $`h`$-invariant CFT, arbitrary deformation by the (1,0) operators of $`h_1`$ automatically solve the spectral flow system (9.5). Lie $`h`$ invariance in the Virasoro generators $`L_{\sigma =0}(m)`$ of the untwisted sectors of these orbifolds!
To see this Lie invariance, one must verify the identity
$$\delta _AL^{ab}L^{c(a}f_{cA}^{b)}=iN(L)_A^{ab}=0,A=1,\mathrm{},\text{rank}g$$
(9.30)
which follows (when only $`L^{AB}`$ and $`L^{\alpha ,\alpha }`$ are non-zero) from the form of the structure constants in the Cartan-Weyl basis. The identity (9.30) tells us (see Subsec. 6.1) that the subset $`L^{AB}`$, $`L^{\alpha ,\alpha }0`$ of inverse inertia tensors in $`A(H(d))`$ also has a Lie symmetry
$$\text{Lie}h=\text{Cartan}g$$
(9.31)
in addition to the inner-automorphic invariance $`H(d)`$.
In the nomenclature of Subsec. 6.2, we have shown that $`L_{\sigma =0}(m)=L(m)`$ in (9.28a) describes a large set of doubly-invariant or ($`H`$ and Lie $`h`$)-invariant CFT’s which we will call
$$A(\text{Cartan}g(H(d)))A(H(d)),A(\text{Cartan}g(H(d)))A(\text{Lie}h).$$
(9.32)
It follows that the orbifold Virasoro generators $`L_\sigma (m)`$ in (9.4) or (9.5) describe the orbifolds
$$\frac{A(\text{Cartan}g(H(d)))}{H(d)}\frac{A(H(d))}{H(d)},\frac{A(\text{Cartan}g(H(d)))}{H(d)}\frac{A(\text{Lie}h(H))}{H}$$
(9.33)
by $`H(d)`$ of the doubly-invariant CFT’s $`A(\text{Cartan}g(H(d)))`$.
Returning to the untwisted sectors, the Cartan invariance (9.30) implies that
$$M(L)_A^CM(L)_C^B=M(L)_A^B,M(L)_A^\alpha =0$$
(9.34a)
$$[L(m),H_A(n)]=nM(L)_A^BH_B(m+n)$$
(9.34b)
$$[L(m),D^A(L,d;\sigma )H_A(n)]=nD(L,d;\sigma )^AH_A(m+n)$$
(9.34c)
at least for unitary CFT’s,$`^{\text{References},\text{References}}`$ where (9.34a,9.34a) are special cases of the general properties of Lie $`h`$-invariant CFT’s in (6.1). The relation (9.34c), which follows from (9.25c) and (9.34a,9.34a), says that $`DH`$ is a $`(1,0)`$ operator of the undeformed Virasoro. As noted above, this means that $`D^A(L,d;\sigma )`$ solves the eigenvector condition for $`D`$ in (9.27c). More explicitly, the eigenvector condition in (9.27c) takes the form
$$\sigma \underset{B}{}d^B\{M(L)_B^A\underset{C}{}M(L)_B^CM(L)_C^A\}=0$$
(9.35)
upon substitution of $`D(L,d;\sigma )`$ for $`D`$. Then, using $`M^2=M`$ in (9.34a), we see that the eigevector condition is satisfied identically for arbitrary $`\sigma `$ and $`d`$.
Finally, the equality of the conformal weights
$$\mathrm{\Delta }_0(D=D(L,d;\sigma ))=\frac{1}{2}kD^2(L,d;\sigma )=\widehat{\mathrm{\Delta }}_0(L,d;\sigma )$$
(9.36)
is established by using the definition of $`D(L,d;\sigma )`$ in (9.25c) and $`M^2=M`$ in (9.34a). This completes our check that, at least for unitary $`A(\text{Cartan}g(H(d)))`$, the orbifold systems (9.5) are points in the general spectral flow (9.5).
We close this section with a number of remarks:
$``$ Conditions on the vector $`d`$. The vector $`d`$ is determined in the orbifolds by the choice of $`H(d)`$, but, as we saw in (9.35), the vector $`d`$ is not determined by the eigenvector condition (9.27c) for the spectral flow. This is of course the phenomenon described by Freericks and Halpern in their original study$`^{\text{References}}`$ of inner-automorphic orbifolds as special points of stress-tensor spectral flow.
$``$ Unitarity revisited. In fact, the Lie $`h`$ relations (9.5) – and hence the conclusions above for the orbifolds – can be checked directly, without using unitarity. The Lie $`h`$ relations $`N(L)_A^{bc}=0`$ and $`M_A^\alpha =0`$ (which are necessary for (9.34b)) follow immediately from $`L^{AB},L^{\alpha ,\alpha }0`$. To check the last Lie $`h`$ relation $`M^2=M`$ in (9.34a) we need the reduced Virasoro master equation of the doubly-invariant CFT’s $`A(\text{Cartan}g(H(d)))`$
$$L^{AB}=2k\underset{D}{}L^{AD}L^{DB}\underset{\alpha }{}(L^{\alpha ,\alpha })^2\alpha ^A\alpha ^B+\underset{\alpha ,D}{}L^{\alpha ,\alpha }\alpha ^DL^{D(A}\alpha ^{B)}$$
(9.37a)
$$L^{\alpha ,\alpha }=2(k+\alpha ^2)(L^{\alpha ,\alpha })^2+\underset{\beta +\gamma =\alpha }{}(2L^{\alpha ,\alpha }L^{\gamma ,\gamma })L^{\beta ,\beta }N_\alpha ^2(\beta ,\gamma )$$
(9.37b)
$$c=2k(\underset{A}{}L^{AA}+\underset{\alpha }{}L^{\alpha ,\alpha })$$
(9.37c)
where the structure constants $`N_\gamma (\alpha ,\beta )`$ are defined in (9.1). The inverse inertia tensors of $`A(\text{Cartan}g(H(d)))`$ are controlled entirely by this system because there are no $`H`$-invariance restrictions on $`L^{AB}`$, $`L^{\alpha ,\alpha }0`$. The reduced Virasoro master equation (9.5) is a consistent subansatz of the Virasoro master equation, with the generically-expected number of inequivalent solutions at each level
$$N(g)=2^{n(g)}=2^{\frac{1}{2}((\text{rank}g)^2+\text{dim}g)}$$
(9.38)
where $`n(g)`$ is the number of equations and unknowns in the system.
Using the reduced master equation (9.5) to eliminate terms linear in $`L`$, we find that the relation $`M^2=M`$ in (9.34a) can be reduced to the consistency relation
$$\underset{\alpha ,\beta }{}\alpha ^A\beta ^BL^{\alpha ,\alpha }L^{\beta ,\beta }(\alpha \beta )=\underset{\alpha }{}\alpha ^A\alpha ^B\{2\alpha ^2(L^{\alpha ,\alpha })^2+\underset{\beta +\gamma =\alpha }{}2L^{\alpha ,\alpha }L^{\beta ,\beta }N_\alpha ^2(\beta ,\gamma )$$
$$\underset{\beta +\gamma =\alpha }{}L^{\gamma ,\gamma }L^{\beta ,\beta }N_\alpha ^2(\beta ,\gamma )\}$$
(9.39)
and this relation is in fact an identity because
$$N_\gamma (\alpha ,\beta )=N_\beta (\gamma ,\alpha )=N_\gamma (\alpha ,\beta )$$
(9.40a)
$$\alpha \beta =(\delta _{\alpha +\beta ,0}\delta _{\alpha +\beta ,0})\alpha ^2+\delta _{\alpha +\beta ,\gamma ^{}}N_\gamma ^{}^2(\alpha ,\beta )\delta _{\alpha +\beta ,\gamma }N_\gamma ^2(\alpha ,\beta )$$
(9.40b)
where (9.40b) follows from the Jacobi identity of Lie $`g`$. This completes the check that the Lie $`h`$ relations (9.5) are true independent of unitarity, and we conjecture that the general Lie $`h`$ relations in (6.1) are similarly true independent of unitarity.
$``$ $`K`$-conjugation in spectral flow. We note that orbifold $`K`$-conjugation
$$L_\sigma (m;D(L,d;\sigma ))+\stackrel{~}{L}_\sigma (m;D(\stackrel{~}{L},d;\sigma ))=L_\sigma ^{\widehat{𝔤}(\sigma )}(m;D(L_g,d;\sigma ))$$
(9.41a)
$$L+\stackrel{~}{L}=L_g,c+\stackrel{~}{c}=c_g,[L_\sigma (m;D(L,d;\sigma )),\stackrel{~}{L}_\sigma (n;D(\stackrel{~}{L},d;\sigma ))]=0$$
(9.41b)
holds for the special points of the stress-tensor spectral flow which describe the orbifolds $`A(\text{Cartan}g(H(d)))/H(d)`$. Here the modes of $`L_\sigma ^{\widehat{𝔤}(\sigma )}`$ are given in Eq. (9.28b), and (9.5) is a special case of the general orbifold $`K`$-conjugation in (6.4). Eq. (9.5) is the first observation of $`K`$-conjugation in stress-tensor spectral flow, and $`K`$-conjugation in the general spectral flow (9.5) deserves further study.
$``$ Twisted $`h`$ currents. We return for a moment to the twisted current formulation of the orbifolds $`A(\text{Cartan}g(H(d)))/H(d)`$. Another consequence of the Lie$`h=`$Cartan$`g`$ invariance of $`A(\text{Cartan}g(H(d)))`$ is the orbifold statement
$$[L_\sigma (m),\widehat{H}_A(n)]=nM(L)_A^B\widehat{H}_B(m+n),A=1,\mathrm{},\text{rank}g,\sigma =\rho (1)1$$
(9.42)
where $`L_\sigma (m)`$ is given in (9.4) and $`\widehat{H}_A(m)`$ are the Cartan modes of sector $`\sigma `$. In a left eigenbasis of $`M(L)`$, Eq. (9.42) is a special case of the more general result (6.20).
$``$ Beyond spectral flow. We have seen that $`A(\text{Cartan}g(H(d)))/H(d)`$ is a large subset of inner-automorphic orbifolds which can also be described by stress-tensor spectral flow. We emphasize however that Eqs. (9.3) and (LABEL:unbalancedgroup) contain the Virasoro generators of a much larger class of inner-automorphic orbifolds which apparently (due to the mode-imbalance phenomenon discussed in App. Appendix H.) can not be described by spectral flow.
$``$ SL(2,R) WZW models. We finally note that the stress tensors proposed for SL(2,R) WZW models in Ref. References are special cases (with $`\sigma \alpha d`$ for $`g=`$SL(2,R)) of the inner-automorphic WZW orbifold stress tensors (9.3) or (9.28b).
Acknowledgements
We thank P. Bouwknegt, R. Dijkgraaf, M. Gaberdiel, A. Giveon, H. Ooguri, N. Nekrasov, N. Reshetikin and K. Skenderis for helpful discussions. We also thank J. Evslin for his participation in the early stages of this paper, and C. Schweigert for reading the manuscript.
J. E. W. was supported by the Department of Education, GAANN. M. B. H. was supported in part by the Director, Office of Science, Office of High Energy and Nuclear Physics, Division of High Energy Physics, of the U.S. Department of Energy under Contract DE-AC03-76SF00098 and in part by the National Science Foundation under grant PHY95-14797.
### Appendix A. Incorporation of the Selection Rules
In this appendix we incorporate the solutions of various selection rules to obtain reduced forms of relations in the text.
For example, we may combine the solution (2.29e) of the $``$-selection rule and the Jacobi identity (2.29b) for $``$ to obtain the reduced form of the Jacobi identity
$$_\mu ^{}\{_{n(r)\mu ;n(s)\nu }^{n(r)+n(s),\mu ^{}}(\sigma )_{n(t)\delta ;n(r)+n(s),\mu ^{}}^{n(r)+n(s)+n(t),\gamma }(\sigma )$$
$$+_{n(t)\delta ;n(r)\mu }^{n(r)+n(t),\mu ^{}}(\sigma )_{n(s)\nu ;n(r)+n(t),\mu ^{}}^{n(r)+n(s)+n(t),\gamma }(\sigma )$$
$$+_{n(s)\nu ;n(t)\delta }^{n(s)+n(t),\mu ^{}}(\sigma )_{n(r)\mu ;n(s)+n(t),\mu ^{}}^{n(r)+n(s)+n(t),\gamma }(\sigma )\}=0$$
(A.1)
in which no spectral indices are summed. Similarly the solution of the $``$-selection rule and (2.29c) imply the relation
$$_{n(r)\mu ;n(s)\nu ;(n(r)+n(s)),\delta }(\sigma )=_{n(r)\mu ;(n(r)+n(s)),\delta ;n(s)\nu }$$
(A.2)
among the reduced, twisted totally-antisymmetric structure constants of sector $`\sigma `$. The relations (A.1) and (A.2) are found to guarantee the Jacobi identity of the general twisted current algebra $`\widehat{𝔤}(\sigma )`$ in Eq. (3.1).
Using the solution (4.1e) of the selection rule for the orbifold conjugation matrix $``$, we obtain the reduced form of (4.1c)
$$_\delta _{n(r)\mu }^{n(r),\delta }(\sigma )_{n(r),\delta }^{n(r)\nu }(\sigma )^{}=_\delta _{n(r)\mu }^{n(r),\delta }(\sigma )^{}_{n(r),\delta }^{n(r)\nu }(\sigma )=\delta _\mu ^\nu $$
(A.3)
which is needed to verify Eq. (4.4b).
Similarly, the reduced forms of (4.2a), (4.2b) and (4.2c)
$$𝒢_{n(r)\mu ;n(r),\nu }(\sigma )^{}=\underset{\mu ^{},\nu ^{}}{}_{n(r)\mu }^{n(r),\mu ^{}}(\sigma )_{n(r),\nu }^{n(r)\nu ^{}}(\sigma )𝒢_{n(r),\mu ^{};n(r)\nu ^{}}(\sigma )$$
(A.4a)
$$_{n(r)\mu ;n(r),\nu }^{0\delta }(\sigma )^{}=\underset{\mu ^{},\nu ^{},\delta ^{}}{}_{n(r)\mu }^{n(r),\mu ^{}}(\sigma )_{n(r),\nu }^{n(r)\nu ^{}}(\sigma )_{n(r),\mu ^{};n(r)\nu ^{}}^{0\delta ^{}}(\sigma )_{0\delta ^{}}^{0\delta }(\sigma )^{}$$
(A.4b)
$$^{n(r)\mu ;n(r),\nu }(\sigma )^{}=\underset{\mu ^{},\nu ^{}}{}^{n(r),\mu ^{};n(r)\nu ^{}}(\sigma )_{n(r),\mu ^{}}^{n(r)\mu }(\sigma )^{}_{n(r)\nu ^{}}^{n(r),\nu }(\sigma )^{}$$
(A.4c)
are needed to verify (4.7) and (7.6), and to check the consistency of (4.5).
### Appendix B. The Projectors of $`A(H)/H`$
The projectors $`𝒫(\overline{n}(r);\sigma )`$ onto the $`\overline{n}(r)`$ subspaces of sector $`\sigma `$,
$$𝒫(\overline{n}(r);\sigma )_a^b_\mu U^{}(\sigma )_a^{n(r)\mu }U(\sigma )_{n(r)\mu }^b,𝒫(\overline{n}(r);\sigma )_a^a=\text{dim}[\overline{n}(r)]$$
(B.1a)
$$𝒫(\overline{n}(r)\pm \rho (\sigma );\sigma )=𝒫(\overline{n}(r);\sigma ),𝒫(\overline{n}(r);\sigma )=𝒫(\rho (\sigma )\overline{n}(r);\sigma )$$
(B.1b)
$$𝒫(\overline{n}(r);\sigma )𝒫(\overline{n}(s);\sigma )=\delta _{n(r)}^{n(s)}𝒫(\overline{n}(r);\sigma ),_r𝒫(\overline{n}(r);\sigma )=1$$
(B.1c)
$$\omega (h_\sigma )=_rE_{n(r)}(\sigma )𝒫(\overline{n}(r);\sigma ),[\omega (h_\sigma ),𝒫(\overline{n}(r);\sigma )]=0$$
(B.1d)
$$𝒫(\overline{n}(r);\sigma )_a^c𝒫(\overline{n}(s);\sigma )_b^dG_{cd}(1E_{n(r)}E_{n(s)})=0$$
(B.1e)
$$𝒫(\overline{n}(r);\sigma )_a^d𝒫(\overline{n}(s);\sigma )_b^ef_{de}^f𝒫(\overline{n}(t);\sigma )_f^c(1E_{n(r)}E_{n(s)}E_{n(t)}^{})=0$$
(B.1f)
$$L^{ab}𝒫(\overline{n}(r);\sigma )_a^c𝒫(\overline{n}(r);\sigma )_b^dG_{cd}=L^{ac}G_{cb}𝒫(\overline{n}(r);\sigma )_a^b$$
(B.1g)
are defined for all orbifolds $`A(H)/H`$. The identity in (B.1g) can be proven by using duality transformations to write out $`_{r,s,\mu ,\nu }f(n(r))^{n(r)\mu ;n(s)\nu }𝒢_{n(r)\mu ;n(s)\nu }`$ in two ways, both with and without the selection rules, for arbitrary periodic $`f`$. These projectors were encountered in Sec. 7 and will play a central role in App. Appendix B..
### Appendix C. A Constrained Basis for Twisted Currents
We consider another basis for the twisted currents of $`A(H)/H`$
$$\widehat{J}_a^{n(r)}(m+\frac{n\left(r\right)}{\rho \left(\sigma \right)})_\mu \chi \left(\sigma \right)_{n\left(r\right)\mu }^1U^{}(\sigma )_a^{n(r)\mu }\widehat{J}_{n(r)\mu }(m+\frac{n\left(r\right)}{\rho \left(\sigma \right)})$$
(C.1a)
$$\widehat{J}_a^{n(r)}(m+\frac{n\left(r\right)}{\rho \left(\sigma \right)})^{}=\rho _a^b\widehat{J}_b^{n(r)}(m\frac{n\left(r\right)}{\rho \left(\sigma \right)})$$
(C.1b)
$$𝒫(\overline{n}(s);\sigma )_a^b\widehat{J}_b^{n(r)}(m+\frac{n\left(r\right)}{\rho \left(\sigma \right)})=\widehat{J}_a^{n(s)}(m+\frac{n\left(s\right)}{\rho \left(\sigma \right)})\delta _{n(s)}^{n(r)}$$
(C.1c)
$$a=1,\mathrm{},\text{dim}g,\overline{n}(r)\{0,\mathrm{},\rho (\sigma )1\}$$
(C.1d)
where $`\{𝒫(\overline{n}(r);\sigma )\}`$ is the set of projectors in App. Appendix A.. These twisted currents carry the original Lie algebra index $`a`$, but the basis is overcomplete with the constraints (C.1c).
In this basis, the twisted current algebra and the orbifold stress tensors take the form
$$\widehat{T}_\sigma (z)=_rL^{ab}:\widehat{J}_a^{n(r)}(z)\widehat{J}_b^{n(r)}(z):$$
(C.2a)
$$\widehat{J}_a^{n(r)}(z)=\underset{m}{}\widehat{J}_a^{n(r)}(m+\frac{n\left(r\right)}{\rho \left(\sigma \right)})z^{(m+\frac{n(r)}{\rho (\sigma )})1},\widehat{J}_a^{n(r)}(ze^{2\pi i})=E_{n(r)}(\sigma )\widehat{J}_a^{n(r)}(z)$$
(C.2b)
$$[\widehat{J}_a^{n(r)}(m+\frac{n\left(r\right)}{\rho \left(\sigma \right)}),\widehat{J}_b^{n(s)}(n+\frac{n\left(s\right)}{\rho \left(\sigma \right)})]=𝒫(\overline{n}(r);\sigma )_a^c𝒫(\overline{n}(s);\sigma )_b^d\{(m+\frac{n\left(r\right)}{\rho \left(\sigma \right)})G_{cd}\delta _{m+n+\frac{n(r)+n(s)}{\rho (\sigma )},0}$$
$$+if_{cd}^e\widehat{J}_e^{n(r)+n(s)}(m+n+\frac{n\left(r\right)+n\left(s\right)}{\rho \left(\sigma \right)})\}$$
(C.2c)
where $`L^{ab}`$ is the $`H`$-invariant inverse inertia tensor of the untwisted sector of the orbifold. The Jacobi identity for (C.2c) is verified with the selection rules (B.1e,LABEL:psglett).
### Appendix D. Conversion from $`n(r)`$ to $`\overline{n}(r)`$
In Eq. (3.9), various integers $`n(r)`$ are converted into their corresponding twist class $`\overline{n}(r)`$. To see how this happens, follow the steps
$$\frac{1}{z}(\frac{w}{z})^{\frac{n(r)}{\rho (\sigma )}}\underset{m}{}\theta (m+\frac{n\left(r\right)}{\rho \left(\sigma \right)}0)(\frac{w}{z})^m$$
(D.1a)
$`=`$ $`{\displaystyle \frac{1}{z}}({\displaystyle \frac{w}{z}})^{\frac{\overline{n}(r)}{\rho (\sigma )}}{\displaystyle \underset{m^{}}{}}\theta (m^{}+\frac{\overline{n}\left(r\right)}{\rho \left(\sigma \right)}0)({\displaystyle \frac{w}{z}})^m^{},0{\displaystyle \frac{\overline{n}(r)}{\rho (\sigma )}}<1`$ (D.1b)
$`=`$ $`{\displaystyle \frac{1}{z}}({\displaystyle \frac{w}{z}})^{\frac{\overline{n}(r)}{\rho (\sigma )}}{\displaystyle \underset{m^{}=0}{\overset{\mathrm{}}{}}}({\displaystyle \frac{w}{z}})^m^{}=({\displaystyle \frac{w}{z}})^{\frac{\overline{n}(r)}{\rho (\sigma )}}{\displaystyle \frac{1}{zw}}`$ (D.1c)
where we have used (2.16) and the change of variable $`m^{}=m+n(r)/\rho (\sigma )`$ to obtain (D.1b).
### Appendix E. Direct Computation of the $`\widehat{T}\widehat{J}`$ OPE’s
Orbifold operator products of the general currents $`\widehat{J}_{n(r)\mu }`$ can be obtained in one step by the prescription
$$ran(r)\mu ,\widehat{J}_a^{(r)}\widehat{J}_{n(r)\mu }$$
(E.1)
from the formulas of Appendix A of Ref. References. For the operator product $`\widehat{T}_\sigma \widehat{J}(\sigma )`$, one finds the same OPE’s as in Eq. (5), but the forms of the twisted tensors $`(\sigma )`$ and $`𝒩(\sigma )`$ which result from this computation are:
$$_{n(r)\mu }^{n(s)\nu }(\sigma )=^{n(t)\delta ;n(u)ϵ}(\sigma )_{n(t)\delta ;n(u)ϵ;n(r)\mu }^{n(s)\nu }(\sigma )$$
(E.2a)
$$𝒩_{n(r)\mu }^{n(s)\nu ;n(t)\delta }(\sigma )=^{n(u)ϵ;n(v)\gamma }(\sigma )𝒩_{n(u)ϵ;n(v)\gamma ;n(r)\mu }^{n(s)\nu ;n(t)\delta }(\sigma )$$
(E.2b)
$`_{n(r)\mu ;n(s)\nu ;n(t)\delta }^{n(u)ϵ}(\sigma )`$ $`=`$ $`\delta _{(n(r)\mu }^{n(u)ϵ}𝒢_{n(s)\nu );n(t)\delta }(\sigma )`$ (E.2c)
$`+{\displaystyle \frac{1}{2}}_{n(v)\gamma ;(n(r)\mu }^{n(u)ϵ}(\sigma )_{n(s)\nu );n(t)\delta }^{n(v)\gamma }(\sigma )`$
$$𝒩_{n(r)\mu ;n(s)\nu ;n(t)\delta }^{n(u)ϵ;n(v)\gamma }(\sigma )=\frac{i}{2}\delta _{(n(r)\mu }^{(n(u)ϵ}_{n(s)\nu );n(t)\delta }^{n(v)\gamma )}(\sigma ).$$
(E.2d)
As expected, these results can be obtained by the duality algorithm (2.36) from the untwisted $`TJ`$ OPE’s in (5) and the standard relations among the untwisted tensors$`^{\text{References},\text{References},\text{References},\text{References},\text{References}}`$
$$M(L)_a^b=L^{cd}M_{cda}^b,N(L)_a^{bc}=L^{de}N_{dea}^{bc}$$
(E.3a)
$$M_{abc}^d=\delta _{(a}^dG_{b)c}+\frac{1}{2}f_{e(a}^df_{b)c}^e,N_{abc}^{de}=\frac{1}{2}\delta _{(a}^{(d}f_{b)c}^{e)}.$$
(E.3b)
Using the duality transformations (2.24d,2.24d) and (2.33c) for $`𝒢,`$ and $``$, one finds that the forms of $`(\sigma )`$ and $`𝒩(\sigma )`$ in (LABEL:mcncappgroup) are equivalent to the duality transformations for $`(\sigma )`$ and $`𝒩(\sigma )`$ in (5.4b,5.4b).
### Appendix F. $`A(S_3)/S_3`$ and $`A(S_4)/S_4`$
$``$$`A(S_3)/S_3`$
We begin this appendix by checking the results of Sec. 8 against the results for the permutation orbifolds $`A(S_3)/S_3`$ given earlier in Ref. References. The permutation orbifolds $`A(S_3)/S_3`$ have three sectors named by the partitions $`\stackrel{}{\sigma }=\{1,1,1\}`$, $`\{2,1\}`$ and $`\{3\}`$. The first of these is the untwisted sector (see (8.3)), where the ambient algebra $`g`$ consists of three copies of an untwisted ($`\sigma _0=\sigma _1=\sigma _2=1`$) affine algebra at level $`\widehat{k}_0=\widehat{k}_1=\widehat{k}_2=k`$.
The twisted sector $`\stackrel{}{\sigma }=\{3\}`$ has an order $`\sigma _0=3`$ orbifold affine algebra$`^{\text{References}}`$ whose currents $`\widehat{J}_{a,j=0}^{(\widehat{j})}`$, $`\widehat{j}=0,1,2`$ have orbifold affine level $`\widehat{k}_0=\sigma _0k=3k`$. The stress tensor, central charge and ground state conformal weight of this sector are
$$\widehat{T}_{\{3\}}=\frac{\lambda ^{ab}}{3}:(\widehat{J}_{a0}^{(0)}\widehat{J}_{b0}^{(0)}+\widehat{J}_{a0}^{(1)}\widehat{J}_{b0}^{(1)}+\widehat{J}_{a0}^{(2)}\widehat{J}_{b0}^{(2)}):+l^{ab}:\widehat{J}_{a0}^{(0)}\widehat{J}_{b0}^{(0)}:$$
(F.1a)
$$\widehat{c}(\{3\})=c=6k\eta _{ab}(\lambda ^{ab}+l^{ab}),\widehat{\mathrm{\Delta }}_0(\{3\})=\frac{2k\eta _{ab}\lambda ^{ab}}{9}.$$
(F.1b)
With the identification $`\widehat{J}_a^{(\widehat{j})}\widehat{J}_{a,j=0}^{(\widehat{j})}`$ and the symmetry of the bilinears given in Ref. References, the result (LABEL:z3sec) is recognized as the stress tensor and ground state conformal weight of the sector $`\omega _1`$ given in Ref. References.
Finally, the twisted sector $`\stackrel{}{\sigma }=\{2,1\}`$ has an order $`\sigma _0=2`$ orbifold affine algebra whose currents $`\widehat{J}_{a,j=0}^{(\widehat{j})}`$, $`\widehat{j}=0,1`$ have orbifold affine level $`\widehat{k}_0=\sigma _0k=2k`$, and a commuting order $`\sigma _1=1`$ untwisted affine algebra with currents $`\widehat{J}_{a,j=1}^{(0)}`$ at level $`\widehat{k}_1=k`$. The stress tensor, central charge and ground state conformal weight of this sector are
$`\widehat{T}_{\{2,1\}}`$ $`=`$ $`{\displaystyle \frac{\lambda ^{ab}}{2}}:(\widehat{J}_{a0}^{(0)}\widehat{J}_{b0}^{(0)}+\widehat{J}_{a0}^{(1)}\widehat{J}_{b0}^{(1)}):+\lambda ^{ab}:\widehat{J}_{a1}^{(0)}\widehat{J}_{b1}^{(0)}:`$ (F.2a)
$`+`$ $`l^{ab}:(\widehat{J}_{a0}^{(0)}\widehat{J}_{b0}^{(0)}+\widehat{J}_{a0}^{(0)}\widehat{J}_{b1}^{(0)}+\widehat{J}_{a1}^{(0)}\widehat{J}_{b0}^{(0)}+\widehat{J}_{a1}^{(0)}\widehat{J}_{b1}^{(0)}):`$
$$\widehat{c}(\{2,1\})=c=6k\eta _{ab}(\lambda ^{ab}+l^{ab}),\widehat{\mathrm{\Delta }}_0(\{2,1\})=\frac{k\eta _{ab}\lambda ^{ab}}{8}.$$
(F.2b)
With the identifications $`\widehat{J}_a^{(\widehat{j})}\widehat{J}_{a,j=0}^{(\widehat{j})}`$, $`J_a\widehat{J}_{a,j=1}^{(0)}`$ this is recognized as the stress tensor and ground state conformal weight of sector $`\omega _2`$ in Ref. References. Agreement is obtained in this case because $`\omega _2`$ and our $`\omega (\stackrel{}{\sigma }=\{2,1\})`$ are in the same conjugacy class of $`S_3`$.
$``$$`A(S_4)/S_4`$
We turn next to the permutation orbifolds $`A(S_4)/S_4`$, which have five sectors named by the partitions $`\stackrel{}{\sigma }`$ given in Eq. (8.9). Consider first the sector $`\stackrel{}{\sigma }=\{3,1\}`$ with the identifications
$$\widehat{J}_a^{(\widehat{j})}\widehat{J}_{a,j=0}^{(\widehat{j})},\widehat{j}=0,1,2;J_a\widehat{J}_{a,j=1}^{(0)}$$
(F.3)
where $`\widehat{J}`$ is an orbifold affine algebra of order three and $`J`$ is a commuting integral affine algebra. In this notation, we obtain
$`\widehat{T}_{\{3,1\}}`$ $`=`$ $`{\displaystyle \frac{\lambda ^{ab}}{3}}:(\widehat{J}_a^{(0)}\widehat{J}_b^{(0)}+\widehat{J}_a^{(1)}\widehat{J}_b^{(1)}+\widehat{J}_a^{(2)}\widehat{J}_b^{(2)}):+\lambda ^{ab}:J_aJ_b:`$ (F.4a)
$`+`$ $`l^{ab}:(\widehat{J}_a^{(0)}\widehat{J}_b^{(0)}+\widehat{J}_a^{(0)}J_b+J_a\widehat{J}_b^{(0)}+J_aJ_b):`$
$$\widehat{c}(\{3,1\})=c=8k\eta _{ab}(\lambda ^{ab}+l^{ab}),\widehat{\mathrm{\Delta }}_0(\{3,1\})=\frac{2k\eta _{ab}\lambda ^{ab}}{9}$$
(F.4b)
from Eq. (8.22b). The remaining ground state conformal weights of $`A(S_4)/S_4`$ are easily computed
$$\widehat{\mathrm{\Delta }}_0(\{4\},\{2,2\},\{2,1,1\})=k\eta _{ab}\lambda ^{ab}(\frac{5}{16},\frac{1}{4},\frac{1}{8})$$
(F.5)
from the partitions $`\stackrel{}{\sigma }`$ in Eq. (8.9).
### Appendix G. The Grade Automorphism Group
The grade (inner) automorphism group is defined by
$$d=\frac{1}{h_g}\underset{i=1}{\overset{r}{}}\frac{\lambda _i}{\alpha _i^2},\alpha d=\frac{1}{2h_g}G(\alpha ),\alpha \mathrm{\Delta }(g)$$
(G.1a)
$$G(\alpha )=\underset{i=0}{\overset{r}{}}n_i(\alpha )\text{when }\alpha =\underset{i=0}{\overset{r}{}}n_i(\alpha )\alpha _i$$
(G.1b)
$$N_c=\rho (\sigma =1)=2h_g,\sigma =0,\mathrm{},2h_g1.$$
(G.1c)
Here $`r=\text{rank}g`$, $`h_g`$ is the Coxeter number of $`g`$ and $`G(\alpha )`$ is the grade of $`\alpha \mathrm{\Delta }(g)`$. The order $`\rho (1)`$ of the $`\sigma =1`$ automorphism given in (G.1c) is computed from (2.9) and the fact that $`G(\psi _g)=h_g1`$, where $`\psi _g`$ is the highest root of $`g`$. Using (2.16), one may also compute the twist classes of the twisted currents $`\widehat{J}`$ of sector $`\sigma =1`$
$$\overline{n}_\alpha (\sigma =1)=\{\begin{array}{ccc}2h_gG(\alpha ),\alpha >0\hfill & & \\ G(\alpha ),\alpha <0\hfill & & \end{array}$$
(G.2)
and the other sectors may be similarly analyzed.
Invariance under the grade automorphism group restricts the inverse inertia tensors of the untwisted sectors to the form $`L^{AB}`$, $`L^{\alpha ,\alpha }0`$, so the stress tensors of the orbifolds by the grade automorphism group are included in Eqs. (9.4), (9.5) and (9.5).
### Appendix H. Another Mode Ordering
$`M^{}`$ ordering is mode ordering by the integer part $`m`$ of the mode number ($`m+\frac{n(r)}{\rho (\sigma )}`$),
$$:\widehat{J}_{n(r)\mu }(m+\frac{n\left(r\right)}{\rho \left(\sigma \right)})\widehat{J}_{n(s)\mu }(n+\frac{n\left(s\right)}{\rho \left(\sigma \right)}):_M^{}\theta (m0)\widehat{J}_{n(s)\mu }(n+\frac{n\left(s\right)}{\rho \left(\sigma \right)})\widehat{J}_{n(r)\mu }(m+\frac{n\left(r\right)}{\rho \left(\sigma \right)})$$
$$+\theta (m<0)\widehat{J}_{n(r)\mu }(m+\frac{n\left(r\right)}{\rho \left(\sigma \right)})\widehat{J}_{n(s)\mu }(n+\frac{n\left(s\right)}{\rho \left(\sigma \right)}).$$
(H.1)
With $`M^{}`$ ordering, we find in place of Eqs. (3.9), (3.11) and (3.13c):
$$\widehat{J}_{n(r)\mu }(z)\widehat{J}_{n(s)\nu }(w)=(\frac{w}{z})^{\frac{n(r)}{\rho (\sigma )}}\{[\frac{1}{(zw)^2}+\frac{n(r)/\rho (\sigma )}{w(zw)}]𝒢_{n(r)\mu ;n(s)\nu }(\sigma )$$
$$+\frac{i_{n(r)\mu ;n(s)\nu }^{n(r)+n(s),\delta }(\sigma )\widehat{J}_{n(r)+n(s),\delta }(w)}{zw}\}+:\widehat{J}_{n(r)\mu }(z)\widehat{J}_{n(s)\nu }(w):_M^{}$$
(H.2a)
$`\widehat{T}_\sigma (z)`$ $`=`$ $`{\displaystyle \underset{r,\mu ,\nu }{}}^{n(r)\mu ;n(r),\nu }(\sigma )\{:\widehat{J}_{n(r)\mu }(z)\widehat{J}_{n(r),\nu }(z):_M^{}{\displaystyle \frac{in(r)}{\rho (\sigma )}}_{n(r)\mu ;n(r)\nu }^{0\delta }(\sigma ){\displaystyle \frac{\widehat{J}_{0\delta }(z)}{z}}`$ (H.2b)
$`+{\displaystyle \frac{1}{z^2}}{\displaystyle \frac{n(r)}{2\rho (\sigma )}}(1{\displaystyle \frac{n(r)}{\rho (\sigma )}})𝒢_{n(r)\mu ;n(r),\nu }(\sigma )\}`$
$$L_\sigma (m)=\underset{r,\mu ,\nu }{}^{n(r)\mu ;n(r),\nu }(\sigma )\{\underset{p}{}:\widehat{J}_{n(r)\mu }(p+\frac{n\left(r\right)}{\rho \left(\sigma \right)})\widehat{J}_{n(r),\nu }(mp\frac{n\left(r\right)}{\rho \left(\sigma \right)}):_M^{}$$
$$i\frac{n(r)}{\rho (\sigma )}_{n(r)\mu ;n(r),\nu }^{0\delta }(\sigma )\widehat{J}_{0\delta }(m)+\delta _{m,0}\frac{n(r)}{2\rho (\sigma )}(1\frac{n(r)}{\rho (\sigma )})𝒢_{n(r)\mu ;n(r),\nu }(\sigma )\}.$$
(H.2c)
Curiously, the $`M^{}`$ results (LABEL:MpNOappgroup) can be obtained by the map $`MM^{},`$ $`\overline{n}(r)n(r)`$ from their $`M`$-counterparts in (3.9), (3.11) and (3.13c). The $`M^{}`$ ordered product in (H.2a) is not periodic under $`n(r)n(r)+\rho (\sigma )`$, but the change is compensated by the $`𝒢`$ term so that the operator product on the left is periodic. Similarly, the total summands of (H.2b,LABEL:LmodeMplett) are periodic, so that each $`n(r)`$ can be replaced by $`\overline{n}(r)`$.
Because of its relation to the untwisted affine vacuum $`|0`$, $`M^{}`$ ordering is used to discuss the general inner-automorphic orbifold in Sec. 9.
### Appendix I. The Mode Imbalance Phenomenon
We begin this appendix by using the orbifold induction procedure (9.24a) to rewrite the stress tensors (9.3) of the inner-automorphic orbifolds in terms of untwisted currents. The result is
$$L_\sigma (m)=L_\sigma ^q^{}(m)+L_\sigma ^l^{}(m)+\delta _{m,0}\widehat{\mathrm{\Delta }}_0(L,d;\sigma )$$
(I.1a)
$$L_\sigma ^q^{}(m)\underset{p}{}\{\underset{A,B}{}L^{AB}:H_A(p)H_B(mp):_M^{}$$
$$+\underset{A,\alpha }{}L^{A\alpha }:[H_A(p)E_\alpha (mp+\sigma \alpha d)+E_\alpha (p+\sigma \alpha d)H_A(mp)]:_M^{}$$
$$+\underset{\alpha ,\beta }{}L^{\alpha \beta }:E_\alpha (p)E_\beta (mp+\sigma (\alpha +\beta )d):_M^{}\}$$
(I.1b)
$$L_\sigma ^l^{}(m)\underset{\alpha +\beta =\gamma }{}L^{\alpha \beta }N_\gamma (\alpha ,\beta )(\sigma \alpha d)E_\gamma (m+\sigma \gamma d)+2k\underset{A,\alpha }{}\sigma d_AL^{A\alpha }E_\alpha (m+\sigma \alpha d)$$
$$+2k\underset{A,B}{}\sigma d_AL^{AB}H_B(m)+\underset{\alpha }{}L^{\alpha ,\alpha }(\sigma \alpha d)\alpha H(m)$$
(I.1c)
$$L^{A\alpha }=0\text{ unless }\alpha d,L^{\alpha \beta }=0\text{ unless }(\alpha +\beta )d,\sigma =0,\mathrm{},\rho (1)1$$
(I.1d)
where the quantity $`\widehat{\mathrm{\Delta }}_0(L,d;\sigma )`$ is given in (9.23c). In this form, the arguments (modes) of the root operators are exactly the quantities in the inner parentheses of Eq. (9.3), so (according to the discussion around (9.18)) all the modes here are integers.
In spite of their unfamiliar form, these constructions are Virasoro generators when their inverse inertia tensors satisfy the H-invariance conditions (I.1d) and the Virasoro master equation (2.3c). Although the current modes were balanced (i.e. summed to $`m`$) when the generators were written in terms of the twisted currents (see Eq. (9.3)), we see in the form (LABEL:unbalancedgroup) a generic mode imbalance (when $`L^{A\alpha }`$ and $`L^{\alpha \beta }`$, $`\alpha +\beta 0`$ are non-zero) in which the modes of the currents of $`L_\sigma (m)`$ do not necessarily sum to $`m`$. It follows that the generic inner-automorphic orbifold can not be described by the stress-tensor spectral flow (9.5), in which the current modes sum to $`m`$. (The mode imbalance phenomenon is avoided however for the special case when only $`L^{A\alpha }`$ and $`L^{\alpha \beta }`$ are non-zero, so that, as discussed in the text, these special inner-automorphic orbifolds can be described by spectral flow.)
A similar mode imbalance is found for all the stress tensors of all the permutation orbifolds when the orbifold induction procedure (9.24b) is used to express these stress tensors in terms of untwisted current modes. As an example, we mention the form
$`L_\sigma (m)`$ $`=`$ $`{\displaystyle \underset{r=0}{\overset{\rho (\sigma )1}{}}}[{\displaystyle \underset{j,l=0}{\overset{\frac{\lambda }{\rho (\sigma )}1}{}}}^{raj;r,bl}(\sigma ){\displaystyle \underset{p}{}}:J_{aj}(\rho (\sigma )p+r)J_{bl}(\rho (\sigma )(mp)r):_M`$ (I.2)
$`+{\displaystyle \underset{j=0}{\overset{\frac{\lambda }{\rho (\sigma )}1}{}}}^{raj;r,bj}(\sigma )\{{\displaystyle \frac{ir}{\rho (\sigma )}}f_{ab}^cJ_{cj}(\rho (\sigma )m)+\delta _{m,0}{\displaystyle \frac{r}{2\rho (\sigma )}}(1{\displaystyle \frac{r}{\rho (\sigma )}})k\rho (\sigma )\eta _{ab}\}]`$
obtained from (3.15) for the sectors of the orbifold $`A(_\lambda )/_\lambda `$.
### Appendix J. The OVME and $`A(𝔻_\lambda )/𝔻_\lambda `$
In Ref. References it was shown that every solution of the orbifold Virasoro master equation$`^{\text{References}}`$ (OVME) at order $`\lambda `$ is a sector of a permutation orbifold of type $`A(𝔻_\lambda )/_\lambda `$. Moreover, it was asserted in Ref. References that the sectors described by the OVME also occur in the permutation orbifolds $`A(𝔻_\lambda )/𝔻_\lambda `$.
To see this we begin with (the inverse of) the first ($`\sigma =1`$) duality transformation,$`^{\text{References}}`$ which constructs a particular $`𝔻_\lambda `$-invariant CFT described by $`L`$ from any particular solution $`_{OVME}=`$ of the OVME at order $`\lambda `$:
$$L_{IJ}^{ab}=\lambda \underset{r=0}{\overset{\lambda 1}{}}_r^{ab}U(1)_{r;I}U(1)_{r;J},_r^{ab}=_r^{ba}=_r^{ba}L_K^{ab}=L_K^{ba}=L_K^{ba}$$
(J.1a)
$$U^{}(1)_{I;r}=\frac{1}{\sqrt{\lambda }}e^{\frac{2\pi irI}{\lambda }},\rho (1)=\lambda ,I,J,r=0,\mathrm{},\lambda 1.$$
(J.1b)
The solution $``$ of the OVME was identified in Ref. References as a sector of the orbifold $`A(𝔻_\lambda )/_\lambda `$ because of the $`𝔻_\lambda `$ symmetry of $`L`$ in (J.1a) and the fact that $`U^{}(1)`$ satisfies
$$\omega (h_1)_I^JU^{}(1)_{J;r}=U^{}(1)_{I;r}e^{\frac{2\pi ir}{\lambda }},\omega (h_1)_I^J=\delta _{I+1,J\text{ mod }\lambda },h_1_\lambda $$
(J.2)
which is the $`_\lambda `$-eigenvalue problem at $`\sigma =1`$.
Consider next the permutation orbifold $`A(𝔻_\lambda )/𝔻_\lambda `$, where the untwisted sector is described by the same $`𝔻_\lambda `$-symmetric $`L`$. The $`2\lambda `$ elements of the group $`𝔻_\lambda `$ are
$$\{h\}=\{r^\sigma ,sr^\sigma ;\sigma =0,\mathrm{},\lambda 1\},r^\lambda =s^2=1,rs=sr^1$$
(J.3)
and $`𝔻_\lambda `$ acts by permuting the currents according to
$$J_{aI}^{^{}}=\omega (h)_{IJ}J_{aJ},h𝔻_\lambda $$
(J.4a)
$$\omega (r^\sigma )_I^J=\delta _{I+\sigma ,J\text{ mod }\lambda },r^\sigma _\lambda 𝔻_\lambda ;\omega (sr^\sigma )_I^J=\delta _{I+\sigma ,J\text{ mod }\lambda }.$$
(J.4b)
The action of the element $`\omega (r^{\sigma =1})𝔻_\lambda `$ in (J.4b) is the same as the action of the element $`\omega (h_{\sigma =1})_\lambda `$ in (J.2). Then, the first duality transformation into the first twisted sector of $`A(𝔻_\lambda )/𝔻_\lambda `$
$$_r^{ab}=\frac{1}{\lambda }\underset{I,J=0}{\overset{\lambda 1}{}}L_{IJ}^{ab}U^{}(1)_{I;r}U^{}(1)_{J;r},_r^{ab}=_r^{ba}=_r^{ba}L_K^{ab}=L_K^{ba}=L_K^{ba}$$
(J.5)
gives the original solution $``$ of the OVME. This establishes the assertion of Ref. References.
### Appendix K. The Setup for Outer-automorphic Orbifolds
In our notation, the outer automorphism$`^{\text{References},\text{References}}`$ groups of simple $`g`$ are described by
$$\alpha _i^{}=\tau (\alpha _i)\alpha ^{}=\underset{i=1}{\overset{\text{rank}g}{}}n_i\tau (\alpha _i):\omega _\alpha ^\beta =\xi _\alpha \delta _{\tau (\alpha ),\beta },\omega _A^B=\underset{i=1}{\overset{\text{rank}g}{}}\lambda _{iA}\tau (\alpha _i)^B$$
(K.1a)
$$\xi _{\alpha _i}=1,\xi _\alpha \xi _\alpha =1,\xi _\alpha \xi _\beta \xi _\gamma ^1N_{\tau (\gamma )}(\tau (\alpha ),\tau (\beta ))=N_\gamma (\alpha ,\beta )$$
(K.1b)
where $`\alpha ,\beta ,\gamma \mathrm{\Delta }(g)`$, $`\{\lambda _i\}`$ are the fundamental weights of $`g`$ and we have taken $`\alpha ^2=2`$.
For $`SU(3)`$, the outer automorphism group is a $`_2`$
$$\rho (\sigma =1)=2,\tau (\alpha _{1,2})=\alpha _{2,1},\xi _{\pm (\alpha _1+\alpha _2)}=1$$
(K.2)
and solution of the $`H`$-eigenvalue problem gives eight twisted currents $`\widehat{J}=\{\widehat{H}_{n(r)=0,1},`$ $`\widehat{E}_{n(r)=0,1}^{\pm \alpha _i},`$ $`\widehat{E}_{n(r)=1}^{\pm (\alpha _1+\alpha _2)}\}`$ which, as expected, satisfy the outer-automorphically twisted affine Lie algebra $`A_2^{(2)}`$. The inverse inertia tensors of the outer automorphic invariant CFT’s on $`SU(3)`$ satisfy the $`H`$-invariance conditions
$$L^{\alpha _1\alpha _1}=L^{\alpha _2\alpha _2},L^{\alpha _1,\alpha _2}=L^{\alpha _2,\alpha _1}$$
(K.3a)
$$L^{\alpha _1+\alpha _2,\pm \alpha _1}=L^{\alpha _1+\alpha _2,\pm \alpha _2},L^{(\alpha _1+\alpha _2),\pm \alpha _1}=L^{(\alpha _1+\alpha _2),\pm \alpha _2}$$
(K.3b)
and corresponding conditions on $`L^{AB},L^{A\alpha }`$. The reader is encouraged to work out the outer-automorphic orbifolds in further detail.
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# Transition between Compressible and Incompressible States in Infinite-Layer Fractional Quantum Hall Systems
## I Introduction
Although the fractional quantum Hall effect is originally a two-dimensional phenomenon , multi-layer systems have been studied as natural development . Multi-layer systems show some unusual effects because of new additional degrees of freedom, namely the inter-layer Coulomb interaction and the inter-layer electron tunneling. For example, in the case of the simplest double-layer system, the ”even denominator” fractional quantum Hall effect with the Landau level filling $`\nu =1/2=1/4+1/4`$ has been observed , and is interpreted in terms of Halperin’s $`\mathrm{\Psi }_{331}`$ state .
One may then naturally ask what happens in infinite multi-layer systems. Experimentally, the integer quantum Hall effect is observed in a 30-layer GaAs/AlGaAs superlattice system and also in organic conductors $`\left(\mathrm{TMTSF}\right)_2\mathrm{X}`$ and BEDT-TTF salts . Motivated by these works we study the ground state of the infinite-layer system under a strong magnetic field and discuss quantum phase transitions between incompressible and compressible states.
We can have the following naive expectation in advance of quantitative analysis: In the limit of infinite inter-layer distance $`d_s`$, layers are completely independent of each other and the Laughlin state is realized for appropriate fillings, e.g. $`\nu =1/q`$ where $`q`$ is an odd integer. However for $`d_sl_B`$ where $`l_B`$ is the magnetic length, other fractional quantum Hall states can be realized because of the inter-layer Coulomb interaction. Furthermore three-dimensional states are also realizable with increasing electron tunnelings between layers. The fractional quantum Hall states are incompressible, but the last three-dimensional state is compressible because the Bloch states are formed perpendicular to layers and the Fermi surface should appear. Therefore a transition between incompressible and compressible states is expected at a certain value of $`d_s`$, provided that $`d_s`$ is a variable parameter.
Based on the above expectation we introduce three typical trial states for the ground state with the lowest Landau level filling $`\nu =1/3`$ and compute energies of them in order to find which one is the ground state. Under a strong magnetic field we restrict the single-particle states within the lowest Landau level, and consider only the Coulomb energy and tunneling energy. We calculate the Coulomb energy by the variational Monte Carlo (VMC) method for the two-dimensional trial states , while the Hartree-Fock approximation is used for the three-dimensional trial state.
## II Model
We consider a system with the area $`𝒮`$ of each layer and the height $`L_z`$ as shown in Fig. 1. The total number $`N_z`$ of layers is given by $`N_z=L_z/d_s`$. Then, with the total number $`N`$ of electrons in the system, we define the electron number per layer $`N_{}`$ by $`N_{}=N/N_z`$. Taking into account the layer structure of the system, we assume the positive charge distribution
$$n_b\left(𝐫\right)=\underset{l=1}{\overset{N_z}{}}\frac{N_{}}{𝒮}\delta \left(zld_s\right).$$
(1)
This positive charge density corresponds to the uniform distribution in each layer and the discrete distribution along the layer. Under a magnetic field $`𝐁=B𝐞_z`$ the Landau level degeneracy $`m_D`$ is given by $`m_D=𝒮/2\pi l_B^2`$ where $`l_B=\left(c\mathrm{}/eB\right)^{1/2}`$ is the magnetic length. The average filling factor of each layer $`\nu `$ is defined as $`\nu =N_{}/m_D`$. Under the strong magnetic field with $`\nu 1`$ we may deal with only the lowest Landau level (LLL) explicitly and neglect higher levels. We will discuss possible ground states in the subspace of the LLL in the following.
Now we write down the Hamiltonian of the system with $`N\mathrm{}`$ as
$$H=\underset{i=1}{\overset{N}{}}H_{0i}+\underset{i<j}{}\frac{e^2}{|𝐫_i𝐫_j|}\underset{i=1}{\overset{N}{}}𝑑𝐫\frac{e^2}{|𝐫_i𝐫|}n_b\left(𝐫\right)+\frac{1}{2}𝑑𝐫𝑑𝐫^{}\frac{e^2}{|𝐫𝐫^{}|}n_b\left(𝐫\right)n_b\left(𝐫^{}\right).$$
(2)
Here the one-body part is given by
$$H_{0i}=\frac{1}{2m_{}^{}}\left(𝐩_i+\frac{e}{c}𝐀_{}\left(𝐫_i\right)\right)^2+\frac{1}{2m_{}^{}}p_{zi}^2+V_0\left(z_i\right),$$
(3)
where $`m_{}^{}`$ and $`m_{}^{}`$ are the effective masses in the layer and along the $`z`$-direction respectively. The squared effective charge $`e^2`$ is defined by use of the dielectric constant $`\epsilon `$ as $`e^2=e^2/\epsilon `$. $`V_0\left(z\right)`$ is the binding potential of electrons to layers (like a multi square well potential). We choose two types of configuration as follows: one is the circular shape and the other is the rectangular shape. These shapes are chosen for convenience of calculations.
### A Circular Geometry
Let us consider a pile of infinitely large layers ($`𝒮=\mathrm{}`$) with the circular symmetry and take the symmetric gauge. For the ($`x,y`$)-space this geometry is the same as the one that Laughlin chose . We adopt for this geometry the periodic boundary condition in the $`z`$-direction. This geometry is convenient for description of two-dimensional states and also convenient for the VMC calculation . If we assume that electrons do not transfer between layers, then the one-particle state is written as
$$\psi _{lm}\left(𝐫\right)=\left(2^{m+1}\pi m!l_B^{2m+2}\right)^{1/2}\xi ^m\mathrm{exp}\left[\frac{|\xi |^2}{4l_B^2}\right]\varphi _0\left(zld_s\right),$$
(4)
where $`\xi =xiy`$ is the complex coordinate in the ($`x,y`$)-space and $`\varphi _0\left(zld_s\right)`$ defined by (15) is the normalized wave function localized at the layer $`l`$. Since the uniform positive charge is inconvenient for finite particles which we actually study numerically, we take an alternative positive charge density $`\stackrel{~}{n}_b\left(𝐫\right)`$ which has the same distribution as that of electrons. Namely with $`P\left(\left\{𝐫_i\right\}\right)=|\psi \left(\left\{𝐫_i\right\}\right)|^2`$ being the probability distribution function of electrons, we require
$$\stackrel{~}{n}_b\left(𝐫\right)=\underset{j=1}{\overset{N}{}}\left(\underset{i}{}d𝐫_i\right)\delta \left(𝐫𝐫_j\right)P\left(\left\{𝐫_i\right\}\right).$$
(5)
This positive charge distribution tends to the uniform one in the infinite-size limit.
We now discuss the Coulomb potential with periodic boundary condition along the $`z`$-axis. Taking account of mirror images we obtain the potential as
$$\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\frac{e^2}{\left[\left(𝐫_{}𝐫_{}^{}\right)^2+\left(zz^{}+nL_z\right)^2\right]^{1/2}}.$$
(6)
The summation over $`n`$ is divergent logarithmically. This divergence should be canceled by the attractive part in the final result. In this paper we replace it by the following one:
$$𝒱\left(𝐫𝐫^{}\right)=\frac{e^2}{\left[\left(𝐫_{}𝐫_{}^{}\right)^2+\left(\frac{L_z}{\pi }\mathrm{sin}\left(\frac{zz^{}}{L_z}\pi \right)\right)^2\right]^{1/2}}.$$
(7)
This potential has the right periodicity along the $`z`$-axis and becomes equivalent to the Coulomb potential in the limit of $`L_z\mathrm{}`$. The Hamiltonian with finite $`N`$ is represented as follows:
$$H=\underset{i=1}{\overset{N}{}}H_{0i}+\underset{i<j}{}𝒱\left(𝐫_i𝐫_j\right)\underset{i=1}{\overset{N}{}}𝑑𝐫𝒱\left(𝐫_i𝐫\right)\stackrel{~}{n}_b\left(𝐫\right)+\frac{1}{2}𝑑𝐫𝑑𝐫^{}𝒱\left(𝐫𝐫^{}\right)\stackrel{~}{n}_b\left(𝐫\right)\stackrel{~}{n}_b\left(𝐫^{}\right).$$
(8)
The expectation value is given by
$`H`$ $`=`$ $`N\mathrm{\Omega }+{\displaystyle \underset{i<j}{}}{\displaystyle \left(\underset{i=1}{\overset{N}{}}d𝐫_i\right)𝒱\left(𝐫_i𝐫_j\right)P\left(\left\{𝐫_i\right\}\right)}`$ (9)
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i,j}{}}{\displaystyle \left(\underset{i=1}{\overset{N}{}}d𝐫_id𝐫_i^{}\right)𝒱\left(𝐫_i𝐫_j^{}\right)P\left(\left\{𝐫_i\right\}\right)P\left(\left\{𝐫_i^{}\right\}\right)},`$ (10)
where $`\mathrm{\Omega }`$ is a constant representing the single-particle energy. It is given by
$$\mathrm{\Omega }=\frac{1}{2}\omega _c+𝑑z\varphi _0^{}\left(z\right)\left(\frac{p_z^2}{2m_{}^{}}+V_0\left(z\right)\right)\varphi _0\left(z\right),$$
(11)
where we have assumed that the tunneling along the $`z`$-direction is suppressed.
### B Rectangular Geometry
Next we consider a pile of rectangular ($`𝒮=L_x\times L_y`$) layers. We impose the periodic boundary condition in each ($`x,y`$ and $`z`$) direction and choose the Landau gauge. We obtain $`m_D=L_xL_y/2\pi l_B^2`$. In the case of three-dimensional states, this geometry is convenient. When there are tunnelings, the one electron state is written as
$$\varphi _{jk}\left(𝐫\right)=\varphi _j\left(𝐫_{}\right)\varphi _k\left(z\right),$$
(12)
where $`j`$ ($`=1,2,\mathrm{},m_D`$) indicates the degree of freedom in the LLL, and $`k=2\pi n/L_z`$ ($`n=N_z/2,N_z/2+1,\mathrm{},N_z/21`$) is the wave number along the $`z`$-direction. The $`(𝐫_{})`$-dependent part $`\varphi _j\left(𝐫_{}\right)`$ is the Landau wave function :
$$\varphi _j\left(𝐫_{}\right)=\left(\frac{1}{L_yl_B\sqrt{\pi }}\right)^{1/2}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\mathrm{exp}\left[i\kappa _{j+m_Dn}y\frac{\left(xX_{j+m_Dn}\right)^2}{2l_B^2}\right],$$
(13)
where $`\kappa _j=2\pi j/L_y`$ and $`X_j=l_B^2\kappa _j`$. For the $`z`$-dependent part $`\varphi _k\left(z\right)`$ we make the Bloch function by use of $`\varphi _0`$ as
$$\varphi _k\left(z\right)=\frac{1}{\sqrt{N_z}}\underset{n_z=\mathrm{}}{\overset{\mathrm{}}{}}\underset{n=1}{\overset{N_z}{}}e^{iknd_s}\varphi _0\left(znd_sn_zL_z\right),$$
(14)
where we take the following localized function:
$$\varphi _0\left(z\right)=\left(\frac{1}{ϵ^2d_s^2\pi }\right)^{1/4}e^{z^2/2ϵ^2d_s^2},$$
(15)
with $`ϵ`$ being positive infinitesimal. Under the tight-binding approximation $`\varphi _{jk}\left(𝐫\right)`$ becomes the eigenstate of the one-particle Hamiltonian $`H_{0i}`$. We introduce the creation and annihilation operators, $`a_{jk}^{}`$ and $`a_{jk}`$, corresponding to $`\varphi _{jk}\left(𝐫\right)`$. Then the total Hamiltonian of this system is given in second quantized form by
$``$ $`=`$ $`{\displaystyle \underset{j}{}}{\displaystyle \underset{k}{}}\left(E_k{\displaystyle \frac{N}{2V}}{\displaystyle \frac{e^2d_s^2}{\pi }}2\zeta \left(2\right)\right)a_{jk}^{}a_{jk}`$ (16)
$`+{\displaystyle \frac{1}{2V}}{\displaystyle \underset{𝐪0}{}}v\left(q\right)\left[\rho \left(𝐪\right)\rho \left(𝐪\right)\rho \left(0\right)\mathrm{exp}\left[{\displaystyle \frac{l_B^2q_{}^2}{2}}{\displaystyle \frac{ϵ^2d_s^2q_z^2}{2}}\right]\right].`$ (17)
Here $`E_k=\mathrm{\Omega }2t\mathrm{cos}kd_s`$ is the eigenvalue of the one-electron state with the inter-layer transfer $`t`$, $`\zeta \left(n\right)`$ is the Riemann’s zeta function, $`v\left(q\right)=4\pi e^2/q^2`$ is the Fourier coefficient of the Coulomb potential, and $`\rho \left(𝐪\right)`$ is the density operator defined by
$$\rho \left(𝐪\right)=𝑑𝐫\psi ^{}\left(𝐫\right)\psi \left(𝐫\right)e^{i𝐪𝐫},$$
(18)
where $`\psi \left(𝐫\right)=_{jk}a_{jk}\varphi _{jk}\left(𝐫\right)`$ is the electron field operator projected onto the LLL.
## III Trial States
We introduce three trial states for the ground states with $`\nu =1/3`$ of the Hamiltonians (8) and (17). As mentioned in the previous section, we take (8) for states without the tunneling (two-dimensional states) and (17) for states with the tunneling (three-dimensional states). Here we remark on the index of position coordinates. The position coordinates in the Hamiltonian (8) have the index $`i`$ ($`=1,2,\mathrm{},N`$) of electrons. However when we refer to two-dimensional states, we make the position coordinates have the layer index additionally because electrons belonging to different layers are distinguishable. Hence the position coordinates have the two indices $`l`$ and $`p`$ where $`l`$ ($`=0,1,\mathrm{},N_z1`$) is the layer index and $`p`$ ($`=1,2,\mathrm{},N_{}`$) is the particle index in a layer. Indices $`\{i\}`$ and $`\{l,p\}`$ are related to each other by $`i=N_{}l+p`$.
The first trial state is a two-dimensional one without the tunneling represented by
$$\mathrm{\Psi }_{\nu =1/3}^L=\underset{l=1}{\overset{N_z}{}}\left[\chi _{\nu =1/3}^L\left(\left\{\xi _{lp}\right\}\right)\right]\mathrm{\Phi }\left(\left\{\left\{z_{lp}\right\}\right\}\right),$$
(19)
where $`\chi `$ and $`\mathrm{\Phi }`$ describe the motion in the ($`x,y`$)-space and along the $`z`$-direction respectively. They are given by
$$\chi _{\nu =1/3}^L\left(\left\{\xi _{lp}\right\}\right)=\underset{p<p^{}}{}\left(\xi _{lp}\xi _{lp^{}}\right)^3\mathrm{exp}\left[\frac{1}{4l_B^2}\underset{p=1}{\overset{N_{}}{}}|\xi _{lp}|^2\right],$$
(20)
$$\mathrm{\Phi }\left(\left\{\left\{z_{lp}\right\}\right\}\right)=\underset{l=1}{\overset{N_z}{}}\underset{p=1}{\overset{N_{}}{}}\varphi _0\left(z_{lp}ld\right).$$
(21)
Here $`(\{\{z_{lp}\}\})`$ denotes $`(z_{11},\mathrm{},z_{1N_{}};\mathrm{};z_{N_z1},\mathrm{},z_{N_zN_{}})`$ and $`(\{\xi _{lp}\})`$ does $`(\xi _{l1},\xi _{l2},\mathrm{},\xi _{lN_{}})`$. We call the state (19) the ”multi-Laughlin state”. It is the direct product of the Laughlin state with $`\nu =1/3`$ in each layer and is expected to be stable for a large inter-layer distance: $`d_sl_B`$.
The second trial state is also a two-dimensional one without the tunneling. It is represented by
$$\mathrm{\Psi }_{\nu =1/3}^{\left(111\right)}=\chi _{\nu =1/3}^{\left(111\right)}\mathrm{\Phi }\left(\left\{\left\{z_{lp}\right\}\right\}\right),$$
(22)
where
$$\chi _{\nu =1/3}^{\left(111\right)}=\underset{l=1}{\overset{N_z}{}}\underset{p<p^{}}{}\left(\xi _{lp}\xi _{lp^{}}\right)\underset{p,p^{}}{}\left(\xi _{lp}\xi _{l+1p^{}}\right)\mathrm{exp}\left[\frac{1}{4l_B^2}\underset{p=1}{\overset{N_{}}{}}|\xi _{lp}|^2\right].$$
(23)
We call this state the ”(111) state” . This state is expected to be stable for $`d_sl_B`$ where the inter-layer Coulomb interaction becomes comparable to the intra-layer one. The numbers (111) indicate the exponents of the Jastrow factor which controls the strength of the electron correlation. The first ”1” is the exponent of the Jastrow factor between a given layer and the nearest lower layer, the second ”1” is that within the layer, and the last ”1” is that between the layer and the nearest upper layer. Generally we can define ($`\lambda \mu \lambda `$) state for our system. For example (131) state is possible for $`\nu =1/5`$, and (212) state is also possible for the same $`\nu `$. Physically (212) state should be unstable because the inter-layer correlation becomes larger than the intra-layer one in this state. In the case of $`\nu =1/3`$ only the (111) state is allowed within the subspace of Halperin-type states.
The third state is expected to be stable at $`d\begin{array}{c}<\hfill \\ \hfill \end{array}l_B`$ when electrons transfer between layers. We call such a state the ”itinerant state” and define it as follows. With the tunneling, Bloch waves are formed along the $`z`$-direction. Then the simplest candidate for the ground state is the one with the smallest kinetic (tunneling) energy. Namely the cosine band in the $`z`$-direction is filled up to the Fermi level , with the filling of the Landau level for each wave number being unity below the Fermi level and zero above it. We can write it by use of electron creation operators as
$$|\mathrm{itinerant}=\underset{j=1}{\overset{m_D}{}}\underset{|k|k_F}{}a_{jk}^{}|0,$$
(24)
where $`k_F`$ is the Fermi wave number given by $`k_F=\nu \pi /d_s`$. Although we assume $`\nu =1/3`$ in this paper, the itinerant state can be defined for all $`\nu [0,1]`$.
## IV Ground-State Energies
### A Multi-Laughlin State and (111) State
We use the VMC method to calculate the expectation value of the Coulomb energy numerically for the multi-Laughlin state and the (111) state. Technical details are explained in Appendix. In the case of the multi-Laughlin state, $`H`$ are reduced to the sum of the two-dimensional result simply because all layers are independent of each other, and the inter-layer Coulomb energy is exactly canceled. We obtain
$`N\epsilon _L=H_LN\mathrm{\Omega }`$ $`=`$ $`N_{}{\displaystyle \underset{p<p^{}}{}}{\displaystyle \left(\underset{p}{}d𝐫_{1p}\right)\frac{e^2}{|𝐫_{1p}𝐫_{1p^{}}|}|\chi _{\nu =1/3}^L\left(\left\{\xi _{1p}\right\}\right)|^2}`$ (25)
$`{\displaystyle \frac{1}{2}}N_{}{\displaystyle \underset{p,p^{}}{}}{\displaystyle \left(\underset{p}{}d𝐫_{1p}\right)\frac{e^2}{|𝐫_{1p}𝐫_{1p^{}}^{}|}|\chi _{\nu =1/3}^L\left(\left\{\xi _{1p}\right\}\right)|^2|\chi _{\nu =1/3}^L\left(\left\{\xi _{1p}^{}\right\}\right)|^2}.`$ (26)
On the other hand such cancellations do not occur for the (111) state, and the calculation becomes more complicated because of the inter-layer correlation. The energy is given by
$`N\epsilon _{\left(111\right)}`$ $`=`$ $`H_{\left(111\right)}N\mathrm{\Omega }`$ (27)
$`=`$ $`\left({\displaystyle \underset{l=l^{}}{}}{\displaystyle \underset{p<p^{}}{}}+{\displaystyle \underset{l<l^{}}{}}{\displaystyle \underset{p,p^{}}{}}\right){\displaystyle \left(\underset{l}{}\underset{p}{}d𝐫_{lp}\right)𝒱\left(𝐫_{lp}𝐫_{l^{}p^{}}\right)|\mathrm{\Psi }_{\nu =1/3}^{\left(111\right)}\left(\left\{\left\{𝐫_{lp}\right\}\right\}\right)|^2}`$ (28)
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{l,l^{}}{}}{\displaystyle \underset{p,p^{}}{}}{\displaystyle \left(\underset{l}{}\underset{p}{}d𝐫_{lp}d𝐫_{lp}^{}\right)𝒱\left(𝐫_{lp}𝐫_{l^{}p^{}}^{}\right)|\mathrm{\Psi }_{\nu =1/3}^{\left(111\right)}(\{\{𝐫_{lp}\}\})|^2|\mathrm{\Psi }_{\nu =1/3}^{\left(111\right)}(\{\{𝐫_{lp}^{}\}\})|^2}.`$ (29)
Figure 2 shows the size dependence of the Coulomb energy of the multi-Laughlin state. The extrapolated value is in good agreement with that in the literature . The dependence of the Coulomb energy for the (111) state on the inter-layer distance is shown in Fig.3. The calculation was performed for three cases of finite size systems: $`N_{}=20`$, $`30`$ and $`42`$ with $`N_z=5`$. We find that the energy decreases with decreasing inter-layer distance. The dependence on $`N_{}`$ is small in the region of small inter-layer distance.
### B Itinerant State
We calculate the energy of the itinerant state for the Hamiltonian (17) with use of the Hartree-Fock approximation. Since we take the homogeneous trial state, the mean field is given by
$`\mathrm{\Delta }(jk;j^{}k^{})`$ $`=`$ $`a_{jk}^{}a_{j^{}k^{}}=\mathrm{itinerant}|a_{jk}^{}a_{j^{}k^{}}|\mathrm{itinerant}`$ (30)
$`=`$ $`\delta _{jj^{}}\delta _{kk^{}}\theta \left(k_F|k|\right).`$ (31)
The mean field for the interaction part of the Hamiltonian (17) consists of the direct and the exchange interactions. Of these the energy due to the direct term is written as
$`{\displaystyle \frac{1}{2V}}{\displaystyle \underset{𝐪0}{}}v\left(q\right)\rho \left(𝐪\right)\rho \left(𝐪\right)`$ $`=`$ $`{\displaystyle \frac{N^2}{2V}}{\displaystyle \underset{𝐪0}{}}v\left(q\right)\delta _{𝐪_{},𝐠}\delta _{q_z,h}`$ (32)
$`=`$ $`{\displaystyle \frac{N^2}{2V}}{\displaystyle \frac{e^2d_s^2}{\pi }}2\zeta \left(2\right),`$ (33)
where we used the relation $`\rho \left(𝐪\right)=N\delta _{𝐪_{},𝐠}\delta _{q_z,h}`$ valid in the macroscopic limit. Here we have $`𝐠=(\frac{2\pi }{L_x}m_Dn_x,\frac{2\pi }{L_y}m_Dn_y)`$ and $`h=\frac{2\pi }{d_s}n_z`$ with $`n_x,n_y,n_z`$ integers. The direct term cancels exactly with the Coulomb energy of the positive charge because of homogeneity of the itinerant state within each layer. Therefore the Hartree-Fock Hamiltonian is given by
$`^{HF}`$ $`=`$ $`{\displaystyle \underset{j}{}}{\displaystyle \underset{k}{}}E_ka_{jk}^{}a_{jk}`$ (34)
$`{\displaystyle \frac{1}{2V}}{\displaystyle \underset{j_1,j_2}{}}{\displaystyle \underset{k_1,k_2}{}}a_{j_1k_1}^{}a_{j_1k_1}\theta \left(k_F|k_2|\right){\displaystyle \underset{𝐪0}{}}v\left(q\right)M(j_1k_1,j_2k_2;𝐪)M(j_2k_2,j_1k_1;𝐪),`$ (35)
where the matrix element for the exchange mean-field is given by
$`M(jk,j^{}k^{};𝐪)`$ $`=`$ $`{\displaystyle 𝑑𝐫\mathrm{\Phi }_{jk}^{}\left(𝐫\right)\mathrm{\Phi }_{j^{}k^{}}\left(𝐫\right)e^{i𝐪𝐫}}`$ (36)
$`=`$ $`{\displaystyle \underset{h}{}}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\delta _{k^{},kq_z+h}\delta _{jj^{}+m_Dn,\frac{L_y}{2\pi }q_y}\mathrm{exp}\left[{\displaystyle \frac{l_B^2q_{}^2}{4}}{\displaystyle \frac{ϵ^2d_s^2q_z^2}{4}}+iq_xX_j\right].`$ (37)
The energy of the itinerant state is obtained as
$`E`$ $`=`$ $`m_D{\displaystyle \underset{|k|k_F}{}}E_k`$ (38)
$`{\displaystyle \frac{m_D}{2V}}{\displaystyle \underset{𝐪0}{}}v\left(q\right){\displaystyle \underset{k}{}}{\displaystyle \underset{h}{}}\theta \left(k_F|k|\right)\theta \left(k_F|kq_z+h|\right)\mathrm{exp}\left[{\displaystyle \frac{l_B^2q_{}^2}{2}}{\displaystyle \frac{ϵ^2d_s^2q_z^2}{2}}\right].`$ (39)
The region of summation over $`k`$ in the above formula is shown in Fig. 4. In order to derive the sum $`S`$ we analyze the cases as follows: (1) $`S=0`$ for $`q_zh>2k_F`$, (2) $`S=\frac{L_z}{2\pi }(k_F(q_zhk_F))`$ for $`0<q_zh2k_F`$, (3) $`S=\frac{L_z}{2\pi }(q_zh+k_F(k_F))`$ for $`2k_F<q_zh0`$, and (4) $`S=0`$ for $`q_zhk_F`$. As a result we summarize the four cases into a single expression given by
$$S=\frac{L_z}{2\pi }2k_F\left(1\frac{|q_zh|}{2k_F}\right)\theta \left(2k_F|q_zh|\right).$$
(40)
Therefore the energy per particle which contains the kinetic term is given in the thermodynamic limit by
$`\epsilon ^{HF}`$ $`=`$ $`\alpha {\displaystyle \frac{2t}{\nu \pi }}\mathrm{sin}\left(\nu \pi \right)`$ (41)
$`{\displaystyle \frac{1}{2\left(2\pi \right)^3}}{\displaystyle \underset{h}{}}{\displaystyle _{0+}^{\mathrm{}}}2\pi q_{}𝑑q_{}{\displaystyle _{h2k_F}^{h+2k_F}}𝑑q_zv\left(q\right)\left(1{\displaystyle \frac{|q_zh|}{2k_F}}\right)\mathrm{exp}\left[{\displaystyle \frac{l_B^2}{2}}q_{}^2{\displaystyle \frac{ϵ^2d_s^2}{2}}q_z^2\right]`$ (42)
$`=`$ $`\alpha +\epsilon _{kin}+\epsilon _{int}^{HF}.`$ (43)
At this stage we let $`ϵ0`$. It can be shown that the contribution of large $`q_z`$ ($`1/ϵd_s`$) to the integration is negligible. To see this we observe that $`v\left(q\right)=1/\left(q_{}^2+q_z^2\right)`$ becomes $`𝒪\left(ϵ^2\right)`$ and summation over $`h`$ converges in the limit of $`ϵ0`$. Therefore, $`\epsilon _{int}^{HF}`$ becomes
$$\epsilon _{int}^{HF}=\frac{1}{2\left(2\pi \right)^2}_{0+}^{\mathrm{}}q_{}𝑑q_{}_{2k_F}^{2k_F}𝑑q_z\underset{h}{}\frac{4\pi e^2}{q_{}^2+\left(q_z+h\right)^2}\left(1\frac{|q_z|}{2k_F}\right)\mathrm{exp}\left(\frac{l_B^2}{2}q_{}^2\right).$$
(44)
With use of the following formula:
$$\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\frac{1}{x^2+\left(y+2\pi n\right)^2}=\frac{\mathrm{sinh}x}{2x\left(\mathrm{cosh}x\mathrm{cos}y\right)},$$
(45)
we can accomplish the summation over $`h`$ as
$$\underset{h}{}\frac{1}{q_{}^2+\left(q_z+h\right)^2}=\frac{d_s^2\mathrm{sinh}d_sq_{}}{2d_sq_{}\left(\mathrm{cosh}d_sq_{}\mathrm{cos}d_sq_z\right)}.$$
(46)
Consequently we obtain the following result:
$$\epsilon _{int}^{HF}=\frac{e^2}{l_B}\frac{1}{2\pi }\frac{1}{\overline{d}_s}_0^{\mathrm{}}𝑑s_0^{2\pi /3}𝑑t\left(1\frac{t}{2\pi \nu }\right)\frac{\mathrm{sinh}s}{\mathrm{cosh}s\mathrm{cos}t}\mathrm{exp}\left(\frac{s^2}{2\overline{d}_s^2}\right).$$
(47)
where $`\overline{d}_s=d_s/l_B`$ denotes the inter-layer distance in units of the magnetic length. Figure 5 shows $`\epsilon _{int}^{HF}`$ as a function of the inter-layer distance. We remark that $`\epsilon _{int}^{HF}`$ decreases with decreasing $`d_s`$. This result originates from the fact that the inter-layer Coulomb interaction becomes more effective as layers get closer.
### C Comparison of Ground-State Energies
We first compare the Coulomb energy per particle of the three states. Figure 6 shows the results. The energy of the (111) state is obtained by the VMC calculation for 5 layers with 42 particles in each layer. For the multi-Laughlin and the itinerant states the result is shown for the infinite-size system. We remark the energy of the multi-Laughlin state does not depend on the inter-layer distance since all layers are independent of each other and electrons feel other layers neutral. In the region of $`d_s\begin{array}{c}>\hfill \\ \hfill \end{array}l_B`$ the Coulomb energy of the multi-Laughlin state is the lowest. This is evident because the system becomes a pile of almost independent two-dimensional layers in such a region. In the region of $`d_s\begin{array}{c}<\hfill \\ \hfill \end{array}l_B`$ the (111) state has the lowest Coulomb energy. This result originates from the fact that the (111) state gains the inter-layer Coulomb energy at the cost of the intra-layer one in the region where intra- and inter-layer interactions are comparable. With respect to the itinerant state, the Coulomb energy decreases in the region of $`d_s\begin{array}{c}<\hfill \\ \hfill \end{array}l_B`$. There is however no region where it is the lowest.
So far we have considered only the Coulomb energy. We must take into account the tunneling energy of the itinerant state in order to compare the total energies of the three states. Let us assume that the transfer energy $`t`$ depends on the inter-layer distance as
$$t=A\mathrm{exp}\left(\alpha d_s\right),$$
(48)
where $`A`$ and $`\alpha `$ are constants. Then as shown in Fig. 7 we find a region of $`d_s`$ where the itinerant state has the lowest total energy. Here we put appropriate values into $`A`$ and $`\alpha `$, namely $`A=0.2e^2/l_B,\alpha =1.0/l_B`$ for a case shown as $`t_1`$, and $`A=0.3e^2/l_B,\alpha =0.6795/l_B`$ for another case shown as $`t_2`$. The parameters for $`t_2`$ with $`B=10T`$ are chosen to reproduce the band width of $`4t=2.5m\mathrm{eV}=0.179e^2/l_B`$ at $`d_s=0.226\AA =2.8l_B`$ which seems appropriate for a GaAs/AlGaAs superlattice .
On the basis of the comparison we propose a phase diagram as shown in Fig. 8 for the ground state of the present system in the plane of $`t`$ and $`d_s`$. If we assume the relation (48) between $`t`$ and $`d_s`$ in the diagram, possible states follow a curve in Fig. 8 as $`d_s`$ is varied. Depending on the parameters $`A`$ and $`\alpha `$, there are either single transition ($`t_1`$) or double transitions ($`t_2`$). Consequently the phase diagram Fig. 8 is interpreted as follows.
1. In the region $`d_sl_B`$, the system is a pile of independent two-dimensional layers and the incompressible multi-Laughlin state is stable.
2. In the region of $`d_s\begin{array}{c}<\hfill \\ \hfill \end{array}l_B`$ and the tunneling energy is small ($`t\begin{array}{c}<\hfill \\ \hfill \end{array}0.1e^2/l_B`$), the (111) state which is also incompressible is stabilized instead of the multi-Laughlin state by the inter-layer correlation.
3. When the tunneling energy is large, the compressible itinerant state is stabilized by competition of the tunneling energy and the Coulomb energy.
## V Discussion
We compared the result of the (111) state for a finite size system and that of the multi-Laughlin and the itinerant states for the infinite size system. The finite size ($`N_{}`$ and $`N_z`$) correction should be taken into account in the (111) state for more quantitative discussion. While we have no information with respect to the dependence on $`N_z`$, the following remark is given from Fig. 3 about the dependence on $`N_{}`$. When one increases $`N_{}`$ the Coulomb energy much increases in the region of $`d_sl_B`$. While for $`d_sl_B`$ the Coulomb energy looks almost independent of $`N_{}`$ within the error bars. Since the boundaries of the (111) state occurs with $`d_s/l_B<2`$ in Fig. 8, the finite-size correction should not change the results qualitatively.
It should be possible to observe transitions in actual systems by adding a pressure to the system or tilting the magnetic field. Since adding a pressure corresponds to control of $`d_s`$, a change of parameters along a curve like $`t_1`$ or $`t_2`$ in the diagram Fig. 8 is expected. On the other hand tilting the magnetic field corresponds to control of $`t`$ independent of $`d_s`$ because the horizontal component of the magnetic field works on electrons to localize in the vertical direction and reduces $`t`$. Then the parameter change along a vertical straight line at a certain $`d_s`$ in Fig. 8 is expected.
The transition between incompressible and compressible states is analogous to the Mott transition in the Hubbard model. Namely the incompressible multi-Laughlin and (111) states correspond to an insulator, while the compressible itinerant state corresponds to a metal. The transition between incompressible and compressible states occurs due to competition of the transfer (tunneling) and the on site (intra-layer) Coulomb energy. In our model we estimate that the critical transfer energy is almost $`0.1e^2/l_B`$.
Although we discussed only three trial states for $`\nu =1/3`$ in the present study, we cannot exclude the possibility of other states. In terms of two-dimensional incompressible states for $`\nu =1/3`$, it is difficult to construct states more stable than the multi-Laughlin and the (111) states within the simple family of Jastrow-type functions. One may however ask whether other states with three-body correlation or even more correlations are more favorable energetically. For three-dimensional compressible states, the itinerant state is taken to be the simplest homogeneous state. Then inhomogeneous states like CDW should also be discussed for more refined argument. Furthermore we have neglected the spin degree of freedom which may be relevant to some systems. For example the SDW state, which accounts for the QHE of $`\left(\mathrm{TMTSF}\right)_2\mathrm{X}`$ , can be a relevant state depending on parameters in the system.
## Acknowledgments
One of the authors (S. S) would like to thank H. Yokoyama and S. Tokizaki for instruction of the VMC method and useful discussion, and M. Arikawa for pertinent suggestions.
## VMC Method for the (111) State
We describe the method of numerical calculation for the (111) state in this appendix. The expectation value of the Coulomb energy for the (111) state is represented by
$`N\epsilon _{\left(111\right)}`$ $`=`$ $`{\displaystyle \left(\underset{l}{}\underset{p}{}d𝐫_{lp}\right)\left(\underset{l=l^{}}{}\underset{p<p^{}}{}+\underset{l<l^{}}{}\underset{p,p^{}}{}\right)𝒱\left(𝐫_{lp}𝐫_{l^{}p^{}}\right)|\mathrm{\Psi }_{\nu =1/3}^{\left(111\right)}\left(\left\{\left\{𝐫_{lp}\right\}\right\}\right)|^2}`$ (49)
$`{\displaystyle \frac{1}{2}}{\displaystyle \left(\underset{l}{}\underset{p}{}d𝐫_{lp}d𝐫_{lp}^{}\right)\underset{l,l^{}}{}\underset{p,p^{}}{}𝒱\left(𝐫_{lp}𝐫_{l^{}p^{}}^{}\right)|\mathrm{\Psi }_{\nu =1/3}^{\left(111\right)}\left(\left\{\left\{𝐫_{lp}\right\}\right\}\right)|^2|\mathrm{\Psi }_{\nu =1/3}^{\left(111\right)}\left(\{\{𝐫_{lp}^{}\}\}\right)|^2}`$ (50)
$`=`$ $`V_1+V_2.`$ (51)
We can obtain the numerical value of $`\epsilon _{\left(111\right)}`$ for finite size (finite particles) systems by the VMC method. In evaluating the second term with minimum statistical errors, we separate the variables in it and reduce the expectation value of the quantity between two bodies to that of one body. For this purpose we introduce the Fourier transform of the Coulomb potential $`\stackrel{~}{𝒱}`$, one body distribution function $`f`$, and its Fourier transform $`g`$ as follows:
$$𝒱\left(𝐫𝐫^{}\right)=\frac{1}{\left(2\pi \right)^2}𝑑𝐪_{}\stackrel{~}{𝒱}_{\overline{z}\overline{z}^{}}\left(q_{}\right)e^{i𝐪_{}\left(𝐫_{}𝐫_{}^{}\right)},$$
(52)
$$\stackrel{~}{𝒱}_{\overline{z}\overline{z}^{}}\left(q_{}\right)=\frac{4\pi e^2}{q_{}}\mathrm{exp}\left[q_{}\frac{L_z}{\pi }\mathrm{sin}\left(\frac{|\overline{z}\overline{z}^{}|}{N_z}\pi \right)\right],$$
(53)
$$f\left(r_{}\right)=\left(\underset{l}{}\underset{p}{}d𝐫_{lp}\right)\delta \left(𝐫_{lp}𝐫_{}\right)|\mathrm{\Psi }_{\nu =1/3}^{\left(111\right)}|^2,$$
(54)
$`g\left(q_{}\right)`$ $`=`$ $`{\displaystyle 𝑑𝐫_{}f\left(r_{}\right)e^{i𝐪_{}𝐫_{}}}`$ (55)
$`=`$ $`{\displaystyle \left(\underset{l}{}\underset{p}{}d𝐫_{lp}\right)e^{i𝐪_{}𝐫_{lp}}|\mathrm{\Psi }_{\nu =1/3}^{\left(111\right)}|^2},`$ (56)
where we wrote $`z/d_s`$ as $`\overline{z}`$. Then the second term $`V_2`$ in (51) is arranged into
$`V_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{\left(2\pi \right)^2}}{\displaystyle _0^{\mathrm{}}}N_{}^2{\displaystyle \underset{l,l^{}}{}}2\pi q_{}dq_{}\stackrel{~}{𝒱}_{ll^{}}\left(q_{}\right)g\left(q_{}\right)^2`$ (57)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{NN_{}}{\left(2\pi \right)^2}}{\displaystyle \underset{l=1}{\overset{N_z}{}}}{\displaystyle _0^{\mathrm{}}}2\pi q_{}𝑑q_{}\stackrel{~}{𝒱}_{1l}\left(q_{}\right)g\left(q_{}\right)^2.`$ (58)
We use this formula to obtain the numerical value of $`V_2`$. Figure 9 shows an example of numerical results. The integral with respect to $`q_{}`$ is performed by the trapezoidal rule because only discrete values of $`g\left(q_{}\right)`$ are obtained by the VMC calculation. We cut off the upper region in the integral at a finite value $`q_c(1/l_B)`$ because the integrand converges to zero near $`q_{}q_c`$ and its contribution to the integral is negligible in the region $`q_c<q_{}<\mathrm{}`$.
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# SPHERICALLY SYMMETRIC SIMULATION WITH BOLTZMANN NEUTRINO TRANSPORT OF CORE COLLAPSE AND POST-BOUNCE EVOLUTION OF A 15 M⊙ STAR
## 1. INTRODUCTION
The mechanism of supernova explosions of massive stars is still not satisfactorily understood. Detailed numerical models showed that the hydrodynamic shock, which is launched when the collapsing stellar core bounces abruptly by the stiffening of the equation of state (EoS) at nuclear densities, cannot propagate out promptly but stalls because of energy losses due to photodisintegration of iron-group nuclei and neutrino emission from the shock-heated matter (e.g., Bruenn 1985, 1989a,b; Myra et al. 1987). Early suggestions that energy deposition by neutrinos might cause an explosion reach back to Colgate & White (1966). The modern version of the neutrino-driven “delayed” explosion mechanism is due to Wilson (1985), who found that neutrino energy deposition can revive the stalled shock on a time scale of several hundred milliseconds after bounce (Bethe & Wilson 1985). Because of the complexity of the involved physics and the low efficiency of the neutrino energy transfer it remained unclear for years whether the explosions are sufficiently energetic and whether the delayed mechanism works for a larger range of stellar masses (Wilson et al. 1986; Bruenn 1993). Recognizing that neutron-finger convection in the newly formed neutron star could increase the neutrino luminosities, Wilson & Mayle (1988, 1993) managed to obtain healthy explosions. However, the question of neutron star convection is not finally settled and currently it is not clear whether neutron-finger instabilities or Ledoux convection (Burrows 1987, Pons et al. 1999) or quasi-Ledoux convection (Keil, Janka, & Müller 1996; Janka & Keil 1998) or none (Bruenn et al. 1995; Bruenn & Dineva 1996; Mezzacappa et al. 1998a) occur and how they affect the explosion.
Multi-dimensional hydrodynamic models (Herant et al. 1994; Miller, Wilson, & Mayle 1993; Burrows, Hayes, & Fryxell 1995; Janka & Müller 1996; Mezzacappa et al. 1998b) have demonstrated the existence and the importance of convective overturn in the neutrino-heating layer behind the supernova shock. Driven by a negative entropy gradient which emerges behind the weakening prompt shock and is enhanced by the neutrino energy deposition, the convective motions transport energy from the region of strongest heating to the shock, thus raising the postshock pressure and pushing the shock farther out. At the same time, cold, low-entropy matter is advected downward where it can readily absorb energy from the upstreaming neutrinos. These hydrodynamic instabilities have a bearing on the measured kick velocities of pulsars (Lyne & Lorimer 1994, Cordes & Chernoff 1998) and the anisotropies observed in many supernovae. They are essential to understand the production of radioactive elements in the vicinity of the nascent neutron star and their large-scale mixing into the hydrogen and helium layers of the exploding star (Kifonidis et al. 2000).
All multi-dimensional simulations have so far been carried out with serious simplifications of the neutrino transport. Even the most advanced spherically symmetric post-bounce models have only employed MGFLD (Bruenn 1993, Bruenn et al. 1995) until recently. The significance of an accurate neutrino transport for the delayed explosion mechanism, however, has long been recognized (Janka 1991, Messer et al. 1998, Yamada, Janka, & Suzuki 1999, Burrows et al. 2000). It is therefore a natural step that a new generation of supernova models will employ schemes based on a solution of the Boltzmann equation. In fact, Mezzacappa et al. (2000) have published results for a 13 M star which show that a better transport can make a qualitative change to the outcome of the simulations. However, they considered a model with an exceptionally small iron core of 1.17 M (Nomoto & Hashimoto 1988) and the explosion energy was only $`0.41\times 10^{51}`$erg at a post-bounce time of $`550`$ms. The growth rate of this energy of $`0.05\times 10^{51}`$erg per 100 ms cannot easily be extrapolated in time and will probably not increase the explosion energy significantly, because the density around the mass cut drops rapidly and the heating region is evacuated by the developing bifurcation between neutron star and ejecta.
In this Letter we present results for a Newtonian simulation of a 15 M star with a 1.28 M iron core (Woosley & Weaver 1995) which show that an accurate neutrino transport does not produce an explosion for this star in spherical symmetry.
## 2. NUMERICAL METHODS
We have developed a new transport code which determines the neutrino phase-space distribution by iteratively solving the radiation moment equations for neutrino energy and momentum coupled to the Boltzmann equation. The code takes into account effects due to the motion of the stellar medium to order $`O(v/c)`$ and determines the neutrino quantities in a comoving frame of reference (Rampp 2000; Rampp & Janka 2000, in preparation). It allows a general relativistic treatment, but for comparison with published results we have restricted ourselves to the Newtonian case. The angle dependence of the distribution function is accounted for by the use of a grid of tangent rays which exploits spherical symmetry. Closure of the set of moment equations is achieved by a variable Eddington factor calculated from the solution of the Boltzmann equation, and the integro-differential character of the latter is tamed by making use of the integral moments of the neutrino distribution as obtained from the moment equations. The method is similar to the one described by Burrows et al. (2000). In order to fulfill lepton number conservation, we employ additional moment equations for neutrino number density and number flux. Severe time step restrictions are avoided and proper establishment of equilibrium is ensured by integrating the set of transport equations implicitly in time. The stiff character of the source terms for neutrino energy and lepton number requires a simultaneous implicit update of the temperature and electron fraction of the stellar medium.
The transport is coupled to the hydrodynamics code Prometheus, which integrates the continuity equations for mass, momentum, energy and particle species in a conservative way on a moving radial grid by explicit time stepping. The integration is accurate to second order in space and time. Shocks are treated as local Riemann problems at the zone interfaces (Fryxell, Müller & Arnett 1989). The source terms for energy and momentum due to gravity and neutrinos, and for lepton number due to neutrino emission/absorption are handled by an operator-splitting technique. The stellar background and the neutrinos are evolved on different radial grids and with different time steps, which are constrained by changes per transport step (which is typically larger than the hydrodynamical step) of at most 10% for the neutrino quantities and 5% for the fluid quantities. Interpolation between both grids is done in a conservative manner.
We used the EoS of Lattimer & Swesty (1991) (with nuclear incompressibility modulus of $`K=180`$ MeV) which is extended to densities and temperatures below the regime of nuclear statistical equilibrium by an ideal gas equation of state, corrected for Coulomb-lattice effects, that includes arbitrarily relativistic and degenerate electrons and positrons, photons, and a mixture of predefined nuclear species. Nuclear burning was not taken into account in the present simulation.
The hydrodynamics was solved on a grid with 400 radial zones out to 20000 km, which were moved with the matter of the iron core during collapse to ensure good spatial resolution at all times, and kept fixed later. For the transport we used a Eulerian grid with 210 geometrically spaced radial zones, 230 tangent rays and 27 energy bins geometrically distributed between 0 and 380 MeV, the zone center of the first zone being at 1 MeV. The quality of the energy conservation limits the error in the net energy deposition by neutrinos to $`<5\times 10^{49}`$erg, and lepton number is globally conserved to better than 0.1%.
The present simulation includes only $`\nu _e`$ and $`\overline{\nu }_e`$. The corresponding rates for charged-current and neutral-current reactions with nucleons and nuclei and for neutrino-electron scattering were taken from Bruenn (1985), Mezzacappa & Bruenn (1993) and Bruenn & Mezzacappa (1997). We neglect production and annihilation of $`\nu _e\overline{\nu }_e`$ pairs, which are of minor importance compared to the charged-current reactions with nucleons. A detailed comparison of core-collapse results with published models of Bruenn & Mezzacappa (1997) showed excellent agreement. Disregarding muon and tau neutrinos and antineutrinos and $`\nu _e\overline{\nu }_e`$ pair processes has virtually no effect on the neutrino heating (see Bruenn 1993).
## 3. RESULTS
Figure 1 shows the trajectories of selected mass shells as a function of time. The bounce shock forms 211.6 ms after the onset of the collapse at a radius of 12.5 km with an enclosed mass of $`0.62`$ M. The central density at this time is $`\rho _\mathrm{c}=3.3\times 10^{14}`$g cm<sup>-3</sup> (cf. Bruenn & Mezzacappa 1997). By the rapid accretion of mass (Fig. 2) the shock is pushed out to $`240`$km. When the accretion rate drops significantly at $`120`$ms after bounce, neutrino heating is able to support further shock expansion to a radius of 350 km. After some time, however, the shock retreats again and finally turns into a standing accretion shock around 250 km, still within the collapsing silicon shell of the progenitor star. No indication for the possibility of an explosion was visible when the simulation was terminated at 350 ms after bounce. At this time the shock was stagnant and enclosed a mass of 1.5 M with increasingly negative postshock velocities. The decreasing density in the neutrino-heating region and the decay of the $`\nu _e`$ and $`\overline{\nu }_e`$ luminosities do not give hope for a later rejuvenation of the shock. The overall evolution in our simulation is very similar to model WPE15ls(180)Newt20 in Bruenn et al. (1995), who used MGFLD for the neutrino transport. The most obvious difference is a larger maximum radius of the shock, 350 km compared to only $`240`$km in the calculation by Bruenn et al. (1995). Also, the shock is able to stay near its maximum radius for a longer time and afterwards does not recede as far as in model WPE15ls(180)Newt20.
The $`\nu _e`$ and $`\overline{\nu }_e`$ luminosities and the root-mean-squared energies at 1000 km are shown as functions of time in Fig. 2. The prompt $`\nu _e`$ burst with a peak luminosity of $`3.36\times 10^{53}`$erg s<sup>-1</sup> arrives at this radius only $`6`$ms after core bounce. About 50 ms after bounce the $`\nu _e`$ and $`\overline{\nu }_e`$ luminosities have become roughly equal with a fairly stable value of (2.5–$`3)\times 10^{52}`$erg s<sup>-1</sup>. By the end of our simulation they begin to decrease slowly, different from the mean energies, which show a gradual rise to 11.2 MeV for $`\nu _e`$ and 15.5 MeV for $`\overline{\nu }_e`$.
In Fig. 3 we present profiles of the net energy deposition rate by $`\nu _e`$ and $`\overline{\nu }_e`$, of the electron fraction $`Y_e`$ and of the entropy per baryon for times 330 ms, 380 ms and 561 ms. The gain radius, below which net neutrino cooling and above which net heating occurs, is at 120–140 km (see also Fig. 1). The heating rates peak somewhat outside the gain radius and reach up to $`120`$ MeV s<sup>-1</sup> per baryon. The cooling rate below the gain radius can exceed 200 MeV s<sup>-1</sup> per baryon at late times. Maximum entropies around 13 $`k_\mathrm{B}`$ per nucleon are seen at the end of the simulation, when the density behind the shock is lowest because of the decreasing mass accretion rate. The negative entropy gradient implies potential instability against convective overturn in the region between maximum heating and supernova shock. In this layer $`Y_e`$ climbs to values larger than 0.5 and also develops a negative gradient. Values $`Y_e>0.5`$ were also found by Mezzacappa et al. (2000) in the neutrino-heated ejecta behind the outgoing shock for the successful explosion of a 13 M star. The neutronization of the neutrino-heated medium is determined by the absorption of $`\nu _e`$ on neutrons and of $`\overline{\nu }_e`$ on protons and the inverse processes. It is sensitive to the luminosities and spectra but also to the angular distributions of the neutrinos in the heating region, which govern the efficiency of energy deposition as well as the lepton exchange with the medium. Since $`\nu _e`$ decouple at a larger radius than $`\overline{\nu }_e`$, their distribution is more isotropic in the heating region, leading to a higher probability of $`\nu _e`$ absorption and thus to an increase of $`Y_e`$. This is enhanced by the recombination of $`\alpha `$ particles (Fuller & Meyer 1995).
## 4. CONCLUSIONS
Our spherically symmetric, Newtonian simulation of a 15 M star with a 1.28 M iron core, using a new Boltzmann solver for the neutrino transport, did not give an explosion until 350 ms after core bounce, although the shock reached a larger maximum radius than in a comparable MGFLD simulation of Bruenn et al. (1995). This is probably explained by stronger neutrino heating of the postshock medium with the more accurate Boltzmann transport. Since both simulations were done with the same progenitor, EoS, and neutrino opacities, and excellent agreement during the core-collapse phase was found, uncertainties due to the different numerics seem to be minimized. Although we have only included $`\nu _e`$ and $`\overline{\nu }_e`$ in our simulation, we consider our conclusions as solid, because muon and tau neutrinos would drain energy from the $`\nu _e`$ and $`\overline{\nu }_e`$ luminosities but contribute to the postshock heating only at an insignificant level due to the lack of charged-current interactions. The main effect of adding pair processes would be a weakening of the early shock propagation by additional energy losses. Also general relativity would probably hamper an explosion (Fryer 1999), but the situation is still ambiguous (De Nisco, Bruenn, & Mezzacappa 1998; Baron 1988).
The importance of an accurate $`\nu _e`$ and $`\overline{\nu }_e`$ transport is emphasized by the finding that $`Y_e>0.5`$ in the region of net neutrino energy deposition. This is interesting because $`Y_e0.48`$ was obtained in the neutrino-heated ejecta in supernova models, e.g., by Herant et al. (1994), Burrows et al. (1995) and Janka & Müller (1996), causing a large overproduction of neutron-rich nuclei around $`N=50`$ and $`A90`$ (<sup>88</sup>Sr, <sup>89</sup>Y, <sup>90</sup>Zr). This is in conflict with measured galactic abundances. With values $`Y_e>0.5`$ this problem disappears (Hoffman et al. 1996).
Using their Boltzmann solver for the neutrino transport, Mezzacappa et al. (2000) obtained a successful, but weak explosion in case of a 13 M progenitor with an extraordinarily small iron core of 1.17 M. For a 15 M star with a larger core (and therefore most likely also for more massive progenitors), we cannot confirm a qualitative difference from spherically symmetric simulations with MGFLD transport, although we find important quantitative differences with our more accurate neutrino transport. In order to obtain explosions via the neutrino-heating mechanism, multi-dimensional simulations seem indispensable for stars with typical iron core masses. Convection inside the neutron star (Keil et al. 1996) or lower neutrino opacities — due to suppression relative to the standard description by nucleon correlation effects (e.g., Janka et al. 1996, Burrows & Sawyer 1998, Reddy et al. 1999) — could raise the neutrino emission significantly on the relevant time scale of a few 100 ms after bounce, and convective overturn in the postshock region has been shown by several groups to support the explosion.
Support by Deutsche Forschungsgemeinschaft grant SFB 375 für Astro-Teilchenphysik is acknowledged. The NEC SX-5/3C of the Rechenzentrum Garching was used for the computations.
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# Soliton Gauge States and T-duality of Closed Bosonic String Compatified on Torus
## I Introduction
String duality has been the subject of active research for the last few years. The five consistent perturbative string theories are now known to be related to each other through various duality symmetries. It is believed that they are merely different moduli points of a single underlying theory termed M-theory. The best known string duality is the T-duality which can be understood perturbatively . T-duality relates a string theory in a background with large volume to another string theory in a background with small volume. For example, it has been shown that the Heterotic $`E_8E_8`$ and $`SO(32)`$ theories sit at different points. which are T-dual to each other, of the moduli space of the same Heterotic theory below ten dimension .
For the compatified bosonic string, the discrete T-duality group were shown to be the residual Weyl subgroup of the enhanced Kac-Moody gauge symmetry . On the other hand, it has been known that space-time gauge symmetry of uncompatified string is related to the existence of gauge states in the spectrum . For the 10D Heterotic string, the Heterotic gauge states are responsible for the massless $`E_8E_8`$ or $`SO(32)`$ gauge symmetry and are used to predict the existence of an infinite number of massive Einstein-Yang-Mills type gauge symmetry. For the toy 2D string, the discrete gauge states are responsible for the $`w_{\mathrm{}}`$ symmetry of the Liouville theory. It is thus of interest to understand the gauge state structure of the compatified string theory, and study their relation to the enhanced Kac-Moody gauge symmetry.
In this paper, for simplicity, we will study gauge states of closed bosonic string compatified on torus. In addition to the usual gauge states, we will discover soliton gauge states (SGS) in the spectrum of some moduli points. These gauge states and SGS form a realization of enhanced Kac-Moody gauge symmetry group in the gauge state sector of the spectrum. Since T-duality group is the Weyl subgroup of the enhanced gauge group, SGS can be considered as the origin of the discrete T-duality group. In section II, we derive massless gauge states of bosonic string compatified on $`R^{25}T^1`$ at self-dual point $`R=\sqrt{2}`$, and show that they form a representation of the enhanced $`SU(2)_RSU(2)_L`$ gauge group. In section III, we generalize the calculation to $`R^{26D}T^D`$ and give examples at some moduli points. Section IV is devoted to the discussion of massive SGS. We will find that there is an infinite number of massive SGS which exists at some moduli points. The existence of these massive SGS implies that there is an infinite enhanced gauge symmetry of compatified string theory. Finally a brief conclusion is given in section V.
## II Soliton Gauge State on $`R^{25}T^1`$
In the simplest torus compatification, one coordinate of the string was compatified on a circle of radius $`R`$
$$X^{25}(\sigma +2\pi ,\pi )=X^{25}(\sigma ,\pi )+2\pi R_n$$
(1)
Singlevaluedness of the wave function then restricts the allowed momenta to be $`p^{25}=m/R`$ with $`m,nZ`$. The mode expansion of the compatified coordinate for right (left) mover is
$`X_R^{25}`$ $`=`$ $`{\displaystyle \frac{1}{2}}x^{25}+\left(p^{25}{\displaystyle \frac{1}{2}}nR\right)\left(\tau \sigma \right)+i{\displaystyle \underset{r0}{}}{\displaystyle \frac{1}{r}}\alpha _r^{25}e^{ir\left(\tau \sigma \right)}`$ ()
$`X_L^{25}`$ $`=`$ $`{\displaystyle \frac{1}{2}}x^{25}+\left(p^{25}+{\displaystyle \frac{1}{2}}nR\right)\left(\tau +\sigma \right)+i{\displaystyle \underset{r0}{}}{\displaystyle \frac{1}{r}}\underset{r}{\overset{25}{\stackrel{}{\alpha }}}e^{ir\left(\tau +\sigma \right)}`$ ()
We have normalized the string tension to be $`\frac{1}{4\pi T}=1`$ or $`\alpha ^{}=2`$ .The Virasoro operators can be written as
$`L_0`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(p^{25}{\displaystyle \frac{1}{2}}nR\right)+{\displaystyle \frac{1}{2}}p^{\mu ^2}+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\alpha _n\alpha _n`$ ()
$`\stackrel{}{L_0}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(p^{25}+{\displaystyle \frac{1}{2}}nR)+{\displaystyle \frac{1}{2}}p^{\mu ^2}+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\underset{n}{\overset{}{\alpha }}\underset{n}{\overset{}{\alpha }}`$ ()
and
$`L_m`$ $`=`$ $`{\displaystyle \frac{1}{2}}\alpha _0^2+{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}\alpha _{mn}\alpha _n`$ ()
$`\underset{m}{\overset{}{L}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\underset{0}{\overset{2}{\stackrel{}{\alpha }}}+{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}\underset{mn}{\overset{}{\alpha }}\underset{n}{\overset{}{\alpha }}\text{ }(m0)`$ ()
where
$`\alpha _0^{25}`$ $`=`$ $`p^{25}{\displaystyle \frac{1}{2}}nRp_R^{25}`$ ()
$`\underset{0}{\overset{25}{\stackrel{}{\alpha }}}`$ $`=`$ $`p^{25}+{\displaystyle \frac{1}{2}}nRp_L^{25}`$ ()
and the 25d momentum is $`\alpha _0^\mu =\underset{0}{\overset{\mu }{\stackrel{}{\alpha }}}=p^\mu k^\mu `$. In the old covariant quantization of the theory, in addition to the physical propagating states, there are four types of gauge states in the spectrum
$$I.a|\psi =L_1|\chi \text{ where }L_m|\chi =0,(\underset{m}{\overset{}{L}}\delta _m)|\chi =0,\text{ }(m=0,1,2,\mathrm{})$$
(2)
$$II.a|\psi =(L_2+\frac{3}{2}L_1^2)|\chi \text{ where }(L_m+\delta _m)|\chi =0,(\underset{m}{\overset{}{L}}\delta _m)|\chi =0,\text{ }(m=0,1,2,\mathrm{})$$
(3)
and by interchanging all left and right mover operators, one gets $`I.b`$ and $`II.b`$ states. Type II states are zero-norm gauge states only at critical space-time dimension. We will only calculate type a states. Similar results can be easily obtained for type $`b`$ states. For type $`I.a`$ state, the $`m=0`$ constraint of Eq. (2.10) gives
$$M^2=\frac{m^2}{R^2}+\frac{1}{4}n^2R^2+N+\stackrel{}{N}1$$
(4)
$$N\stackrel{}{N}=mn1$$
(5)
where $`N\underset{n=1}{\overset{\mathrm{}}{}}\alpha _n\alpha _n`$ and $`\stackrel{}{N\text{ }}\underset{n=1}{\overset{\mathrm{}}{}}\underset{n}{\overset{}{\alpha }}\underset{n}{\overset{}{\alpha }}`$. For massless $`M^2=0`$ states, $`N+\stackrel{}{N}=0`$ or $`1`$. The solutions of Eqs. (2.12) and (2.13) are
$$N=0,\stackrel{}{N}=1,m=n=0\text{ (}\text{any R}\text{)}$$
(6)
or
$$N=\stackrel{}{N}=0,m=n=\pm 1,R=\sqrt{2}$$
(7)
Equation (2.15) gives us our first SGS. It is easy to write down the explicit form of $`|\chi `$ and $`|\psi `$, and impose the $`m0`$ constraints of Eq. (2.10). There are also a vector and a scalar gauge states in Eq. (2.14). Similar results can be obtained for the type $`I.b`$ state. In this case, $`m=n=\pm 1`$. There is no type II solution in the massless case. We note that there are massless soliton gauge states only when $`R=\sqrt{2}`$ which is known as self-dual point in the moduli space. The vertex operators of all gauge states are calculated to be
$$k_\mu \theta _\nu X_R^\mu \stackrel{\mathrm{\_}}{}X_L^\nu e^{ikx};\text{ }LR$$
(8)
$`k_\mu X_R^\mu \stackrel{\mathrm{\_}}{}X_L^{25}e^{ikx}`$ ()
$`k_\mu \stackrel{\mathrm{\_}}{}X_L^\mu X_R^{25}e^{ikx}`$ ()
$`k_\mu X_R^\mu e^{\pm i\sqrt{2}X_L^{25}}e^{ikx}`$ ()
$`k_\mu X_L^\mu e^{\pm i\sqrt{2}X_R^{25}}e^{ikx}`$ ()
It is easy to see that the three gauge states of Eqs. (2.18) and (2.20) form a representation of $`SU(2)_R`$ Kac-Moody algebra. Similarly, Eqs. (2.17) and (2.19) form a repersentation of $`SU(2)_L`$ Kac-Moody algebra. The vector gauge states in Eq. (2.16) are responsible for the gauge symmetry of graviton and antisymmetric tensor field. We see that the self-dual point $`R=\sqrt{2}`$ is very special even from the gauge sector point of view.
## III Soliton Gauge State on $`R^{26D}T^D`$
In this section we compatify D coordinates on a D-dimensional torus $`T^D\frac{R^D}{2\pi \mathrm{\Lambda }^D}`$
$$\stackrel{}{X}(\sigma +2\pi ,\pi )=\stackrel{}{X}(\sigma ,\pi )+2\pi \stackrel{}{L}$$
(9)
with
$$\stackrel{}{L}=\underset{i=1}{\overset{D}{}}n_i\left(R_i\frac{\underset{i}{\overset{}{e}}}{\sqrt{2}}\right)\left(\mathrm{\Lambda }^D\right)$$
(10)
where $`\mathrm{\Lambda }^D`$ is a D-dimensional lattice with a basis $`\{R_1\frac{\underset{1}{\overset{}{e}}}{\sqrt{2}},R_2\frac{\underset{2}{\overset{}{e}}}{\sqrt{2}},\mathrm{},R_D\frac{\underset{D}{\overset{}{e}}}{\sqrt{2}}\}`$. We have chosen $`|\underset{i}{\overset{}{e}}|^2=2`$. The allowed momenta $`\stackrel{}{p}`$ take values on the dual lattice of $`\mathrm{\Lambda }^D`$
$$\stackrel{}{p}=\underset{i=1}{\overset{D}{}}m_i\left(\frac{1}{R_i}\sqrt{2}\underset{i}{\overset{}{\stackrel{}{e}}}\right)\left(\mathrm{\Lambda }^D\right)^{}$$
(11)
The basis of $`\left(\mathrm{\Lambda }^D\right)^{}`$ is $`\{\frac{1}{R1}\sqrt{2}\underset{1}{\overset{}{\stackrel{}{e}}},\frac{1}{R_2}\sqrt{2}\underset{2}{\overset{}{\stackrel{}{e}}},\mathrm{},\frac{1}{R_D}\sqrt{2}\underset{D}{\overset{}{\stackrel{}{e}}}\}`$ and we have $`\underset{i}{\overset{}{e}}\underset{i}{\overset{}{\stackrel{}{e}}}=\delta _{ij}`$. The mode expansion of the compatified coordinates is
$`\underset{R}{\overset{}{X}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\stackrel{}{x}+(\stackrel{}{p}{\displaystyle \frac{1}{2}}\stackrel{}{L})(\tau \sigma )+i{\displaystyle \underset{r0}{}}{\displaystyle \frac{1}{r}}\alpha _r^{25}e^{ir\left(\tau \sigma \right)}`$ ()
$`\underset{L}{\overset{}{X}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\stackrel{}{x}+(\stackrel{}{p}+{\displaystyle \frac{1}{2}}\stackrel{}{L})(\tau +\sigma )+i{\displaystyle \underset{r0}{}}{\displaystyle \frac{1}{r}}\alpha _r^{25}e^{ir\left(\tau +\sigma \right)}`$ ()
The right and left momenta are defined to be $`\underset{R}{\overset{}{p}}=(\stackrel{}{p}\frac{1}{2}\stackrel{}{L})`$ and $`\underset{L}{\overset{}{p}}=(\stackrel{}{p}+\frac{1}{2}\stackrel{}{L})`$. It can be shown that the 2D-vector $`(\underset{R}{\overset{}{p}},\underset{L}{\overset{}{p}})`$ build an even self-dual Lorentzian lattice $`\mathrm{\Gamma }_{D,D}`$, which guarantees the string one loop modular invariance of the theory . The moduli space of the theory is
$$\mu =\frac{SO(D,D)}{SO\left(D\right)\times SO\left(D\right)}/O(D,D,Z)$$
(12)
where $`O(D,D,Z)`$ is the discrete T-duality group and $`dim`$ $`\mu =D^2`$. To complete the parametrization of the moduli space, one needs to introduce an antisymmetric tensor field $`B_{ij}`$ in the bosonic string action. This will modify the right (left) momenta to be
$$\underset{R}{\overset{}{p}}=(\underset{B}{\overset{}{p}}\frac{1}{2}\stackrel{}{L})$$
(13)
$$\underset{L}{\overset{}{p}}=(\underset{B}{\overset{}{p}}+\frac{1}{2}\stackrel{}{L})$$
(14)
where
$$\underset{B}{\overset{}{p}}=\underset{i,j}{}(m_i\frac{1}{R_i}\sqrt{2}\underset{i}{\overset{}{\stackrel{}{e}}}n_j\frac{1}{\sqrt{2}R_i}B_{ij}\underset{i}{\overset{}{\stackrel{}{e}}})$$
(15)
We are now ready to discuss the gauge state. As a first step, we restrict ourselves to moduli space with $`B_{ij}=0`$ . For the type $`I.a`$ state, the $`m=0`$ constraint of Eq. (2.10) for massless states gives
$$N+\stackrel{}{N}+\stackrel{2}{\stackrel{}{p}}+\frac{1}{4}\stackrel{2}{\stackrel{}{L}}=1$$
(16)
$$N\stackrel{}{N}=\underset{i}{}m_in_i1$$
(17)
It is easy to see $`N+\stackrel{}{N}=0`$ or $`1`$. For $`N+\stackrel{}{N}=1`$, $`m_i=n_i=0`$, we have trivial gauge state solutions. SGS exists for the case $`N+\stackrel{}{N}=0`$ and the following moduli points
$$R_i=\sqrt{2},e_i^I=\sqrt{2}\delta _i^I\text{ }\left(i=1,2,\mathrm{},d\right)$$
(18)
with $`m_i=n_i=\pm 1`$, and $`m_j=n_j=0`$ for $`d<jD`$. In each case, the gauge states and SGS form a representation of $`SU(2)^d`$ algebra. Similar results can be easily obtained for the type $`I.b`$ SGS. As in section II, there is no massless type II SGS. We now discuss $`B_{ij}0`$ case. For illustration, we choose $`D=2`$. In this case $`B_{ij}=Bϵ_{ij}`$, and one has four moduli parameters $`R_1,R_2,B,`$ and $`\underset{1}{\overset{}{e}}\underset{2}{\overset{}{e}}`$. For type $`I.a`$ state, the $`m=0`$ constraint of Eq. (2.10) gives
$$N+\stackrel{}{N}+\underset{B}{\overset{2}{\stackrel{}{p}}}+\frac{1}{4}\stackrel{2}{\stackrel{}{L}}=1$$
(19)
$$N\stackrel{}{N}=m_1n_1+m_2n_21$$
(20)
SGS exists only for $`N+\stackrel{}{N}=0`$. For the moduli point
$$R_1=R_2=\sqrt{2},B=\frac{1}{2},\underset{1}{\overset{}{e}}=(\sqrt{2},0),\underset{2}{\overset{}{e}}=(\sqrt{\frac{1}{2}},\sqrt{\frac{3}{2}})$$
(21)
one gets six SGS with momenta $`\underset{R}{\overset{}{p}}`$ being the six root vectors of $`SU(3)_R`$. Together with two other trivial gauge states corresponding to $`N=0,\stackrel{}{N}=1`$ , they form the Frenkel-Kac-Segal representation of $`SU(3)_{k=1}`$ Kac-Moody algebra. Note that $`\underset{1}{\overset{}{e}},\underset{2}{\overset{}{e}}`$ are the two simple roots of $`SU(3)`$ and $`\underset{1}{\overset{}{\stackrel{}{e}}}=(\sqrt{\frac{1}{2}},\sqrt{\frac{1}{6}}),\underset{2}{\overset{}{\stackrel{}{e}}}=(0,\sqrt{\frac{2}{3}})`$. The six sets of winding number are $`(m_1,n_1,m_2,n_2)=(1,1,0,0),(1,1,0,0),(0,0,1,1),(0,0,1,1),(1,1,1,0),(1,1,1,0)`$. Similar results can be obtained for type $`I.b`$ SGS. The gauge states (including SGS) thus form a representation of enhenced $`SU(3)_R`$ $`SU(3)_L`$ at the moduli point of Eq. (3.15). In general, we expect that all enhenced Kac-Moody gauge symmetry at any moduli point should have a realization on SGS.
## IV Massive Soliton Gauge State
In this section we derive the massive SGS at the first massive level $`M^2=2`$. We will find that SGS exists at infinite number of moduli points. One can also show that they exist at an infinite number of massive level. The existence of these massive SGS implies that there is an infinite enhanced gauge symmetry structure of compatified string theory. For type $`I.a`$ state, the $`m=0`$ constraint of Eq. (2.10) gives
$$\frac{m^2}{R^2}+\frac{1}{4}n^2R^2+N+\stackrel{}{N}=3$$
(22)
$$N\stackrel{}{N}=mn1$$
(23)
which implies $`N+\stackrel{}{N}=0,1,2,3`$. Equations (4. 1) and (4.2) can be easily solved as following:
1. $`N+\stackrel{}{N}=3:`$
$$m=n=0,N=1,\stackrel{}{N}=2,\text{any R}$$
(24)
2. $`N+\stackrel{}{N}=2:`$
$`mn`$ $`=`$ $`1,N=\stackrel{}{N}=1,R=\sqrt{2}`$ (25)
$`mn`$ $`=`$ $`1,N=0,\stackrel{}{N}=2,R=\sqrt{2}`$ ()
3. $`N+\stackrel{}{N}=1:`$
$`mn`$ $`=`$ $`2,N=1,\stackrel{}{N}=0,R=2,1.\text{ }(Tduality)`$ (26)
$`mn`$ $`=`$ $`0,N=0,\stackrel{}{N}=1,R={\displaystyle \frac{\left|m\right|}{\sqrt{2}}},{\displaystyle \frac{2\sqrt{2}}{\left|m\right|}}.\text{ }(Tduality)`$ ()
4. $`N+\stackrel{}{N}=0:`$
$$mn=1,N=\stackrel{}{N}=1,R=2\pm \sqrt{2}.\text{ }(Tduality)$$
(27)
where we have included a T-duality transformation $`R\frac{2}{R}`$ for some moduli points. Note that Eq. (4.5) tells us that massive SGS exists at an infinite number of moduli point. For type $`II.a`$ state, the $`m=0`$ constraint of Eq. (2.11) gives
$$\frac{m^2}{R^2}+\frac{1}{4}n^2R^2+N+\stackrel{}{N}=2$$
(28)
$$N\stackrel{}{N}=mn2$$
(29)
which implies $`N+\stackrel{}{N}=0,1,2`$. Equations (4.7) and (4.8) can be solved as following:
1. $`N+\stackrel{}{N}=2:`$
$$m=n=0,N=0,\stackrel{}{N}=2,\text{any R}$$
(30)
2. $`N+\stackrel{}{N}=1:`$
$$mn=1,N=0,\stackrel{}{N}=1,R=\sqrt{2}$$
(31)
3. $`N+\stackrel{}{N}=0:`$
$$mn=2,N=\stackrel{}{N}=0,R=2,1.\text{ }(Tduality)$$
(32)
The vertex operators of all SGS can be easily calculated and written down. Similar results can be obtained for type $`b`$ gauge state. One can also calculate propagating soliton states by using the same technique. We summarize the moduli points which exist soliton state and SGS as following:
a. Soliton gauge state :
$$R=\sqrt{2},2\pm \sqrt{2},\frac{\left|m\right|}{\sqrt{2}},\frac{2\sqrt{2}}{\left|m\right|},2,1$$
(33)
b. Soliton state :
$$R=\sqrt{2},2\pm \sqrt{2},\frac{\left|m\right|}{\sqrt{2}},\frac{2\sqrt{2}}{\left|m\right|},\frac{\left|m\right|}{2},\frac{4}{\left|m\right|}$$
(34)
In Eqs. (4.12) and (4.13), $`mZ_+`$. There is one interesting remark we would like to point out by the end of this section. One notes that in the second case of Eq. (4.5), instead of specifying $`M^2=2`$, in general we have
$$\frac{m^2}{R^2}+\frac{1}{4}n^2R^2=M^2$$
(35)
with $`mn=0`$. For say $`R=\sqrt{2}`$, one gets $`M^2=\frac{m^2}{2}\left(n=0\right)`$. This means that we have an infinite number of massive SGS at any higher massive level of the spectrum. One can even explicitly write down the vertex operators of these SGS. We conjecture that the $`w_{\mathrm{}}`$ symmetry of 2D string theory can be realized in these SGS . Other moduli points also consist of higher massive SGS in the spectrum.
## V Conclusion
It is hoped that all space-time symmetry of string theory are due to the existence of gauge state in the spectrum. The Heterotic gauge state for the 10D Heterotic string and discrete gauge state for the toy 2D string are such examples. We have introduced soliton gauge state (SGS) for compatified string in this paper, and have related them to the enhanced Kaluza-Klein Kac-Moody gauge symmetry in the theory. In many cases, especially for the massive states, it is easier to study gauge symmetry in the gauge state sector than in the propagating spectrum directly. Since the discrete T-duality symmetry group for bosonic string is the Weyl subgroup of the enhanced gauge group, it can also be considered as implied by the existence of SGS. It is not clear whether other discrete duality symmetry group can be understood in this way. Finally, it would be interesting to consider more complicated compatification, e.g. orbifold and Calabi-Yau compatifications and study the relation between SGS and duality symmetry. works in this direction is in progress.
## VI Acknowledgments
This research is supported by National Science Council of Taiwan, R.O.C., under grant number NSC86-2112-M009-016.
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# Episodic lithium production by extra-mixing in red giants
## 1 Introduction
In low-mass red giants ($`M2.5M_{}`$) the excursion of the base of the convective envelope (BCE) into the radiative zone separating it from the hydrogen-burning shell (HBS) initiates the first dredge-up episode during which the surface composition experiences modest changes. The best indicators of the 1st dredge-up are the ratio <sup>12</sup>C/<sup>13</sup>C and the Li abundance, both decreasing. A decrease of the C and an increase of N are also predicted (Boothroyd & Sackmann (1999)). In the overwhelming majority of the field red giants (K giants) the surface Li abundance is low. However, there are several percent of K giants which possess surprisingly high Li abundances, sometimes exceeding the initial value of $`\mathrm{log}\epsilon (^7\text{Li})3`$ for Pop. I stars (Wallerstein & Sneden 1982; Hanni 1984; Brown et al. 1989; Berdyugina & Savanov 1994; Da Silva et al. 1995; De la Reza et al. 1996, 1997). It was not until recently when the evolutionary status of these Li-rich giants (LIRGs) became clearer. Owing to the more accurate luminosities estimated through the Hipparcos parallaxes Jasniewicz et al. (1999) have concluded that most of the LIRGs are past the 1st dredge-up red giant branch (RGB) stars (the low <sup>12</sup>C/<sup>13</sup>C ratios observed in the LIRGs support this) and hence the initial abundance of Li in their atmospheres could not have been preserved.
De la Reza et al. (1997), Jasniewicz (1999) and other researchers argue convincingly for an internal mechanism of Li production in LIRGs. In both papers quoted the “<sup>7</sup>Be-transport” mechanism (Cameron & Fowler (1971)) is considered as the most promising internal source of Li. This internal mechanism has recently received new support: Castilho et al. (1999) have reported that Be – as predicted – is very depleted (about 10 times) in both LIRGs they studied. In this Letter we propose a combined scenario of Li-enrichment: a giant planet (or brown dwarf) is engulfed by a red giant; this external source or event activates inside the giant the “<sup>7</sup>Be-mechanism” producing then Li internally. Before presenting our model in Sect. 3, we will discuss existing evidence for extra-mixing in red giants (Sect. 2), because this is crucial for scrutinizing any scenario. Sect. 4 concludes the paper.
## 2 Extra-mixing in red giants
The 1st dredge-up episode ends when the BCE stops its excursion into the star and begins to retreat. According to the standard theory, after this moment the surface composition is not expected to change further on the RGB. However, this expectation is not supported by observations. It appears that in low-mass field red giants the surface abundances of C, N, and Li as well as the <sup>12</sup>C/<sup>13</sup>C ratio, and in globular cluster red giants even O, Na, Mg and Al, do continue to change up to the RGB tip. This discrepancy has been quite satisfactorily explained by red giants models with extra-mixing placed between the HBS and the BCE (see Denissenkov & Tout (2000), and references therein). Such extra-mixing is commonly believed to work only in radiative regions with a nearly uniform composition. Therefore, it is expected to come into play after the HBS, moving outwards in mass, will have erased the H-discontinuity left behind by the retreating BCE. Standard model calculations show that the HBS arrives at the H-discontinuity while the star is still on the RGB only in low-mass stars, which is well confirmed by observations (Charbonnel (1994)).
Quite recently Gratton et al. (1999) have determined Li, C, N, O and Na abundances and <sup>12</sup>C/<sup>13</sup>C ratios for a large sample of field stars with accurate luminosity estimates in the metallicity range $`2`$ \[Fe/H\] $`1`$. In Fig. 1 we compare $`\mathrm{log}\epsilon (^7\text{Li})`$, \[C/Fe\] (\[A/B\] means $`\mathrm{log}[n(\text{A})/n(\text{B})]_{\text{star}}\mathrm{log}[n(\text{A})/n(\text{B})]_{}`$), $`\mathrm{log}^{12}\text{C}/^{13}\text{C}`$ and \[N/Fe\] for these stars with results of calculations obtained with the method and code of Denissenkov & Weiss (1996) in which extra-mixing is modeled by diffusion in a post-processing approach, which uses full stellar evolution models as background models for the parameterized diffusion and nucleosynthesis. One can see that the behaviour of the plotted abundances on the upper-RGB ($`\mathrm{log}L/L_{}>2`$, following the nomenclature of Gratton et al. 1999) is quite well reproduced by the diffusive mixing with a depth $`\delta m_{\text{mix}}=0.12`$ ($`\delta m`$ is a relative mass coordinate such that $`\delta m=0`$ at the bottom of the HBS and $`\delta m=1`$ at the BCE) and a rate $`D_{\text{mix}}=510^8`$ cm$`{}_{}{}^{2}`$s<sup>-1</sup> (for details about method and results, see also Denissenkov & Weiss (1996)).
From the upper panel of Fig. 1 one can infer that (i) most of the Pop. II main sequence stars preserve the initial Li abundance in their atmospheres ($`\mathrm{log}\epsilon (^7\text{Li})2.3`$ for Pop. II stars), (ii) during the 1st dredge-up Li is diluted exactly down to the level predicted by the standard theory (Sackmann & Boothroyd (1999)), and (iii) extra-mixing on the upper-RGB further decreases the surface Li abundance; extra-mixing is therefore a necessary ingredient to explain this behaviour.
Contrary to the field Pop. II giants which show neither O depletion nor Na enhancement, in globular clusters there are star-to-star variations of both O and Na on the RGB. Even more important is the fact that in globular cluster red giants Na anticorrelates with O (Fig. 2, symbols). A summary of the observational status can be found in Sneden (1999). The global anticorrelation of \[Na/Fe\] vs. \[O/Fe\] can be explained by extra-mixing as well (Denissenkov & Weiss (1996)), but in this case deeper mixing is required. In Fig. 2 we compare observational data with a sample calculation (solid line, calculated with $`\delta m_{\text{mix}}=0.06`$, $`D_{\text{mix}}=510^8`$ cm$`{}_{}{}^{2}`$s<sup>-1</sup>). The corresponding evolution of the surface Li abundance for this case is plotted in Fig. 3 (line 3a). We note that in all our cases (Fig. 1, solid lines; Figs. 2 and 3) the calculations with extra-mixing start from the same red giant model with $`M=0.8M_{}`$, $`\mathrm{log}L/L_{}=2.1`$ and a heavy elements content of $`Z=510^4`$ (\[Fe/H\] $`\mathrm{log}Z/Z_{}=1.58`$). Because of the shallower mixing in the models of the field Pop. II stars no Na was produced in that case, as observed (Gratton et al. (1999)).
To conclude, extra-mixing (with specific values for mixing depth and speed) is necessary to explain abundance trends in the majority of field giants and abundance anomalies in a large number of globular cluster giants. We now turn to the even more peculiar effect of Li-richness.
## 3 A solution to the problem of Li-rich giants
Following their first discovery that the “<sup>7</sup>Be-mechanism” can naturally work in luminous intermediate-mass asymptotic giant branch (AGB) stars (Sackmann & Boothroyd (1992)), Sackmann & Boothroyd (1999) have demonstrated that under certain conditions the same process can produce Li on the first giant branch, too. In AGB stars, <sup>7</sup>Be freshly minted in the reaction <sup>3</sup>He($`\alpha ,\gamma )^7`$Be is quickly mixed away to a cooler region (where Li produced in the reaction <sup>7</sup>Be(e$`{}_{}{}^{},\nu )^7`$Li can survive) by ordinary convection whereas in RGB stars some extra-mixing is required for this.
The majority of field LIRGs have circumstellar dust shells (De la Reza et al. (1996)) and a large number of additional LIRGs have been discovered among stars with IR excess. This feature seems to be the only one to distinguish the LIRGs from ordinary K giants and led De la Reza et al. (1996) to propose a scenario linking the high Li abundances in these stars to the evolution of circumstellar shells. In this scenario every low-mass red giant passes through a short phase during which some internal mechanism initiates atmospherical Li enrichment accompanied by a prompt mass-loss event. De la Reza et al. (1996) have calculated evolutionary paths (in the IRAS color-color diagram) of the detached shells and inferred that the whole cycle completes in about $`10^5`$ years, the very fast initial increase of the surface Li abundance (during the first several thousand years) being followed by the much longer period (up to $`10^5`$ years) of Li depletion.
Recently, Siess & Livio (1999) have considered an original external scenario: a red giant engulfs an orbiting body of sub-stellar mass (brown dwarf or giant planet) which has the initial abundance of Li left unprocessed. This body deposits its Li into the giant’s envelope and also causes a shell ejection as a consequence of associated processes (mass accretion near the BCE where the body is expected to dissolve and subsequent thermal expansion of the overlying layers; for details see the cited paper). This scenario has an obvious disadvantage: it cannot account for Li abundances exceeding the initial one.
In this Letter we propose a combined scenario in which engulfing of a giant planet by a red giant initiates the internal “<sup>7</sup>Be-mechanism”: It takes into account results of quite recent publications where for the first time extremely high Li abundances have been measured in cluster giants. These are the stars IV-101 (\[Fe/H\] $`=1.50`$) in the globular cluster M 3 (Kraft et al. (1999)) and T33 (\[Fe/H\] $`=0.58`$) in the metal-poor open cluster Berkeley 21 (Hill & Pasquini (1999)). In both cases a Li abundance of $`\mathrm{log}\epsilon (^7\text{Li})3.0`$ has been reported. The LIRGs IV-101 and T33 are plotted in Fig. 3 in comparison with 5 Li-normal giants from the same studies. One realizes that an episodical Li-enrichment can happen at any time on the upper-RGB, independent of the red giant’s evolutionary state, thus indicating an external source. Fig. 2 supports this conclusion: Both the Li-rich giant IV-101 (open square and arrow) and another Li-normal one (open square) close to it have Na increased and O decreased and fit well to the global \[Na/Fe\] vs. \[O/Fe\] anticorrelation. At the same time extra-mixing with the parameters adjusted to reproduce this anticorrelation (Fig. 2, solid line) fails to make LIRGs (Fig. 3, lines 3a and 3b for two different values of initial Li). It appears that after having been exposed for a rather long time ($`310^7`$ years) to the “ordinary” extra-mixing which is responsible to the Na-O-anomalies, the star IV-101 – but not the other one – experienced something which suddenly changed its extra-mixing parameters to values appropriate for Li-production.
From our model calculations we found that the “<sup>7</sup>Be-mechanism” can efficiently synthesize Li and after that maintain its high abundance for a long time only if $`0.12\delta m_{\text{mix}}0.18`$ and, even more important, only if $`D_{\text{mix}}10^{11}`$ cm$`{}_{}{}^{2}`$s<sup>-1</sup> (Fig. 3, lines 1c and 2). Denissenkov & Tout (2000) have proposed Zahn’s rotation-driven meridional circulation and turbulent diffusion (Zahn (1992); Maeder & Zahn (1998)) as a physical mechanism for extra-mixing in low-mass red giants. It turns out, however, that with Zahn’s mechanism, values of $`D_{\text{mix}}10^{11}`$ cm$`{}_{}{}^{2}`$s<sup>-1</sup> can be obtained only as upper limits for rotation close to the Keplerian one. In a scenario with engulfing a planet such fast rotation is explained naturally as a result of transferring the planet’s orbital angular momentum to the giant’s envelope (Siess & Livio (1999)). The next question then is how to get the correct mixing depth in this scenario.
The dashed lines in Fig. 3 are similar to those shown in Fig. 10 of Sackmann & Boothroyd (1999). They are the result of calculations under the assumption that mixing depth and rate favourable for the Li-production are constant on the upper-RGB. One of them (like our line 2) has even been used to interpret a LIRG near the RGB tip in the globular cluster NGC 362 by Smith et al. (1999). However, such a straightforward interpretation is not so simple because: (i) mixing under these conditions does not produce Na nor deplete O as is observed in IV-101; (ii) it explains neither the Li-depletion immediately following the Li-production nor the rather short time-scale for the whole cycle; (iii) it requires a very unusual, precise and long-term tuning of the mixing parameters; the tuning appears to be unusual because it assumes shallow but extremely fast mixing compared to that reproducing the \[Na/Fe\] vs. \[O/Fe\] anticorrelation; it would be more natural to expect that faster mixing should be deeper as well.
Thus we propose the following explanation of how the correct mixing depth could appear in the engulfing scenario. According to Siess & Livio (1999) the giant planet (or brown dwarf) dissolves near the BCE in a red giant. After that the rotation profile in the radiative zone takes a step-like shape with a steep increase of the angular velocity up to about a local Keplerian value at the point of deepest penetration by the planet. In the course of the subsequent evolution the HBS moves outwards in mass and after $`810^5`$ years will reach the step in the rotation profile. During a time interval of $`810^4`$ years this step will be crossing a zone $`0.06\delta m0.16`$ where and when the Li-production becomes efficient. Thus we do not need to fix the mixing depth to a preferred value. Instead, the natural growth of the helium core assures that suitable depths will be encountered and very fast mixing (due to the planet’s engulfing) produces Li during this passage (Fig. 3, line 3c).
## 4 Conclusion
A great advantage of the proposed solution is that it can account not only for the Li-production but also for the subsequent Li-depletion. Indeed, we find that after the mixing depth has reduced to less than $`0.08`$, Li begins to be destroyed on a time-scale consistent with the results of De la Reza et al. (1996). It should be emphasized that it is even more difficult to deplete Li quickly after its production than to produce it, and our scenario deals with this naturally.
At the same time, however, it cannot be applied to the Li-rich ($`\mathrm{log}\epsilon (^7\text{Li})1.8`$) star V42 in M5 (Carney et al. (1998)), which appears to be a (low-mass) post-AGB star. Due to the timescales, the Li-enrichment cannot have happened already during the first red giant phase. On the other side, since V42 is only as bright as the RGB-tip ($`M_{\mathrm{Bol}}=3.38`$), but hotter and thus smaller, capturing a companion would have happened already on the RGB. Thus, our scenario fails for this star, which otherwise appears to be typical M5 member, showing standard $`\alpha `$-enhancement and even the O-Na-anticorrelation (Carney et al. (1998)). We can only speculate that its Li-overabundance happened (via the Cameron-Fowler mechanism) during the AGB, where additional deep mixing initiated hot bottom burning as is standard in intermediate-mass stars (Sackmann & Boothroyd (1992)). Due to the thin envelope of this star very modest extra mixing might already be sufficient. We will investigate this possibility in forthcoming work.
Our scenario does neither provide a direct link to the dust shell formation. Siess & Livio (1999) have ascribed the shell detachment to an increased mass loss during the planet’s engulfing but in our scenario this event is separated from the Li-enrichment episode by a time interval of $`710^5`$ years. The following two speculations towards a solution of this problem can be envisaged: (1) The mass of the radiative zone is negligible, and the angular velocity inside it scales as $`r^2`$ (Denissenkov & Tout (2000)); hence, the ratio of the centrifugal acceleration to the gravity scales as $`r^1`$; as after engulfing the planet this ratio is expected to become close to unity near the BCE, then during the subsequent inward excursion of the step in the rotation profile it will surely exceed unity somewhere in the radiative zone, which may initiate dynamical processes of the angular momentum transfer outwards; the latter may be responsible for the increased mass loss. (2) The process of planet engulfing itself may be associated with various dynamical and thermodynamical processes, for instance, a deepening of the convective envelope (Siess & Livio (1999)), which may redistribute the material with the high angular momentum throughout the radiative zone; in this case the fast mixing will be able to penetrate the zone of correct mixing depths from the very beginning. Whether such phenomena happen in a real red giant can be verified only by 3D hydrodynamical simulations.
###### Acknowledgements.
We wish to express our gratitude to the referee, R. Kraft, who pointed out the existence of V42. P.A.D. acknowledges the warm hospitality of the staff of the Max-Planck-Institut für Astrophysik where this study was carried out.
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# The Trans-Planckian Problem of Inflationary Cosmology
## I Introduction
The inflationary Universe scenario is the first theory of the very early Universe to provide a mechanism for the production of density fluctuations on scales of cosmological interest based on causal physics (see also Ref. for initial ideas). The key point is that during the period of inflation fixed comoving scales are stretched exponentially compared to the Hubble radius. Thus, the wavelengths corresponding to the present large-scale structure in the Universe and to the measured Cosmic Microwave Background (CMB) anisotropies were equal to the Hubble radius about 50 Hubble expansion times before the end of inflation. This gives rise to the possibility that causal physics acting before that time can generate fluctuations on these scales while they are of sub-Hubble length.
Most current models of inflation are based on weakly self-coupled scalar matter fields minimally coupled to gravity. In this context, quantum vacuum fluctuations provide a causal mechanism for generating fluctuations. In fact, the coupled linear metric and matter fluctuations can be quantized in a unified manner . The problem reduces to the quantization of a free scalar field with a time-dependent mass (see e.g. Ref. for a comprehensive review). An initial vacuum state thus undergoes squeezing during inflation, and this leads to the generation of fluctuations. According to the standard calculations , the predicted spectrum is scale-invariant (modulo a mild deviation from scale-invariance which stems from the time-dependence of the Hubble constant during the inflationary period).
There are good heuristic reasons to expect a scale-invariant spectrum of fluctuations to emerge from inflation. Since de-Sitter space is time-translation-invariant, one should expect the amplitude of the density fluctuations $`\delta M/M`$ to be independent of the scale (labelled by the comoving wavenumber $`n`$) if measured at the time when the corresponding wavelength crosses the Hubble radius $`l_\mathrm{H}`$ during the inflationary period. Since microphysics cannot change the physical amplitude of the mass fluctuations while the wavelength is larger than $`l_\mathrm{H}`$, one therefore expects $`\delta M/M`$ to be independent of $`n`$ when measured at the time $`t_\mathrm{f}(n)`$ when the scale re-enters the Hubble radius in the post-inflationary Friedmann-Robertson-Walker period:
$$\frac{\delta M}{M}(n,t_\mathrm{f}(n))=\mathrm{const}.,$$
(1)
which is the definition of a scale-invariant Harrison-Zel’dovich spectrum .
The time-translation-invariance is, however, broken in the current models of inflation. The calculations are done by picking an initial time $`t_\mathrm{i}`$ (e.g. the beginning of the inflationary period), by choosing a specific state of the quantum fields at this time (e.g. the local Minkowski vacuum state or the Bunch-Davies vacuum ), by evolving this state using the linearized equations of motion, and by finally calculating the correlation functions and expectation values of interest. In this context, the emergence of a scale-invariant spectrum of fluctuations is seen to arise from a subtle cancellation of the wavenumber dependence in the initial state wave function and in the growth factor before Hubble radius crossing, and thus depends explicitly on the initial state chosen. States can be found which do not yield a scale-invariant spectrum. Thus, it is clear that the prediction of a Harrison-Zel’dovich spectrum is not completely generic in current models of inflation.
There is, however, a much more serious potential problem for the claim that current models of inflation based on weakly self-coupled scalar fields generically lead to a scale-invariant spectrum of fluctuations. Most of these models of inflation involve (see e.g. for a recent review) a period of inflation much longer than the 60 e-foldings of inflation required to solve the horizon and flatness problems of standard cosmology. Since wavelengths exponentially redshift during inflation, the physical wavelengths of the modes which correspond to the present large-scale structure in the Universe were, in those models, much smaller than the Planck length at the initial time $`t_\mathrm{i}`$. Thus, the usual computations of the spectrum of fluctuations involve extrapolating weakly self-coupled field theory coupled to classical gravity into a regime where these theories are known to break down.
This problem is analogous to the Trans-Planckian problem for black hole physics (see Ref. for a recent overview). In black hole physics there is an arbitrarily large blue shift when following modes of Hawking radiation at future infinity into the past, and the usual calculations of Hawking radiation seem suspect (see e.g. Ref. for a discussion of this point).
In the case of the black hole problem, it was recently shown by Unruh , Brout et al. , Hambli and Burgess and by Corley and Jacobson that the prediction of a thermal Hawking spectrum of black hole radiation is insensitive to modifications of the physics at the ultraviolet end of the spectrum. In these works, the dispersion relation of the quantum fields was modified (in rather ad-hoc ways) at energies larger than some ultraviolet scale $`k_\mathrm{C}`$, and it was found that the spectrum of radiation at future infinity at wavenumbers much smaller than $`k_\mathrm{C}`$ is insensitive to the modifications considered. In this sense, Hawking radiation from black holes was shown to be an infrared effect.
The obvious question is whether a similar conclusion will hold for the generation of fluctuations in inflationary cosmology. This is the question we will address in this paper. We will consider a free scalar field in an inflationary background \[de Sitter phase of a Friedmann-Robertson-Walker cosmology with scale factor $`a(t)`$\]. This scalar field can represent the scalar metric fluctuations, the gravitational wave mode, or a matter scalar field on the fixed background geometry - the case of most interest for cosmology corresponds to scalar metric fluctuations. We will modify the usual dispersion relation
$$\omega ^2=k^2,k^2\frac{n^2}{a^2},$$
(2)
where $`n`$ and $`k`$ are the comoving and physical wavenumbers, respectively, for values of $`k`$ larger than some cutoff scale $`k_\mathrm{C}`$, and will calculate the predicted spectrum of fluctuations in the modified theory for well-motivated initial quantum states, states which in the unmodified theory coincide with the state usually chosen as the initial state. The modified dispersion relations which we use are the same as the ones used by Unruh and by Corley and Jacobson . As preferred initial states we will use either the state which minimizes the energy density at the initial time $`t_\mathrm{i}`$, following the approach of Brown and Dutton , or a naive generalization of the local Minkowski vacuum.
We find that in the case of Unruh’s dispersion relation, the spectrum of density fluctuations is unchanged in the minimum energy density initial state. However, in the case of the family of dispersion relations generalizing the choice of Corley and Jacobson, the choice of the minimum energy density initial state leads to a spectrum of fluctuations which, depending on the specific member of the family of dispersion relations chosen, may be characterized by a tilt, by an exponential factor, and by superimposed oscillations.
Our work indicates that the prediction of a scale-invariant spectrum in inflationary cosmology depends sensitively on hidden assumptions about super-Planck-scale physics. This has important implications for the attempts to unify fundamental physics and early Universe cosmology. It is now a rather nontrivial question under which conditions a unified theory of all forces such as string or M-theory will lead to a scale-invariant spectrum, assuming for the moment that it does indeed lead to a period of inflation.
The outline of this paper is as follows. In Section II we demonstrate that the growth of linear density fluctuations, gravitational waves and linear scalar matter fluctuations can all be described in terms of the same framework: that of a free scalar field with a time-dependent mass. In Section III we introduce the two classes of modified dispersion relations which will be used in the calculations. The quantization of the scalar field in the time-dependent background and the construction of the minimum energy density initial state are reviewed in Section IV. Section V contains our calculations for both classes of dispersion relations. Our results are summarized and discussed in the final section.
## II Equivalence between cosmological perturbations and a fictitious scalar field
Without loss of generality, the line element for the spatially flat Friedmann-Lemaître-Robertson-Walker (FLRW) background plus the perturbations can be written in the synchronous gauge according to :
$`\mathrm{d}s^2`$ $`=`$ $`a^2(\eta )\{\mathrm{d}\eta ^2+[\delta _{ij}+h(\eta ,𝐧)Q\delta _{ij}`$ (4)
$`+h_l(\eta ,𝐧){\displaystyle \frac{Q_{,i,j}}{n^2}}+h_{\mathrm{gw}}(\eta ,𝐧)Q_{ij}]\mathrm{d}x^i\mathrm{d}x^j\}.`$
In this equation, the dimensionless quantity $`𝐧`$ is the comoving wavevector related to the physical wavevector $`𝐤`$ through the relation $`𝐤𝐧/a(\eta )`$. $`\eta `$ is the conformal time related to the cosmic time $`t`$ by $`\mathrm{d}t=a(\eta )\mathrm{d}\eta `$. The functions $`h`$ and $`h_l`$ represent the scalar sector and $`Q(x^i)`$ is the eigenfunction of the Laplace operator on the flat spacelike hypersurfaces. The function $`h_{\mathrm{gw}}`$ represents the gravitational waves and $`Q_{ij}(x^i)`$ is the eigentensor of the Laplace operator. It is traceless and transverse, namely $`Q_i{}_{}{}^{i}=Q_{ij}{}_{}{}^{,j}=0`$. It is convenient to introduce the background quantity $`\gamma (\eta )`$ defined by $`\gamma \dot{H}/H^2`$, where a dot means differentiation with respect to cosmic time and $`H`$ is the Hubble rate, $`H\dot{a}/a`$. We can also write $`\gamma =1^{}/^2`$, where $`a^{}/a`$ and a prime denotes differentiation with respect to the conformal time.
In the tensor sector, we define the quantity $`\mu _\mathrm{T}`$ by $`h_{\mathrm{gw}}\mu _\mathrm{T}/a`$. Then, the equation of motion is given by :
$$\mu _\mathrm{T}^{\prime \prime }+\left[n^2\frac{a^{\prime \prime }}{a}\right]\mu _\mathrm{T}=0.$$
(5)
Since gravitational waves do not couple to matter, the last equation is valid for every type of matter.
In the scalar sector, it is convenient to work with a residual gauge invariant variable $`\mu _\mathrm{S}`$ defined by $`\mu _\mathrm{S}[a/(\sqrt{\gamma })](h^{}+\gamma h)`$ where we have supposed $`\gamma 0`$. The case $`\gamma =0`$ must be treated separately (see below). The quantity $`\mu _\mathrm{S}`$ is related to the gauge invariant Bardeen potential by $`\mathrm{\Phi }_\mathrm{B}^{(\mathrm{SG})}=[\gamma /(2n^2)][\mu _\mathrm{S}/(a\sqrt{\gamma })]^{}`$ where the subscript ‘$`\mathrm{SG}`$’ means ‘calculated in the synchronous gauge’ . Therefore, knowing the solution for $`\mu _\mathrm{S}`$ permits the calculation of the Bardeen variable. If matter is described by a scalar field (the inflaton), then one can show that $`\mu _\mathrm{S}`$ obeys the equation:
$$\mu _\mathrm{S}^{\prime \prime }+\left[n^2\frac{(a\sqrt{\gamma })^{\prime \prime }}{(a\sqrt{\gamma })}\right]\mu _\mathrm{S}=0.$$
(6)
The case $`\gamma =0`$ corresponds to a scale factor $`a(t)e^{Ht}`$, i.e. to the de Sitter manifold. Then, one can show that the exact solution to the perturbed Einstein equations is $`\mathrm{\Phi }_\mathrm{B}=0`$: there are no density perturbations at all. This is because when the equation of state is $`p=\rho `$, fluctuations of the inflaton are not coupled to fluctuations of the perturbed metric. Coupling occurs only as a result of the violation of the condition $`p=\rho `$.
Observable quantities can be computed when the initial power spectra are known. These are defined in terms of the two-point correlation functions. For the Bardeen potential one has
$`0|\mathrm{\Phi }_\mathrm{B}(\eta ,𝐱)`$ $`\mathrm{\Phi }_\mathrm{B}(\eta ,𝐱+𝐫)|0`$ (8)
$`{\displaystyle _0^+\mathrm{}}{\displaystyle \frac{\mathrm{d}n}{n}}{\displaystyle \frac{\mathrm{sin}nr}{nr}}n^3P_{\mathrm{\Phi }_\mathrm{B}}(\eta ,n),`$
whereas for gravitational waves the correlator is given by
$$0|h_{ij}(\eta ,𝐱)h^{ij}(\eta ,𝐱+𝐫)|0_0^+\mathrm{}\frac{\mathrm{d}n}{n}\frac{\mathrm{sin}nr}{nr}n^3P_\mathrm{h}(\eta ,n),$$
(9)
where we have written $`h_{ij}=h_{\mathrm{gw}}Q_{ij}`$. We are specially interested in modes which are outside the horizon at the end of inflation, i.e. $`n/(aH)1`$. For these modes, the power spectra do not depend on time and can be written as
$$n^3P_{\mathrm{\Phi }_\mathrm{B}}(n)=A_\mathrm{S}n^{n_\mathrm{S}1},n^3P_\mathrm{h}(n)=A_\mathrm{T}n^{n_\mathrm{T}}.$$
(10)
Let us now consider power law inflation models where the scale factor is given by $`a(\eta )=l_0|\eta |^{1+\beta }`$ where $`\beta `$ is a number such that $`\beta 2`$ and $`l_0`$ has the dimension of a length. The advantage of this class of models is that everything can be calculated exactly. In the case $`\beta =2`$ which corresponds to exponential expansion, the length $`l_0`$ is nothing but the Hubble radius, $`l_\mathrm{H}a^2/a^{}`$ . The function $`\gamma `$ is a constant given by $`\gamma =(\beta +2)/(\beta +1)`$ which vanishes for $`\beta =2`$. We see that Eq. (6) now reduces to Eq. (5). The spectral indices can be determined exactly and read
$$n_\mathrm{S}=2\beta +5,n_\mathrm{T}=2\beta +4.$$
(11)
We have the relation $`n_\mathrm{S}1=n_\mathrm{T}`$ which is valid exactly only for power law inflation.
Let us now consider a massless scalar field $`\mathrm{\Phi }(\eta ,𝐱)`$ living in a FLRW spacetime. It is convenient to Fourier decompose the field and to introduce the quantity $`\mu `$ defined according to $`\mathrm{\Phi }(\eta ,𝐱)[1/(2\pi )^{3/2}]d𝐧(\mu /a)e^{i𝐧𝐱}`$. It is easy to show that the Klein-Gordon equation reduces to the following equation for $`\mu `$
$$\mu ^{\prime \prime }+\left[n^2\frac{a^{\prime \prime }}{a}\right]\mu =0.$$
(12)
This equation is exactly the same as Eq. (5) and Eq. (6). Therefore, investigating the properties of cosmological perturbations is equivalent to investigating the properties of a fictitious scalar field $`\mathrm{\Phi }(\eta ,𝐱)`$. In particular, the calculation of the power spectrum of the scalar and tensor perturbations reduces to the computation of the power spectrum of this fictitious scalar field. In the following, we will restrict our considerations to this case, having in mind that, in fact, we will calculate the power spectra of cosmological perturbations.
Let us make a last remark. Although it seems that we have considered only a limited class of models (i.e. power law inflation), the previous analogy is in fact much more general. This is because the slow roll approximation, valid for a wide class of inflationary models, reduces to first order to power law inflation.
## III Time dependent dispersion relations
In this section, we present the two classes of modified dispersion relations that will be used in this article. Let us return to the equation of motion (12). In this equation, the presence of the term $`n^2`$ is due to the differential operator $`\delta ^{ij}_i_j`$ in the Klein-Gordon equation. In Fourier space, this means that
$$\omega ^2=k^2=\frac{n^2}{a^2}.$$
(13)
The dispersion relation is therefore linear in the physical wavenumber $`k`$: $`\omega =k`$. A possible alteration of the high frequency behaviour of the Klein-Gordon equation can be obtained if we require the presence of a nonlinear function $`F(k)`$ such that $`\omega =F(k)`$ which, for physical wavenumbers smaller than a new characteristic scale $`k_\mathrm{C}`$, i.e. $`kk_\mathrm{C}`$, reduces to $`\omega k`$. This means that the $`n^2`$ term in the Klein-Gordon equation should now be replaced with a time dependent $`n_{\mathrm{eff}}^2(\eta )`$ such that
$$n_{\mathrm{eff}}^2=a^2(\eta )F^2(k)=a^2(\eta )F^2[n/a(\eta )].$$
(14)
We see that, in terms of comoving wavenumbers, we obtain a time dependent dispersion relation. In what follows, we will consider two explicit examples for the function $`n_{\mathrm{eff}}`$. Given the modified dispersion relation, Eq. (12) can now be written as
$$\mu ^{\prime \prime }+\left[n_{\mathrm{eff}}^2\frac{a^{\prime \prime }}{a}\right]\mu =0.$$
(15)
Let us analyze this equation in more detail. We can distinguish three regimes. In Region I, the wavelength of a given mode, $`\lambda (\eta )(2\pi /n)a(\eta )`$, is much smaller than the characteristic length: $`\lambda l_\mathrm{C}`$. The nonlinearities in the dispersion relation play an important role and the solution of the equation of motion depends on the particular form of $`F(k)`$. A crucial issue is that the mode no longer behaves as a free wave initially. As a consequence, the choice of initial conditions cannot be done in the usual way. In Region II, the wavelength of the mode is larger than the characteristic length but still smaller than the Hubble radius, $`l_\mathrm{C}\lambda l_\mathrm{H}`$. In this case, one can consider the dispersion relation to be linear, i.e. $`\mathrm{\Omega }(\eta )0`$ and neglect the term $`a^{\prime \prime }/a`$. Therefore, the solution can be expressed as:
$$\mu _{\mathrm{II}}(\eta )=B_1e^{in\eta }+B_2e^{in\eta }.$$
(16)
Finally, in Region III, the mode is outside the Hubble radius: $`\lambda l_\mathrm{H}`$ and the solution (the growing mode) is given by:
$$\mu _{\mathrm{III}}(\eta )=Ca(\eta ),$$
(17)
where $`C`$ is a $`n`$ dependent constant. This constant has to be determined by performing the matching of $`\mu `$ and $`\mu ^{}`$ at the times of transition between regions I and II and regions II and III, $`\eta _1`$ and $`\eta _2`$ respectively. Then, the spectrum can be calculated and reads
$$n^3P_\mathrm{\Phi }=n^3\left|\frac{\mu }{a}\right|^2=n^3|C|^2.$$
(18)
Let us now turn to the first example of a time dependent modified dispersion relation.
### A Unruh’s dispersion relation
The dispersion relation used by Unruh in Ref. , in the context of black holes physics, is:
$$\omega =F(k)k_\mathrm{C}\mathrm{tanh}^{1/p}[\left(\frac{k}{k_\mathrm{C}}\right)^p],$$
(19)
where $`p`$ is an arbitrary coefficient. For large values of the wave number, this becomes a constant $`k_\mathrm{C}`$ whereas for small values this is a linear law as expected. According to Eq. (14), in the context of cosmology, we take
$$n_{\mathrm{eff}}(\eta )=\frac{2\pi a(\eta )}{l_\mathrm{C}}\mathrm{tanh}^{1/p}[\left(\frac{nl_\mathrm{C}}{2\pi a(\eta )}\right)^p],$$
(20)
where $`l_\mathrm{C}`$ is the characteristic length corresponding to $`k_\mathrm{C}`$. The argument of the hyperbolic tangent can also be rewritten as $`l_\mathrm{C}/\lambda (\eta )`$. This means that when $`\lambda l_\mathrm{C}`$, $`n_{\mathrm{eff}}(\eta )`$ tends to $`n`$.
### B Generalized Corley/Jacobson dispersion relation
The dispersion relation utilized by Corley and Jacobson in Ref. is given by the following expression
$$\omega ^2=F^2(k)k^2\frac{k^4}{k_\mathrm{C}^2}.$$
(21)
In this article, we consider a more general case and write
$$\omega ^2=k^2+k^2\underset{q=1}{\overset{m}{}}b_q\left(\frac{k}{k_\mathrm{C}}\right)^{2q},$$
(22)
where the $`b_q`$ are, a priori, arbitrary coefficients. Let us suppose that the previous sum only contains the last term. The physics depends on the sign of $`b_m`$. If $`b_m`$ is negative, then $`\omega `$ vanishes for $`k=k_\mathrm{C}|b_m|^{2m}`$. Beyond this point, the dispersion relation becomes complex. The Corley/Jacobson case corresponds to $`m=1`$ and $`b_1=1`$. In the context of cosmology, the previous ansatz gives rise to the following function $`n_{\mathrm{eff}}(\eta )`$
$$n_{\mathrm{eff}}^2(\eta )=n^2+n^2\underset{q=1}{\overset{m}{}}\frac{b_q}{(2\pi )^{2q}}\left(\frac{l_\mathrm{C}}{a}\right)^{2q}n^{2q}.$$
(23)
Again, when $`\lambda l_\mathrm{C}`$ then the effective comoving wavenumber simply reduces to $`n`$. On the other hand, when $`\lambda l_\mathrm{C}`$, one has
$$n_{\mathrm{eff}}^2\frac{b_m}{(2\pi )^{2m}}\left(\frac{l_\mathrm{C}}{a}\right)^{2m}n^{2m+2}.$$
(24)
The different dispersion relations used in this article are displayed in Fig. (1) together with the dispersion relation considered in Ref. denoted “KG”.
## IV Quantization of a massive scalar field
The aim of this section is to develop a Lagrangian and Hamiltonian formalism for the system described above. We will show that considering a time-dependent dispersion relation is equivalent to giving a time-dependent mass to the fictitious scalar field. The main purpose of this section is to discuss the initial conditions. As already mentioned previously, when the wavelength of a mode is smaller than the critical length $`l_\mathrm{C}`$, the mode does not behave as a free wave because the dispersion relation in this region is no longer $`\omega =k`$. As a consequence, it is no longer possible to impose the usual initial condition at $`\eta =\eta _\mathrm{i}`$, i.e. $`\mu e^{in(\eta \eta _\mathrm{i})}/\sqrt{2n}`$. Another method must be used. Following Ref. , we will choose the state which initially minimizes the energy density of the field.
### A Lagrangian and Hamiltonian formalisms
We now study a massive fictitious scalar field $`\mathrm{\Phi }`$ whose action is given by
$`S`$ $`=`$ $`{\displaystyle }\mathrm{d}\eta {\displaystyle _{R^{3+}}}\mathrm{d}𝐧[\mu _𝐧^{}\mu _𝐧^{^{}}+{\displaystyle \frac{a^{}_{}{}^{}2}{a^2}}\mu _𝐧\mu _𝐧^{}`$ (26)
$`{\displaystyle \frac{a^{}}{a}}(\mu _𝐧^{}\mu _𝐧^{}+\mu _𝐧\mu _𝐧^{^{}})n_{\mathrm{eff}}^2\mu _𝐧\mu _𝐧^{}].`$
In this equation, the scalar field has been Fourier expanded according to
$$\mathrm{\Phi }(\eta ,𝐱)=\frac{1}{(2\pi )^{3/2}}\frac{1}{a(\eta )}d𝐧\mu _𝐧(\eta )e^{i𝐧𝐱},$$
(27)
and $`\mu _𝐧(\eta )`$ denotes the complex Fourier component of the field. We can easily check that the Lagrange equation of motion for the quantity $`\mu _𝐧(\eta )`$ leads to Eq. (15).
We are now in a position where we can pass to the Hamiltonian formalism. Our first move is to perform the following time-dependent transformation
$$\mu _𝐧(\eta )\frac{1}{N(n,\eta )}\psi _𝐧(\eta ),$$
(28)
where $`N(n,\eta )`$ is a time-dependent factor which will be fixed below. Next, the action given in Eq. (26) expressed in terms of the new variable $`\psi _𝐧(\eta )`$ takes the form
$`S`$ $`=`$ $`{\displaystyle }\mathrm{d}\eta {\displaystyle _{R^{3+}}}\mathrm{d}𝐧[{\displaystyle \frac{1}{N^2}}\psi _𝐧^{}\psi _𝐧^{^{}}+{\displaystyle \frac{1}{N^2}}({\displaystyle \frac{N^{}}{N}}+{\displaystyle \frac{a^{}}{a}})^2\psi _𝐧\psi _𝐧^{}`$ (31)
$`{\displaystyle \frac{1}{N^2}}\left({\displaystyle \frac{N^{}}{N}}+{\displaystyle \frac{a^{}}{a}}\right)\left(\psi _𝐧^{}\psi _𝐧^{}+\psi _𝐧\psi _𝐧^{^{}}\right)`$
$`{\displaystyle \frac{1}{N^2}}n_{\mathrm{eff}}^2\psi _𝐧\psi _𝐧^{}].`$
We can now calculate the conjugate momentum to $`\psi _𝐧(\eta )`$. Its definition is $`p_𝐧\overline{}_𝐧/\psi _𝐧^{}(\eta )`$ where $`\overline{}_𝐧`$ is the Lagrangian density (the bar indicates that one calculates the Lagrangian in Fourier space) which one can deduce from the previous equation. The conjugate momentum reads
$$p_𝐧=\frac{1}{N^2}\left(\psi _𝐧^{}\frac{a^{}}{a}\psi _𝐧\right)\frac{N^{}}{N^3}\psi _𝐧.$$
(32)
The Hamiltonian can be determined using the following relation
$$\overline{}_𝐧p_𝐧\psi _𝐧^{^{}}+p_𝐧^{}\psi _𝐧^{}\overline{}_𝐧.$$
(33)
Inserting the expressions of the Lagrangian and of the the conjugate momentum in this definition, we obtain
$`\overline{}_𝐧`$ $`=`$ $`N^2p_𝐧p_𝐧^{}+{\displaystyle \frac{(aN)^{}}{aN}}\left(\psi _𝐧p_𝐧^{}+\psi _𝐧^{}p_𝐧\right)`$ (35)
$`+{\displaystyle \frac{1}{N^2}}n_{\mathrm{eff}}^2\psi _𝐧\psi _𝐧^{}.`$
The explicit quantization can now be carried out. We express the Fourier component $`\psi _𝐧`$ and its conjugate momentum $`p_𝐧`$ in terms of creation and annihilation operators, satisfying the usual commutation relation $`[c_𝐧,c_𝐫^{}]=\delta (𝐧𝐫)`$, according to
$$\psi _𝐧\sqrt{\mathrm{}}\left(c_𝐧+c_𝐧^{}\right),p_𝐧\frac{\sqrt{\mathrm{}}}{2i}\left(c_𝐧c_𝐧^{}\right).$$
(36)
The Hamiltonian operator is obtained by plugging the previous expressions into Eq. (35) and requiring that ‘$`\mathrm{}\omega /2`$’ be present in each mode, which fixes the normalization factor $`N`$ to be
$$N^2=\mathrm{\hspace{0.17em}2}\omega (\eta ),$$
(37)
where $`\omega `$ is the ‘comoving frequency’ defined by $`\omega (\eta )n_{\mathrm{eff}}`$. Although we use the same notation for convenience, this frequency should not be confused with the physical frequency which appears in Eqns. (19) and (21) and which can obtained by multiplying the comoving frequency by a factor $`1/a`$. The Hamiltonian reads
$`H=`$ $`{\displaystyle _{R^3}}\mathrm{d}𝐧[{\displaystyle \frac{\mathrm{}\omega }{2}}(c_𝐧c_𝐧^{}+c_𝐧c_𝐧^{})+{\displaystyle \frac{i\mathrm{}}{2}}{\displaystyle \frac{(a\sqrt{\omega })^{}}{a\sqrt{\omega }}}(c_𝐧^{}c_𝐧^{}`$ (39)
$`c_𝐧c_𝐧)].`$
This Hamiltonian has the usual structure. The first term is just a collection of harmonic oscillators whereas the second term represents the interaction between the background and the perturbations. This term is responsible for the phenomenon of particle creation, which is a squeezing effect. In a static spacetime, the pump function $`(a\sqrt{\omega })^{}/(a\sqrt{\omega })`$ vanishes and the interaction part of the Hamiltonian disappears. The field operator can be expressed as
$`\mathrm{\Phi }(\eta ,𝐱)`$ $`=`$ $`{\displaystyle \frac{\sqrt{\mathrm{}}}{a(\eta )}}{\displaystyle \frac{1}{(2\pi )^{3/2}}}{\displaystyle }{\displaystyle \frac{\mathrm{d}𝐧}{\sqrt{2\omega (\eta )}}}[c_𝐧(\eta )e^{i𝐧𝐱}`$ (41)
$`+c_𝐧^{}(\eta )e^{i𝐧𝐱}].`$
The time evolution of the creation and annihilation operators and therefore of the quantum scalar field is calculated by means of the Heisenberg equation:
$$i\mathrm{}\frac{\mathrm{d}}{\mathrm{d}\eta }c_𝐧(\eta )=[c_𝐧,H].$$
(42)
Using the form of the Hamiltonian derived previously, one gets the following equations of motion
$`i\mathrm{}{\displaystyle \frac{\mathrm{d}c_𝐧}{\mathrm{d}\eta }}`$ $`=`$ $`\mathrm{}\omega (\eta )c_𝐧+i\mathrm{}{\displaystyle \frac{(a\sqrt{\omega })^{}}{a\sqrt{\omega }}}c_𝐧^{},`$ (43)
$`i\mathrm{}{\displaystyle \frac{\mathrm{d}c_𝐧^{}}{\mathrm{d}\eta }}`$ $`=`$ $`\mathrm{}\omega (\eta )c_𝐧^{}+i\mathrm{}{\displaystyle \frac{(a\sqrt{\omega })^{}}{a\sqrt{\omega }}}c_𝐧.`$ (44)
The solution of these equations is a Bogoliubov transformation which can be written as
$`c_𝐧(\eta )`$ $`=`$ $`u_n(\eta )c_𝐧(\eta _\mathrm{i})+v_n(\eta )c_𝐧^{}(\eta _\mathrm{i}),`$ (45)
$`c_𝐧^{}(\eta )`$ $`=`$ $`u_n^{}(\eta )c_𝐧^{}(\eta _\mathrm{i})+v_n^{}(\eta )c_𝐧(\eta _\mathrm{i}),`$ (46)
where we have introduced two new functions $`u_n(\eta )`$ and $`v_n(\eta )`$. These functions satisfy $`|u_n(\eta )|^2|v_n(\eta )|^2=1`$ in order for the commutation relation given to be preserved in time. Let us notice that $`u_n`$ and $`v_n`$ do not depend on the vector $`𝐧`$ but only on its modulus $`n`$. Inserting the previous equations in Eqns. (43) and (44), one obtains the equation of motion for these two functions
$`i\mathrm{}{\displaystyle \frac{\mathrm{d}u_n}{\mathrm{d}\eta }}`$ $`=`$ $`\mathrm{}\omega (\eta )u_n+i\mathrm{}{\displaystyle \frac{(a\sqrt{\omega })^{}}{a\sqrt{\omega }}}v_n^{},`$ (47)
$`i\mathrm{}{\displaystyle \frac{\mathrm{d}v_n}{\mathrm{d}\eta }}`$ $`=`$ $`\mathrm{}\omega (\eta )v_n+i\mathrm{}{\displaystyle \frac{(a\sqrt{\omega })^{}}{a\sqrt{\omega }}}u_n^{}.`$ (48)
The functions $`u_n`$ and $`v_n`$ can be re-expressed in terms of three other arbitrary functions $`r_n(\eta )`$, $`\theta _n(\eta )`$ and $`\phi _n(\eta )`$. Following this path would lead to the squeezed state formalism. However, we will not need it in this article.
### B Fixing the initial conditions
The previous considerations permit to fix the initial value of the mode function $`\mu _n(\eta _\mathrm{i})`$ and its derivative $`\mu _n^{}(\eta _\mathrm{i})`$ for any choice of function $`\mathrm{\Omega }(\eta )`$, i.e. for any time dependent dispersion relation.
It is straightforward to check that the function
$$\mu _n\frac{1}{N(n,\eta )}(u_n+v_n^{})=\frac{1}{\sqrt{2\omega }}(u_n+v_n^{}),$$
(49)
satisfies Eq. (15)<sup>*</sup><sup>*</sup>*It should be noticed that $`\mu _n`$ is not exactly the mode function introduced before. It is dimensionless (instead of dimension $`\sqrt{\mathrm{}c}`$) and depends only on the modulus $`n`$. In the same manner, we now deal with a ‘new’ function $`\psi _nN\mu _n`$.. From Eqns. (45) and (46), we see that the initial conditions for the two function $`u_n`$ and $`v_n`$ are given by: $`u_n(\eta =\eta _\mathrm{i})=1`$ and $`v_n(\eta =\eta _\mathrm{i})=0`$. Therefore, the initial value of the mode function $`\mu `$ can be written as:
$$\mu (\eta =\eta _\mathrm{i})=\frac{1}{\sqrt{2\omega (\eta _\mathrm{i})}}=\frac{1}{\sqrt{2n_{\mathrm{eff}}}}.$$
(50)
Let us now turn to the determination of $`\mu ^{}(\eta =\eta _\mathrm{i})`$. It will be found by the requirement that the energy density is minimized. The stress energy tensor can be obtained from the action (26) with the help of the standard definition. In terms of the Fourier components $`\psi _n`$, the energy density reads
$`\rho `$ $`=`$ $`{\displaystyle \frac{\mathrm{}}{4\pi ^2a^4}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\mathrm{d}n}{N^2}}[\psi _n^{}\psi _n^{^{}}{\displaystyle \frac{(aN)^{}}{aN}}(\psi _n\psi _n^{^{}}+\psi _n^{}\psi _n^{})`$ (53)
$`+{\displaystyle \frac{a^2}{a^2}}\psi _n\psi _n^{}+{\displaystyle \frac{N^2}{N^2}}\psi _n\psi _n^{}+n_{\mathrm{eff}}^2\psi _n\psi _n^{}`$
$`+2{\displaystyle \frac{a^{}N^{}}{aN}}\psi _n\psi _n^{}].`$
We now define the functions $`x(\eta )`$ and $`y(\eta )`$ as the real and imaginary parts of the ratio $`\psi _n^{}/\psi _nx+iy`$, respectively. Then, the initial energy density can be expressed in terms of $`x_\mathrm{i}x(\eta =\eta _\mathrm{i})`$, $`y_\mathrm{i}y(\eta =\eta _\mathrm{i})`$ and the Wronskian $`W(n)\mu _n^{}\mu _n^{}\mu _n^{^{}}\mu _n`$ which is a time independent quantity (as can be checked in calculating $`\mathrm{d}W(n)/\mathrm{d}\eta `$ and using the equation of motion for $`\mu _n`$)
$`\rho `$ $`=`$ $`{\displaystyle \frac{\mathrm{}}{4\pi ^2a^4}}{\displaystyle _0^{\mathrm{}}}\mathrm{d}n{\displaystyle \frac{W(n)}{2iy_\mathrm{i}}}[x_\mathrm{i}^2+y_\mathrm{i}^22{\displaystyle \frac{(aN)^{}}{aN}}x_\mathrm{i}+{\displaystyle \frac{a^2}{a^2}}`$ (54)
$`+`$ $`n_{\mathrm{eff}}^2+{\displaystyle \frac{N^2}{N^2}}+2{\displaystyle \frac{a^{}N^{}}{aN}}],`$ (55)
where $`N`$ and $`a`$ are also evaluated at the initial time. Notice that, while deriving the previous equation, we used the fact that the Wronskians of $`\mu _n`$ and $`\psi _n`$ are related by a factor $`N^2`$. The ‘vacuum’ used in this article is defined as the state which initially minimizes the energy density. The variation of the previous expression with respect to $`x_\mathrm{i}`$ and $`y_\mathrm{i}`$ leads to
$`\delta \rho `$ $`=`$ $`{\displaystyle \frac{\mathrm{}}{4\pi ^2a^4}}{\displaystyle _0^{\mathrm{}}}\mathrm{d}n{\displaystyle \frac{W(n)}{2i}}\{{\displaystyle \frac{2}{y_\mathrm{i}}}[x_\mathrm{i}{\displaystyle \frac{(aN)^{}}{aN}}]\delta x_\mathrm{i}`$ (58)
$`+{\displaystyle \frac{1}{y_\mathrm{i}^2}}[y_\mathrm{i}^2x_\mathrm{i}^2+2{\displaystyle \frac{(aN)^{}}{aN}}x_\mathrm{i}{\displaystyle \frac{a^2}{a^2}}n_{\mathrm{eff}}^2`$
$`{\displaystyle \frac{N^2}{N^2}}2{\displaystyle \frac{a^{}N^{}}{aN}}\left]\delta y_\mathrm{i}\right\}.`$
Demanding that $`\delta \rho =0`$, one deduces the initial values of $`x`$ and $`y`$
$$x_\mathrm{i}=\frac{a^{}}{a}(\eta _\mathrm{i})+\frac{N^{}}{N}(\eta _\mathrm{i}),y_\mathrm{i}=\pm n_{\mathrm{eff}}.$$
(59)
These expressions can be simplified. Using the explicit form of the function $`N(n,\eta )`$, one can write
$$\frac{N^{}}{N}=\frac{\omega ^{}}{2\omega }.$$
(60)
At the time $`\eta =\eta _\mathrm{i}`$, it is reasonable to consider $`\lambda l_\mathrm{C}`$ (otherwise, the whole problem studied here would be pointless). Then, for Unruh’s dispersion relation, one finds $`N^{}/Na^{}/2a`$ and for the Corley/Jacobson dispersion relation, one has $`N^{}/Nma^{}/2a`$. In addition, $`a^{}/a`$ is very small in the limit where the conformal time goes to $`\mathrm{}`$ since $`a^{}/a(\eta _\mathrm{i})=(1+\beta )/|\eta _\mathrm{i}|`$ and $`|\eta _\mathrm{i}|1`$. Therefore, one gets that $`\psi _n^{}/\psi _niy_\mathrm{i}`$. On the other hand, we have $`\mu _n=\psi _n/N`$. Combining this formula with the previous one, one obtains $`\mu _n^{}+N^{}/N=iy_\mathrm{i}\mu _n`$. Neglecting again the term $`N^{}/N`$, we finally arrive at
$$\mu ^{}(\eta =\eta _\mathrm{i})=\pm i\sqrt{\frac{n_{\mathrm{eff}}}{2}}.$$
(61)
The initial conditions are now completely fixed and given by Eqns. (50) and (61).
Let us also mention that it is possible to adopt another choice of initial conditions which corresponds to the ‘instantaneous Minkowski vacuum’ at $`\eta =\eta _\mathrm{i}`$, namely
$$\mu (\eta _\mathrm{i})=\frac{1}{\sqrt{2n}},\mu ^{}(\eta _\mathrm{i})=\pm i\sqrt{\frac{n}{2}}.$$
(62)
If the dispersion relation is standard, then $`\mathrm{\Omega }=0`$ and the mode is initially free: locally, it does not feel the curvature of space-time and behaves as it were flat. In this case, the two possible choices of initial conditions discussed above coincide.
Two last comments are in order before ending this section. Let us first remark that the concept of an initial state which minimizes the energy density of the field could be problematic in a region where the dispersion relation becomes complex, as it is the case for the Corley/Jacobson dispersion relation with $`b_m<0`$, since the energy needs not to be bounded from below in such a situation. We are not aware of any more obvious method than the one used here to deal with this case.
Finally, although we have introduced two initial states, it should be clear that the minimizing energy state is the only physical vacuum state. The instantaneous Minkowski vacuum is considered here only to stress the fact that the choice of the initial conditions becomes more crucial than in the standard situation where one can show that a large class of initial states leads to the same spectrum (although, as already mentioned in the introduction, it is possible to find examples which do not belong to this class of initial states ).
## V Analytical solutions
In this section, we calculate the spectrum of fluctuations for the two classes of dispersion relations introduced in Section III. We focus on a fixed comoving wavenumber $`n`$ and proceed as follows. We solve the equation of motion in each of the three regions (defined in Section III) separately. The coefficients of the two fundamental solutions in Region I are fixed by the initial conditions discussed above. Then, we explicitly perform the matching of $`\mu `$ and $`\mu ^{}`$ at the transitions between Region I and Region II, which occurs at a time denoted by $`\eta _1`$, and between Region II and Region III, which occurs at time $`\eta _2`$, to obtain the coefficients of the two fundamental solutions in Region III, from which the spectrum can be calculated. The time $`\eta _2`$ is when the mode crosses the Hubble radius, which is given by
$$l_\mathrm{H}(\eta )=\frac{l_0}{|1+\beta |}|\eta |^{2+\beta }.$$
(63)
Thus, the condition $`l_\mathrm{H}(\eta _2)=\lambda (\eta _2)`$ boils down to
$$|\eta _2|=\frac{2\pi }{n}|1+\beta |.$$
(64)
The geometry of space-time is illustrated in Fig. (2).
We start this section with Unruh’s dispersion relation.
### A Unruh’s case
The equation of motion for the mode function can be written as
$$\mu ^{\prime \prime }+\{\frac{4\pi ^2}{l_\mathrm{C}^2}a^2\mathrm{tanh}^{2/p}[\left(\frac{l_\mathrm{C}}{\lambda (\eta )}\right)^p]\frac{a^{\prime \prime }}{a}\}\mu =0.$$
(65)
This equation can be solved exactly in Region I only if the scale factor is given by $`a(\eta )=l_0/|\eta |`$, i.e. in the case $`\beta =2`$. Fortunately, this corresponds to the de Sitter space-time, the prototypical model of inflationary cosmology. Note that in this case $`l_0`$ is the Hubble radius \[see Eq. (63)\]. In Region I, the hyperbolic tangent is approximatively one since $`l_\mathrm{C}\lambda `$ initially. Therefore, Equation (65) reduces to
$$\mu ^{\prime \prime }+\left(\frac{4\pi ^2l_0^2/l_\mathrm{C}^22}{\eta ^2}\right)\mu =0.$$
(66)
Note that, in fact, the form of this last equation is independent of the precise form (i.e. the hyperbolic tangent) of the dispersion relation in the regime $`l_\mathrm{C}\lambda `$. It is just necessary to assume that $`F(k)`$ goes to a constant. We see that the result depends in an essential way on the dimensionless parameter $`ϵl_\mathrm{C}/l_0`$. At this point, we have assumed nothing about the value of the ratio $`l_\mathrm{C}/l_0`$. However, physically, it is clear that $`ϵ1`$. One would expect the cutoff length to be given by the Planck length ($`l_\mathrm{C}l_{\mathrm{Pl}}`$), whereas $`l_010^5l_{\mathrm{Pl}}`$ if the spectrum of fluctuations is COBE normalized. In this case, we have $`ϵ10^5`$. In the following, we will use an expansion in terms of this parameter. The exact solution of Eq. (66) is
$$\mu _\mathrm{I}(\eta )=A_1|\eta |^{x_1}+A_2|\eta |^{x_2},$$
(67)
where the exponents $`x_1`$ and $`x_2`$ are given by
$$x_{1,2}=\frac{1}{2}\pm \frac{1}{2}\sqrt{9\frac{16\pi ^2}{ϵ^2}}.$$
(68)
It is now time to fix the coefficients $`A_1`$ and $`A_2`$. They are completely determined by the initial conditions (50) and (61). In the approximation where $`l_\mathrm{C}\lambda `$, they are solutions of (note that we do not yet use the fact that $`ϵ`$ is small)
$`A_1|\eta _\mathrm{i}|^{x_1}+A_2|\eta _\mathrm{i}|^{x_2}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\sqrt{{\displaystyle \frac{ϵ}{\pi }}}|\eta _\mathrm{i}|^{1/2},`$ (69)
$`A_1x_1|\eta _\mathrm{i}|^{x_11}+A_2x_2|\eta _\mathrm{i}|^{x_21}`$ $`=`$ $`i\sqrt{{\displaystyle \frac{\pi }{ϵ}}}|\eta _\mathrm{i}|^{1/2}.`$ (70)
The exact solution of this system of equations can be written as
$`A_1`$ $`=`$ $`|\eta _\mathrm{i}|^{1/2x_1}{\displaystyle \frac{1}{1x_1/x_2}}{\displaystyle \frac{1}{2}}\sqrt{{\displaystyle \frac{ϵ}{\pi }}}\left(1\pm {\displaystyle \frac{2i\pi }{ϵx_2}}\right),`$ (71)
$`A_2`$ $`=`$ $`|\eta _\mathrm{i}|^{1/2x_2}{\displaystyle \frac{1}{1x_2/x_1}}{\displaystyle \frac{1}{2}}\sqrt{{\displaystyle \frac{ϵ}{\pi }}}\left(1\pm {\displaystyle \frac{2i\pi }{ϵx_1}}\right).`$ (72)
It is at this point that we use the fact that $`ϵ`$ is small. To first order in a systematic expansion in this parameter we obtain
$`A_1`$ $``$ $`{\displaystyle \frac{i}{8}}\left({\displaystyle \frac{ϵ}{\pi }}\right)^{3/2}\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{2i\pi }{ϵ}}\pm {\displaystyle \frac{2i\pi }{ϵ}}\right)`$ (74)
$`\times \mathrm{exp}\left({\displaystyle \frac{2i\pi }{ϵ}}\mathrm{ln}\left|\eta _\mathrm{i}\right|\right),`$
$`A_2`$ $``$ $`{\displaystyle \frac{i}{8}}\left({\displaystyle \frac{ϵ}{\pi }}\right)^{3/2}\left({\displaystyle \frac{1}{2}}+{\displaystyle \frac{2i\pi }{ϵ}}\pm {\displaystyle \frac{2i\pi }{ϵ}}\right)`$ (76)
$`\times \mathrm{exp}\left({\displaystyle \frac{2i\pi }{ϵ}}\mathrm{ln}\left|\eta _\mathrm{i}\right|\right).`$
We now pursue the calculation for both choices of the sign of the initial conditions. We introduce an index ‘u’ for the upper choice and ‘l’ for the lower choice. This leads to
$`A_1^u`$ $`=`$ $`{\displaystyle \frac{i}{16}}\left({\displaystyle \frac{ϵ}{\pi }}\right)^{3/2}\mathrm{exp}\left({\displaystyle \frac{2i\pi }{ϵ}}\mathrm{ln}\left|\eta _\mathrm{i}\right|\right),`$ (77)
$`A_2^u`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{ϵ}{\pi }}\right)^{1/2}\mathrm{exp}\left({\displaystyle \frac{2i\pi }{ϵ}}\mathrm{ln}\left|\eta _\mathrm{i}\right|\right),`$ (78)
$`A_1^l`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{ϵ}{\pi }}\right)^{1/2}\mathrm{exp}\left({\displaystyle \frac{2i\pi }{ϵ}}\mathrm{ln}\left|\eta _\mathrm{i}\right|\right),`$ (79)
$`A_2^l`$ $`=`$ $`{\displaystyle \frac{i}{16}}\left({\displaystyle \frac{ϵ}{\pi }}\right)^{3/2}\mathrm{exp}\left({\displaystyle \frac{2i\pi }{ϵ}}\mathrm{ln}\left|\eta _\mathrm{i}\right|\right).`$ (80)
Therefore, one has $`A_2^uA_1^u`$, $`A_1^lA_2^l`$ and only one branch of the solution (67) survives. Then, the solution in Region I can be expressed as
$$\mu _\mathrm{I}^{u,l}(\eta )=\frac{1}{2}\sqrt{\frac{ϵ|\eta |}{\pi }}\mathrm{exp}\left(\frac{2i\pi }{ϵ}\mathrm{ln}\left|\frac{\eta }{\eta _\mathrm{i}}\right|\right).$$
(81)
Let us now turn to Region II. As already mentioned above, in this region, the solution is given by
$$\mu _{\mathrm{II}}(\eta )=B_1e^{in\eta }+B_2e^{in\eta }.$$
(82)
The coefficients $`B_1`$ and $`B_2`$ are determined by the matching of this solution with the solution (81) at the time $`\eta _1`$. Continuity of $`\mu `$ and $`\mu ^{}`$ yields
$`inB_1e^{in\eta _1}`$ $`+`$ $`inB_2e^{in\eta _1}`$ (83)
$`=`$ $`{\displaystyle \frac{in}{2}}\left({\displaystyle \frac{ϵ|\eta _1|}{\pi }}\right)^{1/2}\mathrm{exp}\left({\displaystyle \frac{2i\pi }{ϵ}}\mathrm{ln}\left|{\displaystyle \frac{\eta _1}{\eta _\mathrm{i}}}\right|\right),`$ (84)
$`inB_1e^{in\eta _1}`$ $``$ $`inB_2e^{in\eta _1}`$ (85)
$`=`$ $`\pm i\left({\displaystyle \frac{\pi }{ϵ|\eta _1|}}\right)^{1/2}\mathrm{exp}\left({\displaystyle \frac{2i\pi }{ϵ}}\mathrm{ln}\left|{\displaystyle \frac{\eta _1}{\eta _\mathrm{i}}}\right|\right).`$ (86)
The solution can be found easily and reads
$`B_1`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{n}}}(1\pm 1)\mathrm{exp}({\displaystyle \frac{2i\pi }{ϵ}}\mathrm{ln}\left|{\displaystyle \frac{\eta _1}{\eta _\mathrm{i}}}\right|+{\displaystyle \frac{2i\pi }{ϵ}}),`$ (87)
$`B_2`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{n}}}(11)\mathrm{exp}({\displaystyle \frac{2i\pi }{ϵ}}\mathrm{ln}\left|{\displaystyle \frac{\eta _1}{\eta _\mathrm{i}}}\right|{\displaystyle \frac{2i\pi }{ϵ}}),`$ (88)
As a consequence, the solution in Region II also contains only one branch.
Finally, we must solve the mode equation in Region III. As already mentioned, the non-decaying mode is
$$\mu _{\mathrm{III}}=Ca(\eta ).$$
(89)
The coefficient $`C`$ is fixed by the matching of the mode function when the mode crosses the horizon at $`\eta _2`$ To be more precise, we should take the decaying mode in Region III into account and match both $`\mu `$ and $`\mu ^{}`$ at time $`\eta _2`$. This only changes the result by an unimportant constant of order one.. One gets
$$C=\mu _{\mathrm{II}}(\eta _2)\frac{|\eta _2|}{l_0}=\frac{2\pi }{n}\frac{\mu _{\mathrm{II}}(\eta _2)}{l_0}.$$
(90)
Therefore, regardless of the choice of the sign of the initial conditions, we have $`|C|1/n^{3/2}`$ and as a result
$$n^3P_\mathrm{\Phi }n^0.$$
(91)
We see that, when $`\beta =2`$, the final answer is not changed compared to what is obtained without the modification of the dispersion relation, i.e. we get a scale invariant spectrum $`n_\mathrm{S}=1`$ \[see Eq. (11)\].
We now discuss different initial conditions. We adopt the ‘instantaneaous Minkowski’ initial conditions given by Eqns. (62). Of course, the form of the solution in Region I is still the same but, now, the coefficients $`A_1`$ and $`A_2`$ are different. The exact expressions for theses coefficients can now be written as
$`A_1`$ $`=`$ $`\pm \sqrt{{\displaystyle \frac{n}{2}}}{\displaystyle \frac{i}{x_2x_1}}|\eta _\mathrm{i}|^{1x_1}\left(1{\displaystyle \frac{ix_2}{n|\eta _\mathrm{i}|}}\right),`$ (92)
$`A_2`$ $`=`$ $`\pm \sqrt{{\displaystyle \frac{n}{2}}}{\displaystyle \frac{i}{x_1x_2}}|\eta _\mathrm{i}|^{1x_2}\left(1{\displaystyle \frac{ix_1}{n|\eta _\mathrm{i}|}}\right).`$ (93)
In the limit when the parameter $`ϵ`$ is small, an expansion of the previous expressions leads to the following formulas
$`A_1{\displaystyle \frac{1}{2\sqrt{2n}}}|\eta _\mathrm{i}|^{2i\pi /ϵ1/2},A_2{\displaystyle \frac{1}{2\sqrt{2n}}}|\eta _\mathrm{i}|^{2i\pi /ϵ1/2}.`$ (94)
The result does not depend on the choice of the sign of the initial conditions. We see also another crucial difference in comparison with the previous case, see Eq. (74) and (76): this time, the coefficients are of the same order in $`ϵ`$. Therefore, the solution in Region I is now given by a cosine instead of by a pure phase
$$\mu _\mathrm{I}(\eta )=\frac{1}{\sqrt{2n}}\left|\frac{\eta }{\eta _\mathrm{i}}\right|^{1/2}\mathrm{cos}\left(\frac{2\pi }{ϵ}\mathrm{ln}\left|\frac{\eta }{\eta _\mathrm{i}}\right|\right).$$
(95)
The solution in Region II is still given by plane waves. The matching at time $`\eta _1`$ permits the calculation of the coefficients $`B_1`$ and $`B_2`$. They read
$`B_1`$ $`=`$ $`{\displaystyle \frac{1}{2n}}\sqrt{{\displaystyle \frac{\pi }{ϵ|\eta _\mathrm{i}|}}}\mathrm{exp}\left(in\eta _1{\displaystyle \frac{2\pi i}{ϵ}}\mathrm{ln}\left|{\displaystyle \frac{\eta _1}{\eta _\mathrm{i}}}\right|\right),`$ (96)
$`B_2`$ $`=`$ $`{\displaystyle \frac{1}{2n}}\sqrt{{\displaystyle \frac{\pi }{ϵ|\eta _\mathrm{i}|}}}\mathrm{exp}\left(in\eta _1+{\displaystyle \frac{2\pi i}{ϵ}}\mathrm{ln}\left|{\displaystyle \frac{\eta _1}{\eta _\mathrm{i}}}\right|\right).`$ (97)
Again, there is an important difference in comparison with the previous case: both coefficients are now non vanishing. The mode function in Region II can be expressed as
$$\mu _{\mathrm{II}}(\eta )=\frac{1}{n}\sqrt{\frac{\pi }{ϵ|\eta _\mathrm{i}|}}\mathrm{cos}\left(n\eta n\eta _1+\frac{2\pi }{ϵ}\mathrm{ln}\left|\frac{\eta _1}{\eta _\mathrm{i}}\right|\right).$$
(98)
The function is proportional to $`1/n`$ instead of $`1/\sqrt{n}`$. The determination of the constant $`C`$ proceeds as previously and leads to the spectrum
$$n^3P_\mathrm{\Phi }n^1\mathrm{cos}^2\left(\frac{2\pi }{ϵ}+\frac{2\pi }{ϵ}\mathrm{ln}\left|\frac{2\pi }{n\eta _\mathrm{i}}\right|\right).$$
(99)
A few remarks are in order here. Firstly, the difference between (91) and (99) demonstrates that the final result does depend on the choice of the initial conditions. Secondly, the spectral index is now modified and is $`n_\mathrm{S}=0`$ instead of $`n_\mathrm{S}=1`$ previously. Thirdly, oscillations in the spectrum are present. If $`n_1`$ and $`n_2`$ are two wave numbers such that the argument of the cosine differs by a factor $`2\pi p`$ where $`p`$ is an integer then one has $`n_2/n_1=\mathrm{exp}(pϵ)`$. This means that unless $`p`$ is comparable to $`ϵ^1`$, $`n_1`$ and $`n_2`$ are almost equal. Therefore, the oscillations are very rapid.
### B The Corley/Jacobson case
With the dispersion relation (22), the equation of motion becomes
$$\mu ^{\prime \prime }+\mu \left[n^2+n^2\underset{q=1}{\overset{m}{}}\frac{b_q}{(2\pi )^{2q}}\left(\frac{ϵn}{|\eta |^{1+\beta }}\right)^{2q}\frac{a^{\prime \prime }}{a}\right]=0.$$
(100)
This equation is valid for any scale factor of the form $`a(\eta )=l_0|\eta |^{1+\beta }`$. Unlike in Unruh’s case, we do not need to specify $`\beta =2`$.
We now need to discuss the form of the solution in Region I. This crucially depends on the sign of the coefficient $`b_m`$. In the regime we are interested in, i.e. $`l_\mathrm{C}\lambda (\eta _\mathrm{i})`$, one can retain only the dominate term and the dispersion relation can be written as
$$n_{\mathrm{eff}}^2n^2+n^2b_m\left(\frac{l_\mathrm{C}}{\lambda }\right)^{2m}.$$
(101)
This means that if $`b_m`$ is positive, the dispersion relation remains real. If $`b_m`$ is negative the situation is more complicated. For very small value of $`|b_m|`$, the dispersion relation can remain real even in the regime $`l_\mathrm{C}\lambda (\eta _\mathrm{i})`$. However, it seems a bit artificial to fine-tune the value of $`|b_m|`$ such that this actually happens. Without this fine-tuning the dispersion relation certainly becomes complex. This last property should not be considered as a surprise. Indeed there exist many situations in Physics where complex dispersion relations appear. This is for example the case in Hydrodynamics when one describes the damping of a sound wave in a fluid due to viscosity . Then, the dispersion relation is given by $`k=\omega /c+ia\omega ^2`$ where $`a`$ is a factor which depends on the viscosity coefficients. In Cosmology, other examples are Silk damping or damping of density perturbations due to neutrino decoupling . In this paper, we choose to analyze both cases and write $`b_ms|b_m|`$ with $`s=\pm 1`$. Then, from Eqns. (50) and (61), the quantities $`\mu _\mathrm{I}(\eta _\mathrm{i})`$ and $`\mu _\mathrm{I}^{}(\eta _\mathrm{i})`$ take the form
$`\mu _\mathrm{I}(\eta _\mathrm{i})`$ $`=`$ $`{\displaystyle \frac{s^{1/4}}{\sqrt{2b\gamma }}}|\eta _i|^{1/2b/2},`$ (102)
$`\mu _\mathrm{I}^{}(\eta _\mathrm{i})`$ $`=`$ $`\pm is^{1/4}\sqrt{{\displaystyle \frac{b\gamma }{2}}}|\eta _i|^{1/2+b/2},`$ (103)
where we have defined $`b`$ and $`\gamma `$ (not to be confused with the function $`\gamma `$ used in Section II) by the following expressions
$$b1m(1+\beta ),\gamma \frac{\sqrt{|b_m|}}{b(2\pi )^m}ϵ^mn^{m+1}.$$
(104)
From the expressions (102) and (103), we deduce
$$\mu _\mathrm{I}^{}(\eta _\mathrm{i})/\mu _\mathrm{I}(\eta _\mathrm{i})=\pm is^{1/2}b\gamma |\eta _i|^{b1}.$$
(105)
This ratio will out turn to be important in the calculation of the various coefficients determined by the matching procedure. To go further, we need to treat the cases $`s=\pm 1`$ separately.
#### 1 The case $`s=1`$, $`b_m<0`$
In Region I, the equation of motion for the mode function reduces to
$$\mu ^{\prime \prime }+n^2\frac{b_m}{(2\pi )^{2m}}\left(\frac{ϵn}{|\eta |^{1+\beta }}\right)^{2m}\mu =0.$$
(106)
For a negative coefficient $`b_m`$, the exact solution of Eq. (106) can be expressed in terms of modified Bessel functions as follows
$$\mu _\mathrm{I}(\eta )=A_1|\eta |^{1/2}I_\nu (z)+A_2|\eta |^{1/2}K_\nu (z),$$
(107)
where $`\nu 1/(2b)`$ and where the function $`z(\eta )`$ is defined by the following expression $`z(\eta )\gamma |\eta |^b`$. The coefficients $`A_1`$ and $`A_2`$ are determined by the initial conditions given in Eqns. (102) and (103). These coefficients should satisfy the system of equations
$`A_1I_\nu (z_\mathrm{i})+A_2K_\nu (z_\mathrm{i})`$ $`=`$ $`|\eta _\mathrm{i}|^{1/2}\mu _\mathrm{I}(\eta _\mathrm{i}),`$ (108)
$`A_1I_{\nu 1}(z_i)+A_2K_{\nu 1}(z_i)`$ $`=`$ $`{\displaystyle \frac{|\eta _\mathrm{i}|^{1/2b}}{\gamma b}}\mu _\mathrm{I}^{}(\eta _\mathrm{i}),`$ (109)
where $`z_\mathrm{i}`$ denotes the value of $`z(\eta )`$ at time $`\eta =\eta _\mathrm{i}`$. The exact solution for $`A_1`$ and $`A_2`$ can be expressed as
$`A_1`$ $`=`$ $`\gamma |\eta _\mathrm{i}|^{b1/2}\mu _\mathrm{I}(\eta _\mathrm{i})K_{\nu 1}(z_\mathrm{i})`$ (111)
$`\times \left[1{\displaystyle \frac{|\eta _\mathrm{i}|^{1b}}{\gamma b}}{\displaystyle \frac{\mu _\mathrm{I}^{}(\eta _\mathrm{i})}{\mu _\mathrm{I}(\eta _\mathrm{i})}}{\displaystyle \frac{K_\nu (z_\mathrm{i})}{K_{\nu 1}(z_\mathrm{i})}}\right],`$
$`A_2`$ $`=`$ $`\gamma |\eta _\mathrm{i}|^{b1/2}\mu _\mathrm{I}(\eta _\mathrm{i})I_{\nu 1}(z_\mathrm{i})`$ (113)
$`\times \left[1+{\displaystyle \frac{|\eta _\mathrm{i}|^{1b}}{\gamma b}}{\displaystyle \frac{\mu _\mathrm{I}^{}(\eta _\mathrm{i})}{\mu _\mathrm{I}(\eta _\mathrm{i})}}{\displaystyle \frac{I_\nu (z_\mathrm{i})}{I_{\nu 1}(z_\mathrm{i})}}\right].`$
In the derivation of the previous expressions, we used the exact equation: $`(I_\nu K_{\nu 1}+I_{\nu 1}K_\nu )(z)=1/z`$. Since, when $`l_\mathrm{C}\lambda (\eta _i)`$, the argument $`z_\mathrm{i}`$ is large we can now rewrite these equations using the asymptotic formulas for Bessel functions of large arguments . Notice that it is necessary to go to the second order in the expansion of the modified Bessel functions. We obtain
$`A_1`$ $``$ $`\left({\displaystyle \frac{\pi }{2}}\right)^{1/2}\gamma ^{1/2}\mu _\mathrm{I}(\eta _\mathrm{i})|\eta _\mathrm{i}|^{b/21/2}e^{z_\mathrm{i}}`$ (115)
$`\times \left[1\pm 1\pm {\displaystyle \frac{2\nu 1}{2\gamma }}|\eta _\mathrm{i}|^b\right],`$
$`A_2`$ $``$ $`\left({\displaystyle \frac{1}{2\pi }}\right)^{1/2}\gamma ^{1/2}\mu _\mathrm{I}(\eta _\mathrm{i})|\eta _\mathrm{i}|^{b/21/2}e^{z_\mathrm{i}}`$ (117)
$`\times \left[11{\displaystyle \frac{12\nu }{2\gamma }}|\eta _\mathrm{i}|^b\right].`$
For the sake of completeness, we pursue the calculation for both choices of the sign of the initial conditions. Let us again use an index “u” for the upper choice and “l” for the lower choice. We obtain:
$`A_1^u`$ $`=`$ $`2\left({\displaystyle \frac{\pi }{2}}\right)^{1/2}\gamma ^{1/2}\mu _\mathrm{I}(\eta _\mathrm{i})|\eta _\mathrm{i}|^{b/21/2}e^{z_\mathrm{i}},`$ (118)
$`A_2^u`$ $`=`$ $`\left({\displaystyle \frac{1}{2\pi }}\right)^{1/2}\mu _\mathrm{I}(\eta _\mathrm{i})|\eta _\mathrm{i}|^{b/21/2}{\displaystyle \frac{2\nu 1}{2\gamma ^{1/2}}}e^{z_\mathrm{i}},`$ (119)
$`A_1^l`$ $`=`$ $`\left({\displaystyle \frac{\pi }{2}}\right)^{1/2}\mu _\mathrm{I}(\eta _\mathrm{i})|\eta _\mathrm{i}|^{b/21/2}{\displaystyle \frac{12\nu }{2\gamma ^{1/2}}}e^{z_\mathrm{i}},`$ (120)
$`A_2^l`$ $`=`$ $`2\left({\displaystyle \frac{1}{2\pi }}\right)^{1/2}\gamma ^{1/2}\mu _\mathrm{I}(\eta _\mathrm{i})|\eta _\mathrm{i}|^{b/21/2}e^{z_\mathrm{i}}.`$ (121)
The exponential factor always determines the behaviour of the coefficients for any power of $`|\eta _\mathrm{i}|`$. This implies $`A_1^uA_1^l0`$. We also see the following crucial effect: it turns out that for one choice of the sign of the derivative the first term in the squared bracket in Eqns. (115) and (117) cancels whereas for the other choice it is no longer the case. This has as a consequence that the dependence on $`\gamma `$ is not the same. Since $`\gamma `$ depends on $`n`$, the $`n`$ dependence of $`A_2^u`$ and $`A_2^l`$ is not the same. We have $`A_2^u\gamma ^1A_2^l`$.
The second step of the calculation is to perform the matching of the solutions at the time $`\eta =\eta _1`$. This will allow us to express the coefficients $`B_1`$ and $`B_2`$ in terms of the coefficients $`A_1`$ and $`A_2`$. In Region II, the solution is given by plane waves. Therefore, the coefficients $`B_1`$ and $`B_2`$ are now solutions of the equations
$`B_1e^{in\eta _1}+B_2e^{in\eta _1}`$ $`=`$ $`A_1|\eta _1|^{1/2}I_\nu (z_1)`$ (123)
$`+A_2|\eta _1|^{1/2}K_\nu (z_1),`$
$`B_1e^{in\eta _1}B_2e^{in\eta _1}`$ $`=`$ $`{\displaystyle \frac{\gamma b}{in}}A_1|\eta _1|^{b1/2}I_{\nu 1}(z_1)`$ (125)
$`+{\displaystyle \frac{\gamma b}{in}}A_2|\eta _1|^{b1/2}K_{\nu 1}(z_1),`$
where $`z_1`$ is the value of the function $`z(\eta )`$ at $`\eta =\eta _1`$. The exact solution of this system of equations can be easily found and reads
$`e^{in\eta _1}B_1`$ $`=`$ $`{\displaystyle \frac{A_1}{2}}|\eta _1|^{1/2}I_\nu (z_1)\left[1+{\displaystyle \frac{i\gamma b}{n}}|\eta _1|^{b1}{\displaystyle \frac{I_{\nu 1}(z_1)}{I_\nu (z_1)}}\right]`$ (127)
$`+{\displaystyle \frac{A_2}{2}}|\eta _1|^{1/2}K_\nu (z_1)\left[1{\displaystyle \frac{i\gamma b}{n}}|\eta _1|^{b1}{\displaystyle \frac{K_{\nu 1}(z_1)}{K_\nu (z_1)}}\right],`$
$`e^{in\eta _1}B_2`$ $`=`$ $`{\displaystyle \frac{A_1}{2}}|\eta _1|^{1/2}I_\nu (z_1)\left[1{\displaystyle \frac{i\gamma b}{n}}|\eta _1|^{b1}{\displaystyle \frac{I_{\nu 1}(z_1)}{I_\nu (z_1)}}\right]`$ (129)
$`+{\displaystyle \frac{A_2}{2}}|\eta _1|^{1/2}K_\nu (z_1)\left[1+{\displaystyle \frac{i\gamma b}{n}}|\eta _1|^{b1}{\displaystyle \frac{K_{\nu 1}(z_1)}{K_\nu (z_1)}}\right].`$
Much simpler (approximate) formulas can be obtained if one notices that the argument of the Bessel function is a big number $`z_1=\gamma |\eta _1|^b1`$, essentially because $`ϵ`$ is a small number in realistic cases. A very simple estimate allows us to quickly check the validity of this approximation. We take $`m=1`$, $`|b_1|=1`$, $`\beta =2.2`$ which would correspond to a spectral index of $`n_\mathrm{S}=0.6`$ for power law inflation and $`ϵ=10^5`$ as already discussed in the previous section. We can then estimate $`z_1`$ for $`n=4\pi `$ which corresponds to the mode which re-enters horizon today and which consequently mainly determines the value of the CMB quadrupole anisotropy. We find $`z_14.7\times 10^4`$. Therefore, we can again use the asymptotic behaviour of the Bessel function to simplify the previous equations. Putting all these ingredients together, we find In order to be able to neglect the terms proportional to $`A_1`$, we make use of the fact that $`|\eta _i||\eta _1|`$.
$`B_1`$ $``$ $`{\displaystyle \frac{A_2}{2}}\left({\displaystyle \frac{\pi }{2\gamma }}\right)^{1/2}|\eta _1|^{1/2b/2}e^{in\eta _1z_1i\frac{\pi }{4}},`$ (130)
$`B_2`$ $``$ $`{\displaystyle \frac{A_2}{2}}\left({\displaystyle \frac{\pi }{2\gamma }}\right)^{1/2}|\eta _1|^{1/2b/2}e^{in\eta _1z_1+i\frac{\pi }{4}}.`$ (131)
As a consequence, the solution in Region II can be written as
$$\mu _{\mathrm{II}}(\eta )=A_2\left(\frac{\pi }{2\gamma }\right)^{1/2}|\eta _1|^{1/2b/2}e^{z_1}\mathrm{cos}\left(n\eta n\eta _1\frac{\pi }{4}\right).$$
(132)
The last step of the calculation is to perform the matching at $`\eta =\eta _2`$ when the mode leaves the Hubble radius (boundary between Region II and Region III). As already mentioned, in Region III, the non-decaying solution is the super-Hubble function given by
$$\mu _{\mathrm{III}}(\eta )=Ca(\eta ).$$
(133)
Repeating the same procedure as for Unruh’s case, the spectra can easily be calculated and read
$`n^3P_\mathrm{\Phi }^l`$ $``$ $`n^{2\beta +4}e^{2(z_\mathrm{i}z_1)}\mathrm{cos}^2\left(n\eta _2n\eta _1{\displaystyle \frac{\pi }{4}}\right),`$ (134)
$`n^3P_\mathrm{\Phi }^u`$ $``$ $`n^{2\beta +22m}e^{2(z_\mathrm{i}z_1)}\mathrm{cos}^2\left(n\eta _2n\eta _1{\displaystyle \frac{\pi }{4}}\right).`$ (135)
We see that the spectrum depends explicitly on the initial conditions chosen. We can check that the tilt is correct by noticing that $`n^3P_\mathrm{\Phi }A_2^2`$, $`\gamma n^{m+1}`$ and using the relation between $`A_2^u`$ and $`A_2^l`$ already mentioned. From now on, we concentrate on the lower case which corresponds to an unmodified tilt and study the expression of the corresponding spectrum in more details (for convenience, we drop the subscript ‘l’). First, as mentioned above, we see that the power-law part is not modified in comparison with the usual case, i.e. the spectral index is still $`n_\mathrm{S}=2\beta +5`$. Secondly, there are oscillations in the spectrum since the argument of the cosine can be written as (considering for simplicity that $`|b_m|=1`$)
$$2\pi |1+\beta |\left(\frac{ϵ}{2\pi }\right)^{1/(1+\beta )}n^{(2+\beta )/(1+\beta )}\frac{\pi }{4}.$$
(136)
However, contrary to Unruh’s case with Minkowski initial conditions, no logarithmic dependence is present. Interestingly enough, for $`\beta =2`$, the oscillations disappear. The most important part concerns the exponential factor. The factor $`z_\mathrm{i}z_1`$ is equal to $`z_\mathrm{i}z_1=\gamma |\eta _\mathrm{i}|^b(1|\eta _1|^b/|\eta _\mathrm{i}^b)\gamma |\eta _\mathrm{i}|^b=z_\mathrm{i}`$ since we have $`|\eta _\mathrm{i}||\eta _1|`$. The factor $`z_\mathrm{i}`$ can be re-written in such a way that the dependence on $`n`$ is explicit
$$z_\mathrm{i}=\frac{\sqrt{|b_m|}}{b(2\pi )^m}ϵ^m|\eta _\mathrm{i}|^{1m(1+\beta )}n^{m+1}.$$
(137)
The important factor in this expression is $`ϵ^m|\eta _\mathrm{i}|^{1m(1+\beta )}`$ since the others ones are of order one. It can be re-written as
$$ϵ^m|\eta _\mathrm{i}|^{1m(1+\beta )}=\left[\frac{l_\mathrm{C}}{a(\eta _\mathrm{i})}\right]^m|\eta _\mathrm{i}|,$$
(138)
and must be considered as large since $`|\eta _\mathrm{i}|1`$ and $`l_\mathrm{C}/a(\eta _\mathrm{i})l_\mathrm{C}/\lambda (\eta _\mathrm{i})1`$, at least for wavenumbers not too different from $`2\pi `$. This means that the influence of the exponential factor is dominant and is responsible for a huge increase of the spectrum at large $`n`$. This is illustrated if we write the spectrum for $`\beta =2`$ and $`m=1`$
$$n^3P_\mathrm{\Phi }e^{An^2},$$
(139)
where $`A1`$. Such a spectrum is almost certainly in contradiction with observations.
We end this subsection with the calculation of the spectrum in the case where the initial state is the minimum energy density state. We restart from the exact expressions for the coefficients $`A_1`$ and $`A_2`$, see Eqns. (111) and (113). Using Eqns. (62), we have
$$\frac{|\eta _\mathrm{i}|^{1b}}{\gamma b}\frac{\mu _\mathrm{I}^{}(\eta _\mathrm{i})}{\mu _\mathrm{I}(\eta _\mathrm{i})}=\pm \frac{i}{\sqrt{b_m}}\left[\frac{\lambda (\eta _\mathrm{i})}{l_\mathrm{C}}\right]^m1,$$
(140)
since, initially, $`l_\mathrm{C}\lambda (\eta _\mathrm{i})`$. As a consequence, we can derive a compact approximate expression for the coefficients $`A_1`$ and $`A_2`$
$$A_10,A_2\frac{1}{\sqrt{2\pi }}\gamma ^{1/2}|\eta _\mathrm{i}|^{b/21/2}\mu _I(\eta _\mathrm{i})e^{z_\mathrm{i}}.$$
(141)
Notice that these formulas are valid for any choice of the sign of $`\mu _\mathrm{I}^{}(\eta _\mathrm{i})`$. The rest of the calculation proceeds as above and leads to ($`|b_m|=1`$)
$`n^3P_\mathrm{\Phi }`$ $`=`$ $`n^{2\beta +4+m}e^{An^{m+1}}`$ (143)
$`\times \mathrm{cos}^2\left[2\pi |1+\beta |\left({\displaystyle \frac{ϵ}{2\pi }}\right)^{\frac{1}{1+\beta }}n^{\frac{2+\beta }{1+\beta }}{\displaystyle \frac{\pi }{4}}\right].`$
The main difference in comparison with the spectrum of the previous section is the presence of a modified tilt. The spectral index is now given by $`n_\mathrm{S}=2\beta +5+m`$.
#### 2 The case $`s=1`$, $`b_m>0`$
When the dispersion relation is real, the solution in Region I can be expressed in terms of usual Bessel functions
$$\mu _\mathrm{I}(\eta )=A_1|\eta |^{1/2}J_\nu (z)+A_2|\eta |^{1/2}J_\nu (z),$$
(144)
where $`\nu `$ and $`z(\eta )`$ have already been defined previously. The coefficients $`A_1`$ and $`A_2`$ are now solutions of the following system of equations
$`A_1J_\nu (z_\mathrm{i})+A_2J_\nu (z_\mathrm{i})`$ $`=`$ $`\mu _\mathrm{I}(\eta _\mathrm{i})|\eta _\mathrm{i}|^{1/2},`$ (145)
$`A_1J_\nu (z_\mathrm{i})+A_2J_\nu (z_\mathrm{i})`$ $`=`$ $`{\displaystyle \frac{\mu _\mathrm{I}^{}(\eta _\mathrm{i})}{\gamma b}}|\eta _\mathrm{i}|^{1/2b}.`$ (146)
Using the relation expressing the Wronskian $`[J_\nu J_{\nu 1}+J_{\nu +1}J_\nu ](z)=2\mathrm{sin}[\pi /(2b)]/(\pi z)`$, and performing some straightforward algebraic manipulations, exact expressions can be easily found. They read
$`A_1`$ $`=`$ $`{\displaystyle \frac{\pi \gamma }{2\mathrm{sin}(\pi \nu )}}|\eta _\mathrm{i}|^{b1/2}\mu _\mathrm{I}(\eta _\mathrm{i})J_{1\nu }(z_\mathrm{i})`$ (148)
$`\times [1{\displaystyle \frac{|\eta _\mathrm{i}|^{1b}}{\gamma b}}{\displaystyle \frac{\mu _\mathrm{I}^{}(\eta _\mathrm{i})}{\mu _\mathrm{I}(\eta _\mathrm{i})}}{\displaystyle \frac{J_\nu (z_\mathrm{i})}{J_{1\nu }(z_\mathrm{i})}}],`$
$`A_2`$ $`=`$ $`{\displaystyle \frac{\pi \gamma }{2\mathrm{sin}(\pi \nu )}}|\eta _\mathrm{i}|^{b1/2}\mu _\mathrm{I}(\eta _\mathrm{i})J_{\nu 1}(z_\mathrm{i})`$ (150)
$`\times [1+{\displaystyle \frac{|\eta _\mathrm{i}|^{1b}}{\gamma b}}{\displaystyle \frac{\mu _\mathrm{I}^{}(\eta _\mathrm{i})}{\mu _\mathrm{I}(\eta _\mathrm{i})}}{\displaystyle \frac{J_\nu (z_\mathrm{i})}{J_{\nu 1}(z_\mathrm{i})}}].`$
These expressions are not valid if $`\nu =1/(2b)`$ is an integer and this particular case must be treated separately. In this article, we assume that this does not happen. Since the solution in Region II is still given by plane waves, the derivation of exact expressions for the coefficients $`B_1`$ and $`B_2`$ proceeds as before. Explicit matching of the mode function and of its derivative leads to
$`e^{in\eta _1}B_1`$ $`=`$ $`{\displaystyle \frac{A_1}{2}}|\eta _1|^{1/2}J_\nu (z_1)\left[1+i{\displaystyle \frac{\gamma b}{n}}|\eta _1|^{b1}{\displaystyle \frac{J_{\nu 1}(z_1)}{J_\nu (z_1)}}\right]`$ (152)
$`+{\displaystyle \frac{A_2}{2}}|\eta _1|^{1/2}J_\nu (z_1)\left[1i{\displaystyle \frac{\gamma b}{n}}|\eta _1|^{b1}{\displaystyle \frac{J_{\nu +1}(z_1)}{J_\nu (z_1)}}\right],`$
$`e^{in\eta _1}B_2`$ $`=`$ $`{\displaystyle \frac{A_1}{2}}|\eta _1|^{1/2}J_\nu (z_1)\left[1i{\displaystyle \frac{\gamma b}{n}}|\eta _1|^{b1}{\displaystyle \frac{J_{\nu 1}(z_1)}{J_\nu (z_1)}}\right]`$ (154)
$`+{\displaystyle \frac{A_2}{2}}|\eta _1|^{1/2}J_\nu (z_1)\left[1+i{\displaystyle \frac{\gamma b}{n}}|\eta _1|^{b1}{\displaystyle \frac{J_{\nu +1}(z_1)}{J_\nu (z_1)}}\right].`$
Having all the relevant exact expressions at our disposal we can now start to do some approximations based on the fact that $`z_\mathrm{i}`$ is a big number. For convenience, we introduce two new definitions \[not to be confused with the functions $`x(\eta )`$ and $`y(\eta )`$ introduced in section IV-B\]
$$x(\eta )z(\eta )+\frac{\pi \nu }{2}\frac{\pi }{4},y(\eta )z(\eta )\frac{\pi \nu }{2}\frac{\pi }{4}.$$
(155)
Then, using the expressions of the Bessel functions for large arguments , we find
$`A_1`$ $``$ $`i\left({\displaystyle \frac{\pi \gamma }{2}}\right)^{1/2}|\eta _\mathrm{i}|^{b/21/2}{\displaystyle \frac{\mu _\mathrm{I}(\eta _\mathrm{i})}{\mathrm{sin}(\pi \nu )}}e^{\pm ix_\mathrm{i}},`$ (156)
$`A_2`$ $``$ $`\pm i\left({\displaystyle \frac{\pi \gamma }{2}}\right)^{1/2}|\eta _\mathrm{i}|^{b/21/2}{\displaystyle \frac{\mu _\mathrm{I}(\eta _\mathrm{i})}{\mathrm{sin}(\pi \nu )}}e^{\pm iy_\mathrm{i}},`$ (157)
where $`x_\mathrm{i}x(\eta _\mathrm{i})`$ and $`y_\mathrm{i}y(\eta _\mathrm{i})`$. The correct matching time is $`|\eta _1|=[nϵ/(2\pi )]^{1/(1+\beta )}b_m^{1/[2m(1+\beta )]}`$, see Ref. , and is equal to the time at which $`\lambda =l_\mathrm{C}`$ if $`b_m=1`$. In the following, for simplicity, we consider $`b_m=1`$. The coefficients $`B_1`$ and $`B_2`$ can be expressed as
$`B_1`$ $``$ $`\left({\displaystyle \frac{1}{2\pi \gamma }}\right)^{1/2}|\eta _1|^{1/2b/2}e^{in\eta _1}(A_1\mathrm{cos}y_1`$ (159)
$`iA_1\mathrm{sin}y_1+A_2\mathrm{cos}x_1iA_2\mathrm{sin}x_1),`$
$`B_2`$ $``$ $`\left({\displaystyle \frac{1}{2\pi \gamma }}\right)^{1/2}|\eta _1|^{1/2b/2}e^{in\eta _1}(A_1\mathrm{cos}y_1`$ (161)
$`+iA_1\mathrm{sin}y_1+A_2\mathrm{cos}x_1+iA_2\mathrm{sin}x_1),`$
where $`x_1x(\eta _1)`$ and $`y_1=y(\eta _1)`$. Our next move is to replace the expressions of $`A_1`$ and $`A_2`$, see Eqns. (156) and (157), in the previous formula. This leads to
$`B_1`$ $`=`$ $`i{\displaystyle \frac{\mu _\mathrm{I}(\eta _\mathrm{i})e^{in\eta _1}}{2\mathrm{sin}(\pi \nu )}}\left|{\displaystyle \frac{\eta _1}{\eta _\mathrm{i}}}\right|^{1/2b/2}e^{\pm ix_\mathrm{i}}(\mathrm{cos}y_1`$ (163)
$`i\mathrm{sin}y_1e^{i\pi \nu }\mathrm{cos}x_1+ie^{i\pi \nu }\mathrm{sin}x_1),`$
$`B_2`$ $`=`$ $`i{\displaystyle \frac{\mu _\mathrm{I}(\eta _\mathrm{i})e^{in\eta _1}}{2\mathrm{sin}(\pi \nu )}}\left|{\displaystyle \frac{\eta _1}{\eta _\mathrm{i}}}\right|^{1/2b/2}e^{\pm ix_\mathrm{i}}(\mathrm{cos}y_1`$ (165)
$`+i\mathrm{sin}y_1e^{i\pi \nu }\mathrm{cos}x_1ie^{i\pi \nu }\mathrm{sin}x_1).`$
Then, the mode function at time $`\eta =\eta _2`$ (which is the relevant quantity for the determination of the constant $`C`$) can be expressed as
$$\mu _{\mathrm{II}}(\eta _2)=i\frac{\mu _\mathrm{I}(\eta _\mathrm{i})}{2\mathrm{sin}(\pi \nu )}\left|\frac{\eta _1}{\eta _\mathrm{i}}\right|^{1/2b/2}e^{\pm ix_\mathrm{i}}.$$
(166)
From this equation, the expression of the spectrum can be easily established and reads
$$n^3P_\mathrm{\Phi }n^{2\beta +4}.$$
(167)
Let us analyze this spectrum in more detail. The first remark is that the tilt is unchanged and that the spectral index is given by the usual expression $`n_\mathrm{S}=2\beta +5`$. The second remark is that the exponential dependence has disappeared. This is due to the fact that, for $`s=1`$, this factor becomes a pure phase. We recover the usual result as pointed out in Ref. .
Let us finally turn to the case where the initial conditions are those which correspond to the instantaneous Minkowski vacuum. Restarting from the exact expressions for the coefficients $`A_1`$ and $`A_2`$, see Eqns. (148) and (150), and using Eqns. (62), we find
$`A_1`$ $``$ $`\left({\displaystyle \frac{\pi \gamma }{2}}\right)^{1/2}|\eta _\mathrm{i}|^{b/21/2}{\displaystyle \frac{\mu _\mathrm{I}(\eta _\mathrm{i})}{\mathrm{sin}(\pi \nu )}}\mathrm{sin}x_\mathrm{i},`$ (168)
$`A_2`$ $``$ $`\left({\displaystyle \frac{\pi \gamma }{2}}\right)^{1/2}|\eta _\mathrm{i}|^{b/21/2}{\displaystyle \frac{\mu _\mathrm{I}(\eta _\mathrm{i})}{\mathrm{sin}(\pi \nu )}}\mathrm{sin}y_\mathrm{i}.`$ (169)
Inserting these equations into the exact formulas giving the coefficients $`B_1`$ and $`B_2`$, we obtain
$`B_1`$ $``$ $`{\displaystyle \frac{\mu _\mathrm{I}(\eta _\mathrm{i})e^{in\eta _1}}{2\mathrm{sin}(\pi \nu )}}\left|{\displaystyle \frac{\eta _1}{\eta _\mathrm{i}}}\right|^{b/2+1/2}`$ (172)
$`\times (\mathrm{sin}x_\mathrm{i}\mathrm{cos}y_1i\mathrm{sin}x_\mathrm{i}\mathrm{sin}y_1\mathrm{sin}y_\mathrm{i}\mathrm{cos}x_1`$
$`+i\mathrm{sin}y_\mathrm{i}\mathrm{sin}x_1),`$
$`B_2`$ $``$ $`{\displaystyle \frac{\mu _\mathrm{I}(\eta _\mathrm{i})e^{in\eta _1}}{2\mathrm{sin}(\pi \nu )}}\left|{\displaystyle \frac{\eta _1}{\eta _\mathrm{i}}}\right|^{b/2+1/2}`$ (175)
$`\times (\mathrm{sin}x_\mathrm{i}\mathrm{cos}y_1+i\mathrm{sin}x_\mathrm{i}\mathrm{sin}y_1\mathrm{sin}y_\mathrm{i}\mathrm{cos}x_1`$
$`i\mathrm{sin}y_\mathrm{i}\mathrm{sin}x_1).`$
We are now in a position where we can write the expression of the mode function at time $`\eta =\eta _2`$. It reads
$`\mu _{\mathrm{II}}(\eta _2)`$ $`=`$ $`{\displaystyle \frac{\mu _\mathrm{I}(\eta _\mathrm{i})}{2\mathrm{sin}(\pi \nu )}}\left|{\displaystyle \frac{\eta _1}{\eta _\mathrm{i}}}\right|^{b/2+1/2}[\stackrel{~}{B}(n)e^{in(\eta _2\eta _1)}`$ (177)
$`+\stackrel{~}{B}^{}(n)e^{in(\eta _2\eta _1)}],`$
where the function $`\stackrel{~}{B}(n)`$ is defined by
$`\stackrel{~}{B}(n)`$ $``$ $`\mathrm{sin}x_\mathrm{i}\mathrm{cos}y_1i\mathrm{sin}x_\mathrm{i}\mathrm{sin}y_1\mathrm{sin}y_\mathrm{i}\mathrm{cos}x_1`$ (179)
$`+i\mathrm{sin}y_\mathrm{i}\mathrm{sin}x_1.`$
Then, one can write $`\stackrel{~}{B}(n)`$ as $`\stackrel{~}{B}(n)|\stackrel{~}{B}|e^{i\psi }`$ and define $`\overline{B}(n)`$ as $`\overline{B}(n)|\stackrel{~}{B}(n)|\mathrm{cos}(n\eta _2n\eta _1+\psi )`$. It follows that the spectrum can be written as
$$n^3P_\mathrm{\Phi }n^{2\beta +4+m}|\overline{B}(n)|^2.$$
(180)
The spectral index is given by $`n_\mathrm{S}=2\beta +5+m`$, i.e. it differs from the standard one but is equal to the spectral index obtained in the case $`s=1`$ for instantaneous Minkowski initial conditions. The factor $`|\overline{B}(n)|^2`$ is of order one and produces a complicated oscillatory pattern.
In conclusion, the resulting spectrum in the case of the Corley/Jacobson dispersion relation is very different from the usual spectrum calculated using an unmodified dispersion relation, and different from what is obtained using Unruh’s relation, even for initial conditions which minimize the energy.
## VI Discussion and Conclusions
We have studied the dependence of the predictions of inflationary cosmology for the spectrum of fluctuations on hidden assumptions about super-Planck-scale physics. The motivation for our work is that in most current models of inflation, the period of exponential expansion lasts so long that at the beginning of inflation, scales of cosmological interest today had a physical wavelength much smaller than the Planck length, and the theories used to compute the spectrum of fluctuations are known to break down on these scales.
We studied the problem by replacing the dispersion relation of a free field theory which is used to compute the spectrum in the standard approaches by a modified dispersion relation, the modifications only being important on length scales smaller than a cutoff length $`l_\mathrm{C}`$ (which we expect to be given by the Planck length). We considered two classes of dispersion relations, based on the ones considered by Unruh (Class A) and by Corley and Jacobson (Class B), respectively, in their studies of the trans-Planckian problem of black hole physics. Admittedly, modifying the physics in this way is a very ad hoc way of taking into account possible effects of super-Planck-scale physics, chosen for mathematical simplicity. We do not want to introduce mode-mode coupling in order to keep the computations simple. However, in order to demonstrate that there is a possible problem for the robustness of the usual predictions of inflation it is sufficient to construct one example of a modified theory which leads to different predictions.
For a non-standard dispersion relation the choice of initial state becomes more difficult. We considered two choices, both of which coincide with the usual initial state in the case of the standard dispersion relation. The first (and better motivated) choice is the state which minimizes the energy density, the second choice is the naive generalization of the ‘local Minkowski vacuum’.
We have shown that for Class A dispersion relations the usual predictions of inflationary cosmology are recovered (in the case of exponential inflation) if the initial state minimizes the energy density. In particular, the spectrum of fluctuations is scale-invariant. If the initial state is chosen to be the ‘local Minkowski vacuum’, then the resulting spectrum has a tilt and superimposed oscillations.
In contrast, for Class B dispersion relations and an initial state which minimizes the energy density, the resulting spectrum of fluctuations is in general not scale-invariant. The precise nature of the spectrum depends sensitively on whether the dispersion relation turns complex or remains real. In the complex case, the spectrum is characterized by an exponential factor (more power in the blue, i.e. $`n_\mathrm{S}>1`$), a tilt (compared to the “standard” predictions) which depends on the precise initial conditions, and superimposed oscillations. The exponent, the tilt, and the precise oscillatory pattern depend on the specific member of the class of dispersion relations chosen. For a spectrum which remains real, the usual result is unchanged.
The reason why for Class A dispersion relations the usual predictions of inflation are maintained is that the time evolution during the period when the mode wavelength is smaller than the cutoff scale is adiabatic. This emerges from our calculations, but an intuitive way of understanding the result is that at all times the effective frequency of the mode is larger than the Hubble rate and the initial vacuum state therefore adjusts itself adiabatically to track the instantaneous vacuum state, thus leading to the same state at time $`\eta _1`$ as in the theory with unmodified dispersion relation <sup>§</sup><sup>§</sup>§We thank Bill Unruh for making this point to us.. For Class B dispersion relations, in contrast, the dispersion relation varies too quickly as a function of time while the scale is smaller than the critical length $`l_\mathrm{C}`$ and hence the evolution is not adiabatic .
Let us now be more quantitative about the previous discussion. In Region I, Eq. (15) can be written as
$$\mu ^{\prime \prime }+n_{\mathrm{eff}}^2\mu =\mu ^{\prime \prime }+a^2(\eta )\omega _{\mathrm{phys}}^2(n,\eta )\mu =0,$$
(181)
where $`\omega _{\mathrm{phys}}`$ is the physical frequency defined by $`\omega _{\mathrm{phys}}(1/a)\sqrt{n^2+a^2\mathrm{\Omega }^2(n,\eta )/l_\mathrm{C}^2}`$. The latter can be considered as constant as long as its characteristic time scale of evolution is small compared to the Hubble time, i.e. as long as we have adiabaticity. Therefore, let us define an “adiabaticity coefficient” $`\alpha `$ according to
$$\alpha (n,\eta )\left|\frac{}{\frac{1}{\omega _{\mathrm{phys}}}\frac{\mathrm{d}\omega _{\mathrm{phys}}}{\mathrm{d}\eta }}\right|,$$
(182)
where we recall that $`a^{}/a`$. When $`\alpha 1`$, adiabaticity is satisfied and Eq. (181) reduces to the equation of motion in the Unruh’s case, Eq. (66). In this situation, we know that the final spectrum is unmodified since there is an exact cancellation of the $`n`$-dependence in the minimizing energy state and in the growth factor before Hubble radius crossing which results in the usual spectrum. The previous argument shows that an unmodified spectrum is expected when $`\alpha 1`$ in the region where the dispersion relation is modified. Let us also note in passing that for the standard case, $`\alpha =1`$, since the time scale of evolution of $`\omega _{\mathrm{phys}}(n,\eta )`$ and of the Hubble rate is the same.
We have calculated the adiabaticity coefficient for the different cases treated in this article. The result is displayed in Fig. (3). When $`\eta `$ goes to $`\mathrm{}`$, we have $`l_\mathrm{C}\lambda `$ whereas $`l_\mathrm{C}\lambda `$ when $`\eta `$ goes to zero.
We see that there exists a clear difference between Unruh’s case and the Corley/Jacobson cases. In the Unruh’s case, $`\alpha `$ goes to infinity when $`l_\mathrm{C}\lambda `$ and adiabaticity is preserved. When $`b_m<0`$, the adiabaticity coefficient reaches zero at the time when $`\omega =0`$ . Then, adiabaticity is progressively re-established. The coefficient $`\alpha `$ goes to infinity and there is a divergence when $`\mathrm{d}\omega /\mathrm{d}\eta =0`$ . In the regime when $`l_\mathrm{C}\lambda `$, $`\alpha `$ goes to one as it should since the various dispersion relations all become similar to the standard one. The previous considerations explain why the final spectrum can be modified in the Corley/Jacobson case with a complex dispersion relation but not in Unruh’s case.
We conclude that it is possible that in models of inflation based consistently on a unified theory at the Planck scale the predictions for fluctuations will not coincide with the usual predictions from our current inflationary Universe models. Generalizing from our results, a crucial issue appears to be whether the evolution of the quantum states corresponding to the fluctuations will be adiabatic on length scales smaller than the Planck scale.
Our results point to the possibility that the interaction between fundamental physics and cosmology may be much richer than hitherto assumed. It is not only a question of if and how fundamental physics leads to inflation; a much richer question is what the specific predictions of the fundamental model of inflation will be, assuming for the sake of argument that such a fundamental model exists.
It is not surprising that super-Planck-scale physics may modify the usual predictions of inflation. One model of the early Universe motivated by string theory, the Pre-Big-Bang Cosmology based on dilaton gravity, leads to a super-exponential period of early evolution in which the Hubble constant is increasing, and where the predicted spectrum of scalar metric fluctuations is not scale-invariant . It would be interesting to analyze the predictions of other models of inflation based on string theory, taking into account the evolution on string scales. One toy model in which this question could be analyzed is the nonsingular Universe based on higher derivative terms in the gravitational action.
In the context of the models studied here, it would be interesting to explore whether the minimum energy density initial state is an attractor in a similar sense that the local Minkowski vacuum is in standard inflationary cosmology .
Note that models of inflation based on a strongly interacting theory (such as the model analyzed in ) do not suffer from the Trans-Planckian problem discussed in this paper. In strongly interacting theories, perturbations are generated at all times at a fixed physical scale, and a scale-invariant spectrum results based on the heuristic arguments mentioned in the Introduction. In such theories, however, the presence of strong interactions makes it hard to calculate the amplitude of the resulting spectrum.
Acknowledgements
We are grateful to Lev Kofman, Dominik Schwarz, Carsten Van de Bruck and in particular Bill Unruh for stimulating discussions and useful comments. We also thank an anonymous referee for useful comments. We acknowledge support from the BROWN-CNRS University Accord which made possible the visit of J. M. to Brown during which most of the work on this project was done, and we are grateful to Herb Fried for his efforts to secure this Accord. One of us (R. B.) wishes to thank Bill Unruh for hospitality at the University of British Columbia during the time when this work was completed. J. M. thanks the High Energy Group of Brown University for warm hospitality. The research was supported in part by the U.S. Department of Energy under Contract DE-FG02-91ER40688, TASK A.
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# Multistability and nonsmooth bifurcations in the quasiperiodically forced circle map
## 1 Introduction
The Arnold circle map \[Arnold65\]
$$x_{n+1}=x_n+\mathrm{\Omega }+\frac{K}{2\pi }\mathrm{sin}(2\pi x_n)\mathrm{mod}\mathrm{\hspace{0.33em}1}$$
(1)
is one of the paradigms for studying properties of nonlinear dynamical systems, both because it is a very simple map, and because of its great physical relevance (see, e.g. \[BBJ85\]). Using the circle map one can model the structure of phase-lockings (devil’s staircase) of a periodically forced nonlinear oscillator \[JBB83, JBB84\] and the current-voltage characteristics of a driven Josephson junction \[BBJ84\]. The phase-locked regions of the Arnold circle map form the well-known Arnold tongues \[Arnold83, Hall84\]. If $`|K|<1`$ there is a unique periodic attractor with a particular rotation number in each tongue.
In this paper we study the structure of the phase-locked regions of the Arnold circle map driven by a rigid rotation with an irrational frequency. This system exhibits different kinds of dynamics, namely quasiperiodic motions with two and three incommensurate frequencies, chaotic attractors, and strange nonchaotic attractors (SNAs). SNAs have a strange geometrical structure, but unlike chaotic attractors they do not exhibit a sensitive dependence to changes in the initial conditions, i.e. their dynamics is not chaotic. They have been found in many quasiperiodically forced systems \[GOPY84\], and also in the quasiperiodically forced circle map \[DGO89, FKP95\].
Previous investigations show that regions of bistability occur in phase-locked regions of the quasiperiodically forced circle map \[GFPS99\] and the phase-locked regions change in shape depending on the strength of the forcing. This change in shape is related to the emergence of SNAs \[FGO97\]. Based on these studies the aim of this paper is twofold. We study regions in parameter space where more than two attractors coexist (pockets of multistability). Secondly, we discuss the relation between these multistable regions and the appearance of strange nonchaotic attractors. In our discussion of the way these attractors are created and destroyed we are led to a description of nonstandard (nonsmooth) bifurcations of the invariant curves.
If the unforced circle map is modified by introducing additional nonlinearities, coexisting attractors with the same rotation number can occur within the phase-locked regions \[McGP96\]. In our system, multistability within the phase-locked regions is induced by the forcing rather than an additional nonlinear term. In quasiperiodically forced systems the coexisting attractors may be either invariant curves or SNAs depending on the strength of the forcing. Our investigation focuses on how these multistable regions appear and disappear under variation of the system’s parameters: the nonlinearity $`K`$ and the forcing amplitude $`\epsilon `$. For the tongue with zero rotation number the multistable regions open and close by smooth saddle-node or pitchfork bifurcations of invariant curves if $`K`$ and $`\epsilon `$ are small. For larger $`K`$ these saddle-node and pitchfork bifurcations become nonsmooth: instead of merging uniformly (smooth bifurcation), the relevant stable and unstable invariant curves appear to collide only in a dense set of points.
The paper is organized as follows. Section 2 recalls important properties of the Arnold circle map relevant for this study and their changes under the influence of quasiperiodic forcing. In particular we discuss the phase-locked region with zero rotation number. Within this phase-locked region we find pockets of multistability with a rather complex bifurcation structure which is analyzed in Sec. 3. Smooth and nonsmooth saddle-node and pitchfork bifurcations, leading to coexisting attractors, are studied in Sec. 4 to get a better understanding of the changes in the bifurcation structure depending on the strength of nonlinearity and forcing. Furthermore, we investigate the transition between smooth and nonsmooth bifurcations and its implications to the dynamics of the system. In the full parameter space we find bifurcations of codimension two. In Sec. 5 we discuss a special codimension-2 point that involves only nonsmooth bifurcations. It turns out that the unfolding of this point is very different from the smooth analog. Finally, in Sec. 6, we briefly discuss phase-locked regions with small, but finite, rotation number. We conclude this paper with a summary in Sec. 7 and an Appendix with details on the numerical computations. For readers with a black and white copy of this article we provide a supplementary website \[OWGFwww00\].
## 2 The Quasiperiodically Forced Circle Map
The quasiperiodically forced circle map is a map on the torus with lift
$`x_{n+1}`$ $`=`$ $`x_n+\mathrm{\Omega }+{\displaystyle \frac{K}{2\pi }}\mathrm{sin}(2\pi x_n)+\epsilon \mathrm{sin}(2\pi \vartheta _n),`$ (2)
$`\vartheta _{n+1}`$ $`=`$ $`\vartheta _n+\omega \mathrm{mod}\mathrm{\hspace{0.33em}1},`$ (3)
where $`\vartheta _n`$ and $`x_n`$ modulo 1 give the coordinates on the torus. The parameter $`\mathrm{\Omega }`$ is the phase shift, $`K`$ denotes the strength of nonlinearity ($`K>0`$), $`\epsilon `$ is the forcing amplitude, and the forcing frequency $`\omega `$ is irrational. Throughout this paper we choose to work with $`\omega =(\sqrt{5}1)/2`$.
### 2.1 The unforced system
Let us recall the behavior of the unforced circle map (1). The dynamics of this map can be either periodic, quasiperiodic, or chaotic, depending on the parameters $`\mathrm{\Omega }`$ and $`K`$. The critical line $`K=1`$ divides the parameter space into two regions. If $`K<1`$ the map is invertible and the motion can only be periodic (phase-locked) or quasiperiodic. For $`K>1`$ the map is noninvertible and chaotic motion is possible.
The rotation number is used to characterize the different kinds of motion. It is defined as
$$\rho (\mathrm{\Omega },K)=\underset{N\mathrm{}}{lim}\frac{x_Nx_0}{N},$$
(4)
where $`x_N`$ is the $`N`$th iterate of (1), starting from $`x_0`$. It can be shown that $`\rho (\mathrm{\Omega },K)`$ does not depend on $`x_0`$ if $`K<1`$. If the rotation number is rational, the attracting motion is periodic, otherwise it is quasiperiodic. For $`K<1`$ the parameter space is split into regions with rational rotation number, the phase-locked regions or Arnold tongues, and regions with irrational rotation number corresponding to quasiperiodic motion. For example, the main tongue, the phase-locked region with zero rotation number, is bounded by the curves $`\mathrm{\Omega }=\pm \mathrm{\Omega }_0(K)`$, where $`\mathrm{\Omega }_0(K)=K/2\pi `$. For any choice of $`\mathrm{\Omega }`$ and $`K`$, with $`|\mathrm{\Omega }|<\mathrm{\Omega }_0`$ and $`K<1`$, there are exactly two fixed points, one is attracting and the other is repelling. At the boundary $`|\mathrm{\Omega }|=\mathrm{\Omega }_0`$ the two fixed points are annihilated in a saddle-node bifurcation. For other rotation numbers $`\rho 0`$ the $`K`$-dependency of the boundary $`\mathrm{\Omega }_\rho (K)`$ is nonlinear, but it is always a curve of saddle-node bifurcations.
### 2.2 The forced system
A variety of behavior is possible in the coupled maps (2)–(3). The rotation number (4) exists, but depends on $`\epsilon `$ and $`\omega `$ in addition to $`\mathrm{\Omega }`$ and $`K`$, and the direct analogs of the periodic and quasiperiodic motion of the uncoupled Arnold map are invariant curves (the graph of a function $`\vartheta x(\vartheta )`$) and motion which is dense on the torus, respectively. In the former case the rotation number is rationally related to $`\omega `$ ($`\rho =r_1+r_2\omega `$ with $`r_i`$ rational, $`i=1,2`$) and in the latter case there is no such rational relation. For both rationally and irrationally related rotation numbers strange nonchaotic attractors (SNAs) may also be possible. An SNA has a strange geometric structure, that is, it can be viewed as the graph of an everywhere discontinuous function $`\vartheta x(\vartheta )`$, but the dynamics on the attractor is not chaotic, because typical Lyapunov exponents in the $`x`$–direction are negative (there is always a zero Lyapunov exponent in the $`\vartheta `$–direction).
If the rotation number is rationally related to $`\omega `$ then the motion is said to be phase-locked and the regions of parameter space in which the motion is phase-locked are analogous to the Arnold tongues of the unforced map. On the boundaries of the phase-locked regions we expect to see saddle-node bifurcations. There is an additional complication in the forced maps \[FKP95\] in that the saddle-node bifurcations may be smooth (two invariant curves converge uniformly from inside the phase-locked region) or nonsmooth. In the smooth saddle-node bifurcation the nontrivial Lyapunov exponent in the $`x`$direction goes to zero at the bifurcation point. In the nonsmooth saddle-node bifurcation the two invariant curves appear to collide only on a dense set of points. Moreover, the typical nontrivial Lyapunov exponent remains negative. These nonsmooth saddle-node bifurcations seem to be associated with the appearance of SNAs outside the phase-locked region \[DGO89, FKP95, Glendinning98\].
It can be shown that one mechanism of the appearance of SNAs is related to changes in the shape of the phase-locked regions \[FGO97, GFPS99\]. For the unforced map the width of a phase-locked region increases monotonically with increasing nonlinearity $`K`$; this is no longer the case for positive forcing amplitude $`\epsilon `$. Moreover, as Fig. 1 shows, for fixed $`K`$ the width of the phase-locked region oscillates as $`\epsilon `$ increases. In particular, there are certain values of $`\epsilon `$ for which the width of the phase-locked region becomes extremely small. Unfortunately, using only numerical methods, we cannot decide whether the region actually closes or not at those $`\epsilon `$-values. For more details on numerical computations we refer to App. A.1.
For small fixed nonlinearity $`K`$ the boundary $`\mathrm{\Omega }_0(\epsilon ,K)`$ of the phase-locked region can be approximated by (the modulus of) a Bessel function of order zero using first order perturbation theory \[GFPS99\]. Numerical simulations also revealed regions of bistability in the vicinity of the zeroes of the Bessel function, where the width of the phase-locked region is very small. The bistability regions are bounded by saddle-node bifurcations of invariant curves, which has been confirmed by second order perturbation theory \[GW99\].
The study in \[GFPS99, GW99\] only applies for $`K`$ close to $`0`$. We wish to study what happens to the phase-locked region with zero rotation number for larger $`K`$. However, we restrict our considerations to the invertible case $`K<1`$, so that chaos is ruled out. We find that the regions of bistability contain other regions where even more attractors coexist. In the following we describe how these regions appear and disappear as a parameter varies. We also study smooth and nonsmooth bifurcations and make some remarks on the appearance of SNAs.
The majority of the rest of this paper describes the results of numerical simulations of Eqs. (2)–(3). As such, the reader should bear in mind that our conclusions are based on numerical observations and may turn out to be misleading in places. We have made every effort to avoid such problems (see the Appendix) and believe that the phenomena reported are sufficiently interesting and mathematically intractable to merit this numerical investigation, even if we remain uncertain of some of the outcomes. The reader is encouraged to maintain a healthy scepticism throughout.
## 3 The Internal Structure of the Main Tongue for $`K=0.8`$
In the simplest case the boundary $`|\mathrm{\Omega }|=\mathrm{\Omega }_0(\epsilon ,K)`$ represents the disappearance of two invariant curves, a stable and an unstable one. However, for $`|\mathrm{\Omega }|<\mathrm{\Omega }_0(\epsilon ,K)`$ more than two invariant curves may exist that disappear before this boundary is crossed. Such pockets of multistability are found near local minima of $`\mathrm{\Omega }_0(\epsilon ,K)`$, cf. the region of bistability predicted by perturbation analysis for small $`K`$ \[GFPS99, GW99\]. For example, Fig. 2(a) shows a cross-section of Fig. 1 at $`K=0.8`$, with both positive and negative sides of the boundary of the phase-locked region. The outer boundary is the function $`|\mathrm{\Omega }|=\mathrm{\Omega }_0(\epsilon ,0.8)`$. Extra curves are drawn marking the boundaries of pockets of multistability, which is best seen in the enlargements Figs. 2(b), 4 and 5. These pockets of multistability can be considered as overlaps of different “bubbles” with the same rotation number as in Fig. 2(a). For better visualization we have chosen different colors for bifurcations of different pairs of invariant curves.
### 3.1 Bifurcations for $`K=0.8`$ in the first region of overlap
The first overlap is enlarged in Fig. 2(b). The orange and purple curves enclose a rhombus shaped region where two attracting and two repelling invariant circles exist. This region of bistability is bounded by curves of saddle-node bifurcations that end in pitchfork bifurcations on the line $`\mathrm{\Omega }=0`$. Note that if $`\mathrm{\Omega }=0`$ the map has a symmetry ($`xx,\vartheta \vartheta +1/2`$) which implies that the rotation number in the $`x`$–direction is always zero and pitchfork bifurcations should be expected. The bifurcation sequence for $`\mathrm{\Omega }=0`$ is sketched in Fig. 3(a) where each circle is represented as a point and $`\epsilon `$ increases along the horizontal axis. The bottom and top lines are identical, representing the modulo 1 computations.
The purple and orange dots are the pitchfork bifurcations that mark the crossing of the purple and orange curves in Fig. 2(b) along $`\mathrm{\Omega }=0`$, respectively.
### 3.2 Bifurcations for $`K=0.8`$ in the second region of overlap
Figure 4(a) shows a detail of the second overlap. This picture is very similar to Fig. 2(b), but the bifurcation diagram along $`\mathrm{\Omega }=0`$ in Fig. 3(b) reveals a more complex structure. Let us first discuss Fig. 3(b). The first bifurcation (light blue dot) is the same pitchfork bifurcation as the purple dot in Fig. 3(a): the stable circle becomes unstable, creating two new stable circles. As $`\epsilon `$ increases two other stable circles are born in a pair of saddle-node bifurcations (red dots). We now have four different attractors. Note that these saddle-node bifurcations happen at the same values of $`\epsilon `$ due to the symmetry along the line $`\mathrm{\Omega }=0`$, as referred to earlier.
At the purple dots the two attracting circles from the pitchfork bifurcation disappear in a pair of saddle-node bifurcations. Note that this pair of saddle-node bifurcations is connected to the purple pitchfork bifurcation of Fig. 3(a) via the purple curve off $`\mathrm{\Omega }=0`$ in Fig. 2(a). We are now left with two attractors and two repellors. These last two attractors do not disappear in a pitchfork bifurcation as in Fig. 3(a) (orange dot). Instead, they disappear in a pair of saddle-node bifurcations (dark blue dots) with two repellors that are born in a pitchfork bifurcation (green dot) for slightly smaller $`\epsilon `$. Note that this pitchfork bifurcation is subcritical, as opposed to the supercritical orange one in Fig. 3(a).
The unfolding of these bifurcations inside the phase-locked region with $`\mathrm{\Omega }0`$ is shown in detail in Figs. 4 and 5. The curves are colored according to the colors of the bifurcations in Fig. 3(b). As expected, for $`\mathrm{\Omega }0`$ pairs of saddle-node bifurcations no longer happen at the same values of $`\epsilon `$. They form two different curves that cross each other exactly at $`\mathrm{\Omega }=0`$. We already mentioned earlier that the purple curves connect all the way left in a pitchfork bifurcation on the line $`\mathrm{\Omega }=0`$ in the first overlap. The light blue curves start in the pitchfork bifurcation at $`\mathrm{\Omega }=0`$ and become the outer boundary $`|\mathrm{\Omega }|=\mathrm{\Omega }_0(\epsilon ,0.8)`$ once they cross the purple curves. The red curves form swallowtails with the purple curves on the right side and the dark blue curves on the left side; see Fig. 4(b). Finally, the dark blue curves form swallowtails with red and green curves; see Fig. 5.
## 4 The Structure of Bifurcations in the $`(\epsilon ,K)`$-plane
The bifurcation structure depends on the strength of the nonlinearity $`K`$. We study the two-parameter dependence only on the cross-section $`\mathrm{\Omega }=0`$, because the unfolding inside the phase-locked region with $`\mathrm{\Omega }0`$ is similar to that discussed in Sec. 3. As expected, the regions of overlap change shape with $`K`$. In particular, only the regions with no more than two attractors persist for small $`K`$ and moderate $`\epsilon `$. This is shown in Fig. 6 with a cross-section at $`\mathrm{\Omega }=0`$ in the $`(\epsilon ,K)`$-plane of the second overlap; see App. A.2 for details on how this picture was generated.
In the following sections we describe the bifurcations in more detail. Section 4.1 discusses the sequence of bifurcations that happen as $`K`$ decreases. For small $`K`$ (less than approximately $`0.8`$ for the bifurcations we have looked at), saddle-node and pitchfork bifurcations happen via a uniform collision of invariant curves: at the moment of bifurcation, two (saddle-node) or three (pitchfork) curves merge at each value of $`\vartheta `$. We call these bifurcations smooth bifurcations. For $`K`$ close to 1 the attractors may become extremely wrinkled, which gives rise to nonsmooth bifurcations: at the moment of bifurcation the invariant curves now collide only in a dense set of $`\vartheta `$-values. The nonsmooth pitchfork and nonsmooth saddle-node bifurcations are described in detail in Secs. 4.2 and 4.3, respectively.
### 4.1 Smooth bifurcations for $`\mathrm{\Omega }=0`$ in the second region of overlap
Figure 6 indicates that at most two attractors exist for small $`K`$ and moderate $`\epsilon `$. This means that for fixed small $`K`$, the bifurcation portrait looks like Fig. 3(a). Hence, as we decrease $`K`$ from $`K=0.8`$ to 0, the extra pairs of saddle-node bifurcations (see Fig. 3(b)) need to disappear somehow. It turns out that the last pitchfork bifurcation, the green dot in Fig. 3(b), “absorbs” these saddle-node bifurcations one by one. In doing so, the pitchfork bifurcation switches from subcritical to supercritical and vice versa (a standard codimension-2 bifurcation). A sketch of this process along the line $`\mathrm{\Omega }=0`$ is shown in Fig. 7(a)–(c).
In the $`(\mathrm{\Omega },\epsilon )`$-plane the picture changes as follows. The first swap, Fig. 7(a), comes about as the pitchfork point and the intersection point of the dark blue curves on $`\mathrm{\Omega }=0`$ collapse; see Fig. 5(b). When $`K`$ decreases, these points move closer together, causing the slope of the dark blue curves to become steeper and the ends of the swallowtail to move closer to $`\mathrm{\Omega }=0`$. Upon collision the green curves disappear and the dark blue curves end in a supercritical pitchfork bifurcation.
In the second swap the dark blue curves disappear in a similar way via a collision of the pitchfork point and the intersection point of the two red curves, making the pitchfork subcritical again. Note that, in order for this to happen, the pitchfork point crosses the intersection point of the two purple curves; compare Figs. 4(b) and 5(a). Figures 7(a)–(b) show why this is a crossing and not a collision: the pair of purple saddle-node bifurcations happens “far out” in state space from the pitchfork bifurcation. Therefore, the crossing is only a crossing in this projection on the $`(\mathrm{\Omega },\epsilon )`$-plane.
The third swap, Fig. 7(c), is identical to the first, causing the disappearance of the red curves. In this bifurcation diagram at most two attractors coexist, which is the desired situation for $`K`$ small.
### 4.2 Nonsmooth pitchfork bifurcations
For $`K`$ close to 1 the situation is more complicated, because some of the invariant curves are very wrinkled and the pitchfork bifurcation becomes nonsmooth. This bifurcation has been found by Sturman in a similar map. Let us discuss what happens for $`\mathrm{\Omega }=0`$, along the line $`K=0.9`$ as we approach the pitchfork bifurcation by decreasing $`\epsilon `$, starting in the yellow region in Fig. 6. Figure 8(a) shows all invariant curves just before the bifurcation. The two stable invariant curves (black and blue) correspond to the two outer branches of the pitchfork, the unstable invariant curve (red) relates to the inner branch of the pitchfork separating the two outer ones. The fourth invariant curve (green) is also unstable, but it is “far away” and does not take part in the bifurcation. As we decrease $`\epsilon `$ towards the bifurcation point the three invariant curves (blue, black and red) approach each other, but due to their wrinkled structures they appear to collide only in a dense set of $`\vartheta `$-values instead of merging uniformly as in a smooth bifurcation. This indicates that at the moment of bifurcation the attractor is an SNA. Numerical evidence suggests that this SNA persists and smoothes out to an invariant curve over a small $`\epsilon `$-interval; see for example the red attractor in Fig. 8(b) for $`\epsilon `$ below the bifurcation value. It is possible that we see the reverse of fractalization, a mechanism for the appearance of SNAs reported in \[NK96\]. However, if we use the method of rational approximations for testing whether the attractor is an SNA, we get conflicting results; see App. A.3 for more details. We remark that we get these conflicting numerical results only for the nonsmooth pitchfork bifurcation. In any case, after a further decrease in $`\epsilon `$ the attractor is clearly a smooth invariant curve.
It is important to note that the nonsmooth pitchfork bifurcation is uniquely defined as the moment of collision of three invariant curves and the locus of bifurcation lies on a curve in the ($`\epsilon ,K`$)-plane. The process of fractalization is a gradual process where the moment of transition from an invariant curve to an SNA is not well-defined numerically. We wish to emphasize that it is, therefore, completely unclear whether the set of parameter pairs $`(\epsilon ,K)`$ with $`\mathrm{\Omega }=0`$ that exhibit SNAs after the nonsmooth pitchfork bifurcation has zero or finite size.
Since the boundary of the region of bi- or multistability for $`\mathrm{\Omega }=0`$ is given by smooth and nonsmooth pitchfork bifurcations there is a codimension-2 point in the ($`\epsilon ,K`$)-plane where the smooth and the nonsmooth pitchfork bifurcation curves meet. An approximation of this codimension-2 point is ($`\epsilon ,K`$) = (1.564, 0.89). In any neighborhood of this point we always find all three kinds of dynamical behaviors: one stable invariant curve, two stable invariant curves, and one SNA; although the latter may only exist on the nonsmooth bifurcation curve itself.
### 4.3 Nonsmooth saddle-node bifurcations
If $`K`$ is small and $`\epsilon `$ moderate then the saddle-node bifurcations observed numerically involve two invariant curves on the torus which converge and destroy each other. At larger values of $`K`$, simulations suggest that two invariant curves touch on an orbit at the bifurcation point, so points of intersection of these sets are dense on the curves. (Strictly speaking, the invariant sets are no longer continuous at the bifurcation point, but we will continue to refer to them as curves.) For quasiperiodically forced circle maps we can distinguish two types of these nonsmooth saddle-node bifurcations: one-sided and two-sided. In the one-sided nonsmooth saddle-node bifurcation, the collisions occur between pairs of invariant curves on the cylinder. An example is shown in Fig. 9(a) where the invariant curves and some of their translates by one in the $`x`$direction are computed close to the bifurcation point. In the two-sided case each stable invariant curve on the cylinder touches both the unstable invariant curve immediately above it and the unstable invariant curve immediately below it. On the torus this implies that at the bifurcation point the attractor is everywhere discontinuous. An example of a two-sided nonsmooth saddle-node bifurcation is shown in Fig. 10(a). These two-sided nonsmooth saddle-node bifurcations are described in \[FKP95\], where it is shown that after the bifurcation (with $`\mathrm{\Omega }=0`$ and $`K`$ fixed) the map has an SNA with unbounded motion in the $`x`$direction (Fig. 10(b)) despite the fact that the rotation number remains zero due to the symmetry of Eqs. (2)–(3) when $`\mathrm{\Omega }=0`$. This implies that the diffusion in the $`x`$direction is extremely slow; see \[FKP95, SFGP99\] for further details. In general, the two-sided nonsmooth saddle-node bifurcation is of codimension two, but it occurs as a codimension-1 phenomenon due to the symmetry if $`\mathrm{\Omega }=0`$. Figure 11 shows the range of dynamics observed in the third region of overlap in the plane $`\mathrm{\Omega }=0`$. Unbounded SNAs are observed in the blue regions and two-sided nonsmooth saddle-node bifurcations occur on the boundary between the white and blue regions.
The unbounded SNA of Fig. 10(b) must contain orbits which are unbounded above and orbits which are unbounded below \[SFGP99\]. This bidirectional diffusive motion of the unbounded SNA with $`\mathrm{\Omega }=0`$ becomes effectively unidirectional if $`|\mathrm{\Omega }|`$ is very small, leading to a nonzero rotation number \[SFGP99\]. This suggests that these unbounded SNAs lie on the boundary of the phase-locked region. Fig. 12 shows this boundary in $`(\epsilon ,K,\mathrm{\Omega })`$\- space. It is clear that the height (i.e. the width in $`\mathrm{\Omega }`$) of the boundary is very small, if not zero, in regions of the $`(\epsilon ,K)`$ plane with $`\mathrm{\Omega }=0`$ which have unbounded SNAs (compare the low plateau on the left of Fig. 12 with the region of unbounded SNAs of Fig. 11). As the height of the boundary becomes non-negligible we observe that the saddle-node bifurcation on the boundary has become one-sided (see Fig. 9), and we believe that it remains one-sided and nonsmooth throughout the red areas of Fig. 12 with non-negligible height. These red regions of the boundary correspond to saddle-node bifurcations with negative nontrivial Lyapunov exponents, and appear to preceed the creation of SNAs with nonzero rotation numbers outside the phase-locked region \[DGO89, Glendinning98\].
If the phase-locked region really has zero height on the plateau, then it is wrong to refer to saddle-node bifurcations on the interior of the plateau: these points would correspond to a transition from an SNA with negative rotation number to an SNA with positive rotation number through an SNA with zero rotation number as $`\mathrm{\Omega }`$ increases through zero. In this full three-parameter unfolding the two-sided nonsmooth saddle-node bifurcations are of codimension two, occurring on curves bounding the plateau and separating parts of the boundary of the phase-locked region with unbounded SNAs from parts with one-sided nonsmooth saddle-node bifurcations.
## 5 Nonsmooth bifurcation points of codimension two
In the full parameter space, and also already in the $`(\epsilon ,K)`$-plane with $`\mathrm{\Omega }=0`$, we expect codimension-2 bifurcations. We have already seen some of these; for example, the transition from yellow to red in Figs. 1 and 12 is the codimension-2 bifurcation curve marking the transition from a smooth saddle-node bifurcation to a nonsmooth one, respectively \[KNPS00\]. Another example is the codimension-2 point, mentioned in Sec. 4.2, where the smooth and nonsmooth pitchfork bifurcation curves meet. Note that there is no curve of codimension-2 points in this case, since the pitchfork bifurcation is restricted to the plane $`\mathrm{\Omega }=0`$. As discussed in the previous section, there is also the two-sided nonsmooth saddle-node bifurcation curve.
In this section we want to draw attention to an interesting codimension-2 point in the plane $`\mathrm{\Omega }=0`$. This point can be seen in Fig. 11 and in the enlargements Figs. 13(a)–(b) as the point where the region with unbounded SNAs (blue) and the bistable region (yellow) touch. It can be characterized as the moment where two nonsmooth (supercritical) pitchfork bifurcations happen simultaneously, i.e. a region of overlap is pulled apart; compare also the sketch in Fig. 7 of \[GFPS99\].
### 5.1 Smooth analog of the nonsmooth codimension-2 bifurcation point
Let us first discuss the smooth analog of this nonsmooth codimension-2 bifurcation point. Suppose for $`\mathrm{\Omega }=0`$ and some $`K<1`$ fixed the bifurcation diagram involves only two supercritical pitchfork bifurcations that occur in the order as shown in Fig. 14(a). Now assume that as we increase $`K`$, this order is switched before we reach $`K=1`$, without changing the type of pitchfork bifurcation from supercritical to subcritical. As shown in Fig. 14(b), this means that we necessarily need extra curves of saddle-node bifurcations.
The smooth analog of the nonsmooth codimension-2 bifurcation point is the point where the two pitchfork bifurcations happen at the same parameter values. In the $`(\epsilon ,K)`$-plane the complete bifurcation diagram should look like Fig. 15. The two supercritical pitchfork bifurcation curves are colored light-blue and magenta. The green and light-green curves are saddle-node bifurcation curves; compare also the colors in Figs. 14(a)–(b). Above the green and light-green curves, but below the light-blue and magenta curves there are four attractors. If we cross either the light-blue or the magenta curve from this region there are three attractors. Above the light-blue and magenta curves there are two attractors.
We distinguish five qualitatively different $`K`$-intervals numbered I–V in Fig. 15. Interval I corresponds to Fig. 14(a) and V to Fig. 14(b). The qualitative behavior in the intervals II–IV is given in Figs. 16(a)–(c), respectively. The codimension-2 point that we are discussing here is the intersection point of the light-blue and magenta pitchfork bifurcation curves. The intersections of the green and light-blue curves, and the light-green and magenta curves are only intersections in this projection onto the $`(\epsilon ,K)`$-plane as can be seen in Fig. 16, where the transition from Fig. 16(a) to (b) marks the crossing of green and light-blue, and the transition from Fig. 16(b) to (c) represents the crossing of light-green and magenta.
### 5.2 A nonsmooth bifurcation point of codimension two
The nonsmooth version of Fig. 15 looks surprisingly simple in contrast; see Fig. 17. The coloring of the bifurcation curves is as in Fig. 15 with the restriction that all curves represent nonsmooth bifurcations. The upper pair of bifurcation curves corresponds to two-sided nonsmooth saddle-node bifurcations and the lower pair to nonsmooth pitchfork bifurcations.
We now have four regions with distinctive dynamics. To the left and to the right of the two pitchfork bifurcation curves there is one attractor. In between the two curves of pitchfork bifurcations there are two attractors. This is illustrated for a particular choice of the parameters in Fig. 18(b). The two attractors (black and blue) are two invariant curves separated by unstable invariant curves (green and red). In contrast to the previously discussed case of the nonsmooth pitchfork bifurcation, both unstable invariant curves are now close to the attracting invariant curves, because we are near the codimension-2 bifurcation point. The two different pitchfork bifurcations are illustrated in Figs. 18(a) and (c). If we decrease the forcing amplitude $`\epsilon `$, the black, blue and green invariant curves disappear in a nonsmooth pitchfork bifurcation to form the green attractor coexisting with the red unstable invariant curve, as shown in Fig. 18(a). This corresponds to a crossing of the light-blue nonsmooth pitchfork bifurcation curve in Fig. 17. We cross the magenta nonsmooth pitchfork bifurcation curve in this figure by increasing $`\epsilon `$. In this case the red, black and blue invariant curves disappear and form the red attractor coexisting with the green unstable invariant curve as in Fig. 18(c). The fourth region between the nonsmooth saddle-node bifurcations is characterized by the existence of unbounded SNAs.
Finally, we remark that the SNA region seems to have a fractal-like structure in the neighborhood of the codimension-2 point; compare Figs. 13(a)–(b). This would imply that the boundary of the phase-locked region is fractal!
## 6 Tongues with Nonzero Rotation Numbers
We studied the structure of the phase-locked region with zero rotation number in great detail. In this section we briefly discuss the geometry of the tongues with nonzero rotation numbers that are close to the main tongue. More precisely, we determine the boundaries of the tongues with rotation numbers $`\rho =1/F_k`$, with $`F_k`$ the $`k`$th Fibonacci number $`F_k=F_{k1}+F_{k2}`$ and $`F_1=F_2=1`$, using the same numerical procedure as before; see App. A.1. Several tongues for $`K=0.99`$ are shown in Fig. 19. The fluctuations of the widths of the tongues with $`\rho >0`$ are due to numerical errors, which are of the same magnitude as the widths themselves. The thick borderline $`\mathrm{\Omega }_0(\epsilon ,K)`$ in this parameter regime corresponds to transitions to SNAs or bifurcations of SNAs of the type discussed in Sec. 4.3.
How these tongues approach the boundary $`\mathrm{\Omega }_0(\epsilon ,K)`$ for fixed $`(\epsilon ,K)`$ as $`k`$ goes to infinity depends in general on $`\epsilon `$ and $`K`$; compare $`\epsilon =2.55`$ and $`\epsilon =2.61`$ in Fig. 19. This behavior can be quantified with a scaling law for the distance $`\mathrm{\Omega }_\rho \mathrm{\Omega }_0`$ between the main tongue and the tongues with nonzero rotation numbers $`\rho `$. We used several parameter pairs $`(\epsilon ,K)`$ in Fig. 20 and our numerical calculations strongly suggest the scaling law $`\rho \mathrm{\Omega }_\rho `$ for $`\epsilon `$ and $`K`$ such that $`\mathrm{\Omega }_0=0`$, and $`\rho \sqrt{\mathrm{\Omega }_\rho \mathrm{\Omega }_0}`$ otherwise. This agrees with the linear scaling in the high-$`\epsilon `$ limit conjectured by Ding et al. , the square-root scaling in the Arnold circle map for $`\epsilon =0`$ and $`K>0`$ (see e.g. \[MacKayTresser84\]) and the trivial linear scaling of the pure rotation for $`\epsilon =K=0`$. The square root scaling is associated with the one-sided saddle-node bifurcation, and the linear scaling is related to the change in the rotation number of the SNAs as in Sec. 4.3. The change between these scalings is associated with a two-sided nonsmooth saddle-node bifurcation point of codimension two at which two one-sided nonsmooth saddle-node bifurcation curves meet (presumably the nonsmooth analog of a cusp bifurcation point). A detailed study of this curve of codimension-2 points could be of interest.
## 7 Summary
We have studied the structure of the phase-locked regions in the quasiperiodically forced circle map. In particular, we have found regions of multistability where several attractors coexist. These regions of multistability appear due to the emergence of additional pairs of invariant curves as a result of saddle-node or pitchfork bifurcations under the variation of the forcing amplitude. As a result, these regions look like overlaps of phase-locked regions with the same rotation number.
Opening and closing of these pockets of multistability are due to saddle-node and pitchfork bifurcations of invariant curves. These bifurcations can be either smooth or nonsmooth depending on the strength of nonlinearity and the forcing amplitude. This is organized by the type of interaction between the stable and unstable invariant curves. In the smooth case these curves approach each other uniformly and then touch uniformly in each value of $`\vartheta `$, so that the bifurcation looks like a simple merging of the invariant curves, analogous to the unforced case. Nonsmooth saddle-node bifurcations appear due to a wrinkled structure of the participating invariant curves, which collide at the bifurcation only in a dense set of $`\vartheta `$-values. The result of this bifurcation is the emergence of a strange nonchaotic attractor. Similar to the nonsmooth saddle-node bifurcation we find a nonsmooth pitchfork bifurcation, but the details of this bifurcation are still unclear. Both for the saddle-node and the pitchfork bifurcation there are codimension-2 points in parameter space marking the transition from a smooth to a nonsmooth bifurcation. The exact determination of the codimension-2 pitchfork bifurcation point and the self–similarity properties in its neighborhood should be possible by using renormalization group techniques.
We have also investigated a nonsmooth codimension-2 bifurcation involving the merging of two nonsmooth pitchfork bifurcations and two nonsmooth saddle-node bifurcations which leads to regions of unbounded SNAs. This has no straightforward analog in smooth bifurcations.
The positions of phase-locked regions in the neighborhood of the region with zero rotation number were also described. Fixing the nonlinearity $`K`$ and the forcing $`\epsilon `$, we found that the rotation number $`\rho `$ of these regions scales linearly with $`\mathrm{\Omega }_\rho `$ whenever the width of the main tongue appears to be zero, and as $`\sqrt{\mathrm{\Omega }_\rho \mathrm{\Omega }_0}`$ otherwise.
Acknowledgements
We thank for the hospitality of the Max Planck Institute for Physics of Complex Systems in Dresden, where this collaboration started during the Workshop and Seminar “Beyond Quasiperiodicity: Structures and Complex Dynamics” (January 1999), and where the work was completed in March 2000. H.O. was supported by an AFOSR/DDRE MURI grant AFS-5X-F496209610471 while she was employed at Caltech in Pasadena, USA. J.W. acknowledges financial support by the EU through a TMR Network on the dynamics of spatially extended systems under contract number FMRXCT960010. P.G. is grateful for a travel grant from the British Council. U.F. acknowledges financial support by the Deutsche Forschungsgemeinschaft (Heisenberg-Program and Sfb 555).
## Appendix: Details on the Numerical Computations
We now explain our numerical computations in more detail and give values for the accuracy parameters. The Appendix is organized such that each section is related to one section in the main text. Appendix A.1 explains how to compute the tongue boundary for the main tongue with zero rotation number and relates to Sec. 2.2. Appendix A.2 discusses the generation of Fig. 6 in Sec. 4. Here, we also discuss how to identify SNAs using the phase sensitivity exponent, which is related to the derivative of $`x_n`$ with respect to the external phase $`\vartheta `$. The numerical issues that arise when determining the nonsmooth pitchfork bifurcation and the accompanying fractalization process, as reported in Sec. 4.2, are described in more detail in App. A.3. In contrast to the method for identifying SNAs in the second region of overlap as described in App. A.2, we used another more efficient method for determining the regions where SNAs exist in the third region of overlap. This method is explained in App. A.4.
### A.1 Numerical computation of the boundary of the phase-locked region
The boundary $`\mathrm{\Omega }=\mathrm{\Omega }_0(\epsilon ,K)`$ of the phase-locked region with $`\rho =0`$ is half its width due to the symmetry. It is approximated by estimating the boundary point $`\mathrm{\Omega }_0`$ on a grid of $`320\times 40`$ points in the $`(\epsilon ,K)`$-plane. Following \[SFGP99\] we determine the rotation number (4) within an accuracy of $`\pm 1/N`$ by averaging over a sample of 25 orbits of length $`N=F_{28}=\mathrm{317\hspace{0.17em}811}`$ (after 1000 preiterations to eliminate the effect of transients), where $`F_k`$ are the Fibonacci numbers $`F_1=F_2=1`$ and $`F_k=F_{k1}+F_{k2}`$. The Fibonacci numbers $`F_k`$ are used since the value $`\vartheta _{F_k}`$ after $`F_k`$ iterations is close to the initial value $`\vartheta _0`$ due to the fact that ratios of Fibonacci numbers are good rational approximants of our irrational driving frequency $`\omega `$. The initial interval $`[\mathrm{\Omega }_0^{},\mathrm{\Omega }_0^+]=[0,0.2]`$ is repeatedly bisected, preserving the relation $`\rho (\mathrm{\Omega }_0^{})<1/N<\rho (\mathrm{\Omega }_0^+)`$ to ensure that $`\mathrm{\Omega }_0[\mathrm{\Omega }_0^{},\mathrm{\Omega }_0^+]`$, until $`\mathrm{\Omega }_0^+\mathrm{\Omega }_0^{}<\mathrm{\Delta }\mathrm{\Omega }=10^5`$. Finally, we choose $`\mathrm{\Omega }_0`$ to be the mean value of $`\mathrm{\Omega }_0^{}`$ and $`\mathrm{\Omega }_0^+`$, or zero if the mean value is smaller than our numerical accuracy $`\mathrm{\Delta }\mathrm{\Omega }`$. Note that, as remarked in the Sec. 2.2, the question of whether $`\mathrm{\Omega }_0`$ really vanishes or is just very small cannot be answered by using only numerical methods.
To distinguish between smooth and nonsmooth saddle-node bifurcations of invariant curves on the boundary of the phase-locked region we compute the nontrivial Lyapunov exponent
$$\lambda (\mathrm{\Omega },\epsilon ,K)=\underset{N\mathrm{}}{lim}\frac{1}{N}\underset{n=0}{\overset{N1}{}}\mathrm{ln}\left|\frac{x_{n+1}}{x_n}\right|_{(x_n,\vartheta _n)}=\underset{N\mathrm{}}{lim}\frac{1}{N}\underset{n=0}{\overset{N1}{}}\mathrm{ln}|1+K\mathrm{cos}2\pi x_n|.$$
(5)
Vanishing $`\lambda `$ (yellow regions in Fig. 1) indicates smooth saddle-node bifurcation while negative $`\lambda `$ indicates nonsmooth saddle-node bifurcation (red regions); see Feudel et al. and Sec. 2.2.
### A.2 The cross-section $`\mathrm{\Omega }=0`$ of the second region of overlap
Figure 6 shows the bifurcation structure at a cross-section $`\mathrm{\Omega }=0`$. The picture was generated as follows. For each grid point $`N=F_{32}=\mathrm{2\hspace{0.17em}178\hspace{0.17em}309}`$ iterations of Eqs. (2)–(3) were computed using 25 different initial conditions $`(x_0,\vartheta _0=0)`$. We take advantage of the fact that inside the tongue with zero rotation number, each attractor is represented as a single-valued function $`x=X(\vartheta )`$, $`\vartheta [0,1)`$. (This function is smooth in the case of a nonstrange attractor and discontinuous everywhere in the case of an SNA.) This means that the number of different attractors is equal to the number of different $`N`$th iterates, i.e., different $`x_N`$-values. By setting a tolerance of $`\pm 10^6`$ the number of attractors was determined numerically.
To identify the emergence of SNAs, we compute the attractors and quantify their smoothness properties using the so-called phase sensitivity exponent introduced by Pikovsky and Feudel . By formally differentiating Eq. (2) with respect to the external phase $`\vartheta `$, we get
$$\frac{x_{n+1}}{\vartheta }=(1+K\mathrm{cos}2\pi x_n)\frac{x_n}{\vartheta }+2\pi \epsilon \mathrm{cos}2\pi \vartheta _n.$$
(6)
When Eq. (6) is iterated together with Eqs. (2)–(3), starting from some initial point $`(x_0,\vartheta _0)`$ and $`x_0/\vartheta =0`$, the phase sensitivity
$$\mathrm{\Gamma }_N=\underset{(x_0,\vartheta _0)}{\mathrm{min}}\underset{0nN}{\mathrm{max}}\left|\frac{x_n}{\vartheta }\right|$$
(7)
diverges like $`N^\mu `$ for large $`N`$ in the case of an SNA. On the other hand, in the case of a smooth attractor it saturates, i.e. the phase sensitivity exponent $`\mu `$ is zero. The criteria for saturation we employ are $`\mathrm{\Gamma }_N<10^{15}`$ and $`\mu <0.25`$ (obtained by fitting the slope in a ln-ln diagram using three different $`N`$s). The black area in Fig. 6 shows the parameter region for which $`\mathrm{\Gamma }_N`$ does not saturate when eight different initial points are iterated for up to $`N=F_{33}=\mathrm{3\hspace{0.17em}524\hspace{0.17em}578}`$ time steps. How much of this black area persists as $`N`$ tends to infinity? Figure 21 provides more numerical results for larger $`N`$ and fixed $`K=0.9`$, showing the minimum number of iterations $`N=F_k`$ (in terms of the index of the Fibonacci number) for which the phase sensitivity saturates versus the forcing amplitude $`\epsilon `$. Although a large number of iterations, $`F_{43}=\mathrm{433\hspace{0.17em}494\hspace{0.17em}437}`$, was used, saturation occurs only for $`\epsilon 1.566`$ and $`\epsilon 1.568`$. Hence the interval $`\epsilon [1.566,1.568]`$ contains the nonsmooth pitchfork bifurcation point, as described in Sec. 4.2, which is approximately $`1.5676`$. In the next section we discuss the numerical determination of the length of this gap.
### A.3 Numerical issues regarding the nonsmooth pitchfork bifurcation
Figure 21 shows the minimum number of iterations $`N=F_k`$ (in terms of the index of the Fibonacci number) for which the phase sensitivity exponent converges, versus $`\epsilon `$. There is no convergence of the phase sensitivity exponent, even for $`N`$ as large as $`F_{43}=\mathrm{433\hspace{0.17em}494\hspace{0.17em}437}`$, in the interval $`\epsilon [1.566,1.568]`$. Close to the right side of this interval, a pitchfork bifurcation occurs, because there is one smooth attractor (white region in Fig. 21) for smaller $`\epsilon `$ and there are two smooth attractors (yellow region) for larger $`\epsilon `$. Hence, we clearly have a gap when we bound the number of iterations by $`F_k=F_{43}`$, but it is unclear whether this gap has nonzero width as $`k\mathrm{}`$.
There are other methods to assess the smoothness properties of the attractor(s). For example, in Fig. 22 we applied the method of rational approximations. This method is based on the approximation of the irrational frequency $`\omega `$ by rational frequencies $`\omega _k=F_{k1}/F_k`$ with $`k`$ and $`\omega =lim_k\mathrm{}\omega _k`$, replacing the quasiperiodically forced map (2)–(3) with a sequence of periodically (with period $`F_k`$) forced maps. The $`F_k`$th iteration of such a map is an orientation-preserving diffeomorphism on a circle depending on $`(\mathrm{\Omega },\epsilon ,K)`$ and on the initial phase $`\vartheta _0`$. The union of all attracting invariant sets of this family of diffeomorphisms with $`\vartheta _0[0,1/F_k)`$, forms the $`k`$th approximation of the attractors of the quasiperiodically forced system. (It is sufficient to consider the subinterval $`[0,1/F_k)`$, since diffeomorphisms with $`\vartheta _0[n/F_k,(n+1)/F_k)`$, $`n=1,2,\mathrm{},F_k1`$, are topologically conjugate.) For smooth attractors there is a number $`k`$ for which the rational approximation of order $`k`$ and larger does not depend sensitively on $`\vartheta _0`$ \[PikovskyFeudel94\].
Figure 22 shows that this is here the case: the dependence on the initial phase for moderate $`k`$ — one or two stable fixed points for all $`\epsilon `$, depending on $`\vartheta _0`$ — disappears as we cross the lower curve towards higher $`k`$ values — one stable fixed point in the white marked interval, corresponding to a single smooth attractor in the quasiperiodically forced map, and two stable fixed points in the yellow marked interval, corresponding to a pair of smooth attractors.
However, as mentioned in \[PikovskyFeudel94\], a vanishing dependence on $`\vartheta _0`$ is a necessary but not sufficient condition for smoothness. We also have to stipulate that the maximum derivative of the attracting sets with respect to $`\vartheta _0`$ is bounded for all $`\vartheta _0[0,1)`$ as $`k\mathrm{}`$. Note that it is nevertheless sufficient to determine the attracting sets only in the subinterval $`[0,1/F_k)`$ since the other parts can be obtained by iterating the map (2)–(3) with the rational frequency $`\omega _k`$. Furthermore, it is elegant to iterate Eq. (6) simultaneously in order to determine the derivative. The maximum derivative obtained by this procedure is an approximation of the phase sensitivity. Using the same criteria for boundedness as before, we surprisingly find that the maximum derivative always saturates at some order $`k`$, shown by the upper curve in Fig. 22. The qualitative features of Figs. 2122 persist if we choose different $`K`$s in the region with large phase sensitivity.
### A.4 SNAs in the third region of overlap
To compute SNAs near the boundary of the tongue in the third region of overlap we use a more effective method than the method described in App. A.2. This method takes advantage of the fact that the corresponding SNAs are unbounded in the $`x`$-direction in the lift of the map. We define the amplitude of an attractor at time $`N`$ as
$$\stackrel{~}{\mathrm{\Gamma }}_N=\underset{(x_0,\vartheta _0)}{\mathrm{min}}\left(\underset{0nN}{\mathrm{max}}x_n\underset{0nN}{\mathrm{min}}x_n\right).$$
For an unbounded SNA the asymptotic behavior of the amplitude is given by $`c\mathrm{ln}N`$ \[FKP95\], in all other cases the amplitude saturates for large $`N`$. The numerical algorithm is almost the same as for the phase sensitivity (7). However, by fitting the slope using four different $`N`$, we find as criteria for saturation $`\stackrel{~}{\mathrm{\Gamma }}_N<6.0`$ and $`c<0.02`$. Therefore, typically far less iterations are necessary, so that a larger maximum number of iterations $`N=F_{40}=\mathrm{102\hspace{0.17em}334\hspace{0.17em}155}`$ and 10 different initial conditions can be used.
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# Asymmetric quantum channel for quantum teleportation
## Abstract
We propose a realistic optimal strategy for continuous-variable teleportation in a realistic situation. We show that any imperfect quantum operation can be understood by a combination of an asymmetrically-decohered quantum channel and perfect apparatuses for other operations. For the asymmetrically-decohered quantum channel, we find some counter-intuitive results; teleportation does not necessarily get better as the channel is initially squeezed more. We show that decoherence-assisted measurement and transformation may enhance the fidelity for the asymmetrically mixed quantum channel.
Quantum teleportation is one of the important manifestations of quantum mechanics. In particular, quantum teleportation of continuous variable states has attracted a great deal of attention because of a high detection efficiency, handy manipulation of continuous variable states , and possibility of application to high-quality quantum communication. Two kinds of protocols have been suggested for continuous variable teleportation; one utilizes the entanglement between quadrature-phase variables and the other between the photon-number sum and the relative phase . Both the protocols employ a squeezed two-mode vacuum for the quantum channel. In this paper, we report how to optimize the quantum teleportation of quadrature-phase variables when the quantum channel and experimental conditions are not perfect.
There are a few obstacles which make the teleportation of quadrature-phase variables imperfect. The perfect quantum teleportation is possible only by a maximally-entangled quantum channel, i.e., by an infinitely squeezed state which is unphysical as it incurs the infinite energy. Moreover, when the quantum channel is exposed to the real world, it is influenced by the environment, which turns the pure squeezed state into a mixture and deteriorates the entanglement property. To maximize the channel entanglement, purification protocols for continuous variable states have been suggested by Parker et al. for partially-entangled pure states and by Duan et al. for mixed Gaussian states. However, the theoretical suggestions have not been realized by experiment. Further, there are other obstacles in experiment such as imperfect detection efficiency at the sending station and imperfect unitary transformation at the receiving station. We show that the imperfect conditions may be absorbed into the imperfect quantum channel while other apparatuses are treated perfect, and find the optimization condition for the teleportation under a given experimental condition. We show that blindly maximizing the initial entanglement of the quantum channel does not necessarily bring about the best teleportation.
Two modes $`a`$ and $`b`$ of the squeezed vacuum are distributed, respectively, to a sending and a receiving station. At the sending station, the original unknown state is entangled with the field mode $`a`$ of the quantum channel by a 50/50 beam splitter. Two conjugate quadrature variables are measured respectively for the two output fields of the beam splitter using homodyne detectors. Upon receiving the measurement results through the classical channel, the other mode $`b`$ of the squeezed vacuum is displaced accordingly at the receiving station. The quantum teleportation of quadrature-phase variables is well described by a phase-space distribution, in particular, the Wigner function that is the Fourier transform of its characteristic function $`C(\eta )\text{Tr}[\rho \widehat{D}(\eta )]`$ for the state of the density operator $`\rho `$. $`\widehat{D}(\eta )`$ is the displacement operator.
The quantum teleportation is completed by a unitary displacement operation at the receiving station. If a field state of its Wigner function $`W(\alpha )`$ is displaced by $`\beta `$, it is represented by the Wigner function $`W(\alpha \beta )`$. In the experiment, the displacement operation is performed using a beam splitter of a high transmittance $`T`$ . To displace a field state of the Wigner function $`W(\alpha )`$, it is injected into the beam splitter while a high intensity coherent state of amplitude $`\beta /\sqrt{1T}`$ is injected to the other input port. The beam splitter operation results in the convolution between the two input states . Remembering that the synthesizing coherent state is the displaced vacuum, we find the Wigner function $`W_d(\gamma \beta )`$ of the output field by
$$W_d(\gamma \beta )=\frac{1}{1T}d^2\alpha W(\alpha )W_{vac}\left(\frac{\gamma \beta \sqrt{T}\alpha }{\sqrt{1T}}\right)$$
(1)
where $`W_{vac}(\alpha )`$ is the Wigner function for the vacuum. We can easily see that displacing a field by a beam splitter of its transmittance $`T`$ is equivalent to unitarily-displacing the field after it is mixed with the vacuum at a beam splitter of the same $`T`$. Note that mixing a field with the vacuum at a beam splitter results in the same dynamics of the field influenced by the vacuum environment. Kim and Imoto found that assuming the coupling of the system with the environment $`\kappa `$ and the exposure time to the environment $`\tau `$, the normalized interaction time $`R1\mathrm{exp}(\kappa \tau )`$ is the same as $`1T`$.
The inefficient detection at the sending station is another factor which degrades the teleportation . When the two photomultipliers of a homodyne detector have the same efficiency $`\eta `$, the imperfect homodyne detector is described by a perfect homodyne detector with a beam splitter in front . A field passes through the beam splitter of the transmittance $`\eta `$ and it is mixed with the vacuum which has been injected into the other input port. The inefficiency of the detection can also be passed to the quantum channel. We will discuss later that inefficiency at the sending station gives an effect not only to the quantum channel but also to the original unknown field to teleport.
We have seen that imperfect operations at the receiving and sending stations can be understood as a combination of the perfect operations with an imperfect mixed quantum channel. Imperfection at the displacement operation is absorbed by the field mode to the receiving station. The field mode to the sending station can absorb inefficiency in the homodyne detection. These considerations lead the quantum channel mixed asymmetrically due to a different condition for each mode of the quantum channel. The study of the teleportation using the asymmetrically-decohered quantum channel is important not only because it can explain the experimental situation but also because it gives novel features and deeper understanding of the nature of entanglement for the continuous variables. If quantum teleportation is used for a quantum communication, it is more likely that the two modes of the quantum channel will undergo different environmental conditions. To the best of our knowledge, the impact of the asymmetric channel on the quantum teleportation has not yet been seriously explored.
The quantum channel, which is initially in the two-mode squeezed vacuum of squeezing factor $`s`$, is influenced by the thermal environments. Assuming that two thermal modes are independently coupled with the quantum channel, the dynamics of the quantum channel is described by a Fokker-Planck equation in the interaction picture . Solving the equation, the time-dependent Wigner function is obtained as
$`W_{ab}(\alpha _a,\alpha _b)=𝒩\mathrm{exp}\{`$ $``$ $`{\displaystyle \frac{2}{m_am_bc_ac_b}}[m_a|\alpha _a|^2+m_b|\alpha _b|^2`$ (2)
$`+`$ $`\sqrt{c_ac_b}(\alpha _a\alpha _b+\alpha _a^{}\alpha _b^{})]\},`$ (3)
where $`m_i=R_i(1+2\overline{n}_i)+T_i\mathrm{cosh}2s`$ and $`c_i=R_i\mathrm{sinh}2s`$, ($`i=a,b`$); $`\overline{n}_i`$ is the average thermal photon number of the environment for the channel mode $`i=a,b`$. The normalized interaction time $`R_i(1T_i)`$ is zero when the quantum channel is not subject to the environment and grows to unity when the channel completely assimilates the environment.
It has been shown that a two-mode Gaussian state is separable when a semi-positive well-defined $`P`$ function can be assigned to it after some local operations . The two-mode squeezed state subject to the thermal environment is separable when
$$(m_a1)(m_b1)c_ac_b.$$
(4)
As a special case, if the channel mode $`b`$ is influenced by the vacuum, i.e., $`\overline{n}_b=0`$, the channel becomes separable when
$$R_a\frac{1}{1+\overline{n}_a}.$$
(5)
For this case, the separation condition depends only on the average thermal photons influencing the channel mode $`a`$, even when the channel is minimally squeezed at the initial instance.
By the ideal teleportation, the original state is recovered at the receiving station. However, when the channel is not maximally entangled the teleportation is not ideal. The fidelity $``$ is defined as $`=\pi d^2\alpha W_o(\alpha )W_r(\alpha )`$, where $`W_o(\alpha )`$ and $`W_r(\alpha )`$ are the Wigner functions, respectively, for the original and teleported states, to show how close the teleported state is to the original state . One of the important assumptions in teleportation is that the original state is unknown so that the teleportation protocol has to be selected to optimize the average fidelity over all the possible original states. However, the average of the fidelity defined above is zero for continuous variables because of the scope of the possible original states . We thus take a subset composed of all coherent states and find the strategy to optimize the teleportation. A coherent state is one of a few manifestly quantum and extremely useful states generated in laboratories. Because the coherent states are nonorthogonal it is impossible to discern them with certainty. Any state can be written as a weighted sum of coherent states.
Let us consider the teleportation using the quantum channel (2). Before any action at the receiving and sending stations, the total state is a product of the original unknown state and the quantum channel. Setting homodyne detectors at the two output ports of the beam splitter, quadrature-phase variables $`p_1`$ and $`q_2`$ are measured at the respective output ports. Upon receiving each pair of measurement results $`g\sqrt{2}(q_2ip_1)`$, the quadrature variables of the channel field $`b`$ is displaced by $`g^{}(g)`$ accordingly. For a coherent original state $`|\alpha `$, we can calculate the measurement-conditioned Wigner function for the teleported state. Using the definition of the fidelity given above, $`(\alpha ,g,g^{})`$ is calculated. The average fidelity for the total set of coherent states is found:
$$\overline{}=\frac{d^2\alpha d^2g(\alpha ,g,g^{})}{d^2\alpha }=_o\frac{d^2g\mathrm{exp}[|gg^{}|^2]}{d^2\alpha }$$
(6)
where $``$ is a channel-dependent factor and
$$_o=\frac{1}{1+\frac{1}{2}(m_a+m_b)\sqrt{c_ac_b}}.$$
(7)
Eq.(6) readily shows that only when the displacement $`g^{}`$ is the same as $`g`$, the average fidelity may have a finite value and $`\overline{}=_o`$ is the optimum average fidelity in this case. Braunstein and Kimble found that the teleportation is optimized when $`g=g^{}`$ for the pure quantum channel and the fidelity $`_o`$ is greater than 1/2 when the pure channel is entangled . Using Eqs. (4) and (7), we find that for the symmetrically decohered mixed channel, i.e., $`m_a=m_b`$, the fidelity is greater than 1/2 as far as the channel is entangled. However, when the channel is asymmetrically decohered, the value 1/2 is not necessarily the critical fidelity for the standard teleportation scheme described above. Another important fact is found that the teleportation does not necessarily get better as the quantum channel is initially squeezed more. To illustrate more clearly, assume $`\overline{n}_b=0`$. For $`R_a=0`$ and $`R_b=1`$, the quantum channel is inseparable as shown in Eq.(5) but Eq.(6) gives the fidelity $`_o=1/(2+\mathrm{sinh}^2s)`$, which means that the more initially squeezed the quantum channel, the smaller the fidelity is.
In Fig. 1, the fidelity $`_o`$ is plotted against initial squeezing $`s`$ for the asymmetric quantum channel where the channel mode $`b`$ is influenced by the vacuum environment for the normalized interaction time $`R_b=0.01`$ and 0.05 while the channel mode $`a`$ is not influenced by an environment. The teleportation via the quantum channel, which has been decohered by the vacuum with interaction time $`R_b=0.01`$, in fact, corresponds to the teleportation with the pure squeezed quantum channel and imperfect displacement using the beam splitter of transmittance $`T=99\%`$. It is clearly seen that even for the quantum channel of seemingly-negligible asymmetry, if the channel is initially squeezed more than a certain degree, the teleportation becomes worse. By the first-derivative of $`_o`$ with regard to $`s`$, we find that the teleportation is optimized when the squeezing is $`\text{e}^{2s}=|t_at_b|/(t_a+t_b)`$ for a fixed channel condition $`t_a`$ and $`t_b`$, where $`t_i=\sqrt{T_i}`$ ($`i=a,b`$). Note that this result does not depend on the temperature of the environments.
For a mixed quantum channel, is the standard teleportation of the orthogonal measurement and the unitary transformation the best strategy? The straightforward answer is beyond the scope of this paper but the following argument gives some hint. Let us consider the asymmetrically mixed channel by taking $`\overline{n}_b=0`$ and $`\overline{n}_a0`$. For given initial squeezing and $`T_a`$, the fidelity is maximized to $`_o=1/[1+(1+\overline{n}_a)(1T_a)]`$ only when $`T_b=(\mathrm{sinh}2s/2\mathrm{sinh}^2s)^2T_a`$. We find that the critical fidelity 1/2 is recovered to coincide with the separation condition under this condition. This shows that some further decoherence may enhance the fidelity . The transformation accompanied by decoherence is no longer unitary which implies that a general transformation may optimize the fidelity. With the same argument, we find that a general measurement may optimize the fidelity for asymmetrically mixed quantum channel.
Why does it happen? One may intuitively expect that the more the quantum channel is squeezed, the better teleportation is. It is true when the channel remains pure and not influenced by an environment. However, when the channel is asymmetrically mixed, the conjecture may be wrong. There are many parameters which influences the teleportation. To make the analysis simple without losing the interesting features, we assume for the rest of the paper that the channel is exposed to low-temperature environments only for short periods of time, i.e., $`R_iT_i/\overline{n}_i`$. The separation condition (4) shows that the quantum channel remains entangled longer as it is initially squeezed more. However, as we have seen in Fig. 1, the fidelity of teleportation can be worse with increasing the initial squeezing. For the short interaction time with low-temperature environments, we write the Wigner function (2) to highlight the EPR correlation of the quantum channel as follows
$`W_{ab}(\alpha _a,\alpha _b)`$ $``$ $`𝒩\mathrm{exp}({\displaystyle \frac{2}{m_am_bc_ac_b}}[\text{e}^{2s}\{(t_bq_a+t_aq_b)^2+t_bp_at_ap_b)^2\}`$ (9)
$`+\text{e}^{2s}\{(t_bq_at_aq_b)^2+(t_bp_a+t_ap_b)^2\}]).`$
As the initial squeezing $`s\mathrm{}`$, the Wigner function $`W_{ab}C\delta (t_bq_a+t_aq_b)\delta (t_bp_at_ap_b)`$. It is clear that the asymmetric channel has the EPR correlation between the scaled quadrature variables; positions $`q_a`$ and $`q_b^{}=(t_a/t_b)q_b`$ and momenta $`p_a`$ and $`p_b^{}=(t_a/t_b)p_b`$. This is somewhat similar to the relation between an original phase space and amplified or dissipated phase space depending on the scale factor $`t_a/t_b`$. When the channel is strongly squeezed the channel entanglement teleports the original photon to the scaled space, which brings about large noise in the teleported state. For a large squeezing, even though the quantum channel is strongly entangled, the strong entanglement results in inefficient teleportation because the entanglement is between the scaled spaces.
The fact that the asymmetrically mixed quantum channel shows quantum correlation in differently-scaled spaces is not inherent in the two-dimensional qubit system where maximally correlated observables remain unchanged under the asymmetric decoherence. Consider a bipartite system in a spin-singlet state which is asymmetrically decohered in the phase-insensitive and isotropic environment as for the continuous variable system we discuss in this paper. Its correlation function for the spin variables along the directions a and b of two subsystems is given by $`\widehat{U}(𝐚)\widehat{\sigma }_z\widehat{U}^{}(𝐚)\widehat{U}(𝐛)\widehat{\sigma }_z\widehat{U}^{}(𝐛)\mathrm{cos}\theta _{ab}`$, where $`\widehat{U}`$ is a unitary operator and $`\theta _{ab}`$ the relative angle of the two vectors a and b. We clearly see that the maximal correlation between anti-parallel is preserved even after the asymmetric decoherence.
An imperfect detection efficiency at the receiving station can also be analyzed as a combination of perfect detection with a beam splitter in front. At the beam splitter, not only the channel state but also the unknown original field are mixed with the vacuum. This is why the average fidelity is maximized to $`=1/[\text{e}^{2s}+1/T_a]`$ for the displacement factor $`g^{}=g/t_a`$ assuming the channel mode $`b`$ is not subject to the environment. This shows that the teleportation is more efficient when the channel is initially squeezed more, which is in agreement with the earlier result .
We have shown that the important experimental errors can be absorbed by an imperfect mixed quantum channel while the experimental operations are assumed to be perfect. Because the experimental error does not occur symmetrically between the receiving and sending stations, it is important to study the influence of the asymmetrically-decohered quantum channel on the teleportation. For the asymmetrically-decohered quantum channel we found that the strong initial squeezing does not always optimize the teleportation because the asymmetric quantum channel has the EPR correlation between the differently-scaled phase spaces. We found that a measurement and transformation accompanied by decoherence optimizes the teleportation.
###### Acknowledgements.
This work was supported in part by the Brain Korea project (D-0055) of the Korean Ministry of Education.
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# Meromorphic groups.
## 1 Introduction
Let $`𝒜`$ be the category of reduced irreducible compact complex spaces. By a Zariski open subset of some $`X𝒜`$ we mean, as usual, the complement of an analytic subset of $`X`$. We understand a meromorphic mapping from $`X`$ to $`Y`$ ($`X,Y𝒜`$) in the sense of Remmert. Roughly speaking it is an analytic subset $`Z`$ of $`X\times Y`$, such that for some (nonempty, so dense) Zariski open subset $`U`$ of $`X`$, $`Z(U\times Y)`$ is the graph of a holomorphic function from $`U`$ to $`Y`$.
Fujiki, in his study of automorphism groups of compact Kähler manifolds, introduces the notion of a “meromorphic group”. As we will be proposing a less restrictive meaning for “meromorphic group” we will refer to Fujiki’s notion as “Fujiki-meromorphic”. A Fujiki-meromorphic group is a complex Lie group $`G`$ which is a Zariski open subset of some compact complex space $`G^{}`$ such that the group operation of $`G`$ extends to a meromorphic mapping from $`G^{}\times G^{}`$ to $`G^{}`$ which is holomorphic on $`(G\times G^{})(G^{}\times G)`$. Let $`𝒞`$ be the full subcategory of $`𝒜`$ consisting of compact complex spaces which are holomorphic images of compact Kähler manifolds. Fujiki proves that if $`G`$ is a Fujiki-meromorphic group in $`𝒞`$ (namely $`G^{}𝒞`$), then $`G`$ is “meromorphically” isomorphic to an extension of a complex torus by a linear algebraic group, generalizing Chevalley’s well-known theorem for algebraic groups. He raises the issue whether this remains true in the more general category $`𝒜`$ and proves it for $`G`$ commutative.
Our proposed definition of a meromorphic group is as follows: $`G`$ is a connected complex Lie group with a finite covering by Zariski open subsets $`U_i`$ of irreducible compact complex spaces $`X_i`$ ($`i=1,..,n`$) such that both the transition maps and the group operation on $`G`$ extend to meromorphic maps between the various $`X_i`$ and their products. Note that if the $`X_i`$ happen to be algebraic varieties then this agrees with the definition of an abstract algebraic group. Complex algebraic groups, complex tori, and Fujiki-meromorphic groups are all meromorphic groups. In fact our results imply that meromorphic groups coincide with Fujiki-meromorphic groups, and moreover have Kähler compactifications.
This paper is informed by model-theoretic concerns, and indeed some model-theoretic results play a role in the proofs. The main point is that $`𝒜`$ considered as a many-sorted structure, whose sorts are the compact complex spaces and whose basic relations are the analytic subsets of various Cartesian products of the sorts, is a structure with quantifier-elimination and finite Morley rank (sort by sort). This was proved by Zilber (although quantifier-elimination was also noted earlier in ). Quantifier-elimination says that the definable sets are precisely the finite unions of locally Zariski closed subsets of various compact complex spaces. It follows that definable functions are precisely “piecewise meromorphic” functions. Moreover definable groups are (up to definable isomorphism) precisely the meromorphic groups, giving another justification for introducing the notion of meromorphic group. A strongly minimal set in $`𝒜`$ is a definable set without infinite, co-infinite definable subsets. A strongly minimal group in $`𝒜`$ is precisely a meromorphic group without proper infinite Zariski-closed subsets. In it was noted that the deep results of apply to strongly minimal sets definable in $`𝒜`$, implying that any strongly minimal definable group $`G`$ is either an ($`1`$-dimensional) algebraic group, or is “modular”: every definable subset of $`G\times \mathrm{}\times G`$ is essentially a translate of a subgroup. Simple complex tori of dimension $`>1`$ are examples of strongly minimal modular groups (see ). For the converse, Hrushovski asked whether strongly minimal modular groups are (necessarily simple) complex tori. In fact Hrushovski outlined to the first author some ideas for proving this, depending however on finding a good compactification of the group. In any case, the question was answered by the second author in for the special case when $`G`$ is itself interpretable in a strongly minimal compact complex manifold. In additional observations about $`𝒜`$ and its model theory were made, including “elimination of imaginaries”. Also, it was asked whether the Chevalley theorem holds for groups definable in $`𝒜`$. We found subsequently that the same question was asked in for Fujiki-meromorphic groups.
We will prove the following results:
###### Theorem 1.1
Suppose $`G`$ is a strongly minimal meromorphic group. Then $`G`$ is meromorphically isomorphic to either a $`1`$-dimensional algebraic group or a simple nonalgebraic torus.
###### Theorem 1.2
Suppose $`G`$ is a connected meromorphic group. Then $`G`$ has a normal connected meromorphic subgroup $`L`$ such that $`L`$ is (meromorphically isomorphic to) a linear algebraic group, and $`G/L`$ is a complex torus. Moreover $`L`$ is unique.
Theorem 1.1 is a special case of Theorem 1.2. Theorem 1.1 will be proved by finding a good compactification of $`G`$ (i.e. showing that $`G`$ is Fujiki-meromorphic) and then (as $`G`$ is commutative) referring to . By again finding a suitable compactification we will first prove Theorem 1.2 for the special case when $`G`$ is an extension of a $`1`$-dimensional linear algebraic group by a simple complex torus. The general case will follow by an induction on dimension, making use of some additional ingredients such as the structure of compact complex spaces with algebraic co-dimension $`1`$, and some model theory of groups of finite Morley rank.
In the next section we give some definitions and recall both complex analytic and model-theoretic notions. In section 3 we carry out compactifications, proving Theorem 1.1 among other things. In section 4 we prove Theorem 1.2. Some additional remarks are made in section 5.
## 2 Preliminaries
For basic results, notions and notation concerning complex spaces and meromorphic maps, we refer the reader to , and . However we will repeat a few crucial definitions and results which we will be relying on. For us $`𝒜`$ denotes the class of reduced irreducible compact complex spaces. We take as given the notion of a holomorphic map $`f`$ from $`X`$ to $`Y`$ where $`X,Y𝒜`$. $`dim(X)`$ denotes the complex dimension of $`X`$. A modification of $`X𝒜`$ is some $`Y𝒜`$ and a surjective holomorphic $`f:YX`$ such that for some proper closed analytic subsets $`A`$ of $`Y`$ and $`B`$ of $`X`$, $`f|(YA):YAXB`$ is biholomorphic. Resolution of singularities says that any $`X`$ has a modification $`(Y,f)`$ such that $`Y`$ is nonsingular (so a connected compact complex manifold). The notion of a meromorphic mapping $`f`$ from $`X`$ to $`Y`$ ($`X,Y𝒜`$) is crucial. Such an object can be defined in various equivalent ways. For $`X`$ irreducible we define $`f`$ to be a function from $`X`$ to the set of subsets of $`Y`$ such that (i) the “graph” of $`f`$, $`\{(x,y)X\times Y:yf(x)\}`$ is an irreducible analytic subset of $`X\times Y`$, and for all $`x`$ in some (dense) Zariski open subset $`U`$ of $`X`$, $`f(x)`$ is a singleton. For a general $`X`$ we say that $`f`$ is meromorphic if each of its restrictions to the irreducible components of $`X`$ is meromorphic. We say that $`f`$ is holomorphic, or defined, at the points in $`U`$. Let $`Z`$ be the graph of $`f`$ as defined above, and $`\pi `$ the projection from $`Z`$ onto $`X`$. Then $`(Z,\pi )`$ turns out to be a modification of $`X`$. The projection of $`Z`$ on the second coordinate is then a holomorphic map from $`Z`$ to $`Y`$ which is said to be ”resolution of indeterminacies” of the meromorphic map $`f`$. From the definition of a meromorphic map one easily derives the following fact.
###### Fact 2.1
Let $`X,Y𝒜`$. Let $`f,g`$ be meromorphic mappings from $`X`$ to $`Y`$. Suppose that for some dense Zariski open subset $`U`$ of $`X`$, $`f`$ and $`g`$ agree on $`U`$. Then $`f=g`$.
Suppose that $`U`$ is a dense Zariski open subset of $`X𝒜`$, and $`f`$ a holomorphic map from $`U`$ into $`Y𝒜`$. By abuse of language we may sometimes say that $`f`$ is meromorphic if there is a meromorphic mapping $`g`$ from $`X`$ to $`Y`$ which agrees with $`f`$ on $`U`$. A natural category which can be associated to $`𝒜`$ is the category whose objects are those complex spaces which are Zariski open subsets of spaces in $`𝒜`$ and whose morphisms are the holomorphic maps which are meromorphic in the sense of the previous sentence. If we restrict our attention to those $`X`$, $`Y`$ which are projective algebraic varieties, this category is exactly that of quasiprojective varieties and morphisms.
It is also natural to consider complex spaces which have a finite covering by Zariski open subsets $`U_i`$ of spaces $`X_i`$ in $`𝒜`$ where the transition maps are meromorphic in the above sense. Morphisms in this category would be holomorphic maps which are meromorphic (in the above sense) when read in each $`U_i`$. What we will call a meromorphic group is exactly a group object in this latter category. Here is the precise definition.
###### Definition 2.2
A meromorphic group is a connected complex Lie group $`G`$, with a finite covering by open subsets $`W_i`$, for $`i=1,..,n`$, and for each $`i`$ a (biholomorphic) isomorphism $`\varphi _i`$ of $`W_i`$ with a Zariski open subset $`U_i`$ of some $`X_i𝒜`$ such that
* For each $`ij`$, $`\varphi _i(W_iW_j)`$ is a Zariski open subset of $`X_i`$, and the induced biholomorphic map between $`\varphi _i(W_iW_j)`$ and $`\varphi _j(W_iW_j)`$ is meromorphic, namely is the restriction of a meromorphic mapping between $`X_i`$ and $`X_j`$.
* For each $`i,j,k`$, $`\{(x,y)U_i\times U_j:\varphi _i^1(x)\varphi _j^1(y)W_k\}`$ is Zariski open in $`X_i\times X_j`$ and the induced holomorphic map ((x,y) goes to $`\varphi _k(\varphi _i^1(x)\varphi _j^1(y))`$) from $`U_i\times U_j`$ to $`U_k`$ is meromorphic, namely is the restriction of a meromorphic mapping between $`X_i\times X_j`$ and $`X_k`$.
Conditions (i) and (ii) can be expressed briefly by saying that the transition maps as well as the group operation are meromorphic when read in the various $`U_i`$ and their Cartesian products.
We say that the covering by the $`W_i`$’s and the isomorphisms with the $`U_i`$’s satisfying (i) and (ii) above give the complex Lie group $`G`$ a meromorphic structure.
If $`G`$ is a meromorphic group as in Definition 2.2, by a meromorphic subgroup $`H`$ of $`G`$ we mean a closed subgroup such that for each $`i`$, $`\varphi _i(HW_i)`$ is the intersection of an analytic subset of $`X_i`$ with $`U_i`$. Clearly $`H`$ has the structure of a meromorphic group.
A holomorphic homomorphism (complex Lie homomorphism) $`f`$ between meromorphic groups $`G_1`$ and $`G_2`$ is meromorphic if when restricted to the charts the map is meromorphic, that is, extends to meromorphic mappings between the relevant compact complex spaces.
So now we have the category of meromorphic groups and meromorphic homomorphisms. The following says that quotient objects exist. It follows from looking at the equivalent category of definable groups and using the elimination of imaginaries result from . This will be explained below.
###### Fact 2.3
Let $`G`$ be a meromorphic group and $`N`$ a normal meromorphic subgroup. Then there is a meromorphic group $`H`$ and a surjective meromorphic homomorphism from $`G`$ to $`H`$ whose kernel is $`N`$.
We now repeat the definition of a Fujiki-meromorphic group and recall what Fujiki proved.
###### Definition 2.4
Let $`G`$ be a complex Lie group.
* A meromorphic compactification of $`G`$ is a compact complex space $`G^{}𝒜`$ which contains $`G`$ as a dense Zariski open subset, such that the group operation $`\mu :G\times GG`$ is meromorphic, i.e. the restriction of a meromorphic mapping $`\mu ^{}`$ say, from $`G^{}\times G^{}`$ to $`G^{}`$
* A Fujiki-compactification of $`G`$ is a meromorphic compactification $`(G^{},\mu ^{})`$ of $`G`$ such that $`\mu ^{}`$ is holomorphic on $`(G\times G^{})(G^{}G)`$.
* $`G`$ is Fujiki-meromorphic if $`G`$ has a Fujiki-compactification.
###### Remark 2.5
(i) A Fujiki-meromorphic group is a meromorphic group.
(ii) A connected compact complex Lie group (i.e. a complex torus) is Fujiki-meromorphic.
(iii)(Remark 2.3 of .) A complex algebraic group is Fujiki-meromorphic.
Following notation of Fujiki :
###### Definition 2.6
We will call the meromorphic group $`G`$ regular if there is a meromorphic homomorphism $`f`$ from $`G^0`$, the connected component of the identity in $`G`$, onto a complex torus $`T`$ such that the kernel $`L`$ of $`f`$ is meromorphically isomorphic to a connected linear algebraic group. (Briefly said: $`G^0`$ is meromorphically an extension of a complex torus $`T`$ by a linear algebraic group $`L`$.)
###### Remark 2.7
* Let $`G`$ be regular and let $`T,L`$ be as above. Then $`L`$ and $`T`$ are unique. In particular $`L`$ is the unique maximal normal connected meromorphic subgroup of $`G^0`$ which is meromorphically isomorphic to a linear algebraic group.
* A regular meromorphic group is Fujiki-meromorphic.
Proof: (i) Suppose $`L_1`$ is a normal connected meromorphic subgroup of $`G`$ which is meromorphically isomorphic to a linear algebraic group. Then $`L_1/L`$ meromorphically embeds in $`T`$. So $`L_1/L`$ is both a complex torus and a linear algebraic group forcing it to be trivial. That is $`L_1`$ is contained in $`L`$.
(ii) As in Remark 2.3 of . $``$
Recall that $`𝒞`$ is the subclass (in fact full subcategory) of $`𝒜`$ consisting of those $`X`$ which are holomorphic images of compact connected Kähler manifolds. We will say that the connected meromorphic group $`G`$ is of type $`𝒞`$ if there is a Fujiki-compactification $`G^{}`$ of $`G`$ which is in $`𝒞`$. Fujiki proves:
###### Fact 2.8
(i) Suppose $`G`$ is Fujiki-meromorphic. Then $`G`$ is regular iff $`G`$ is of type $`𝒞`$.
(ii) Suppose that $`G`$ is commutative and Fujiki-meromorphic. Then $`G`$ is regular.
In the final part of this section we discuss the model theory of compact complex manifolds. We will have to assume the basics of model theory, and a bit more. is a good reference for basic model theory. The first four chapters of (by Bouscaren, Ziegler, Lascar, Pillay) are a useful reference for various aspects of applied and geometric stability theory. is an advanced text on geometric stability. deals with the theory of groups of finite Morley rank. Another good reference for stable groups is .
We consider $`𝒜`$ as a many-sorted first order structure whose sorts are the (reduced, irreducible) compact complex spaces and basic relations the analytic subsets of finite Cartesian products of such things.
###### Fact 2.9
$`Th(𝒜)`$ has quantifier-elimination, elimination of imaginaries and each sort has finite Morley rank. Moreover $`𝒜`$ is $`\mathrm{}_1`$-compact.
Quantifier-elimination was proved in , and independently in . It says that any definable subset of a sort $`X`$ is “analytically constructible”, that is a finite union of intersections of analytic (Zariski closed) subsets and complements of analytic (Zariski open) sets. A characterization of definable functions follows from this: Suppose $`U`$ is a definable set, and $`f`$ a definable function from $`X`$ into some sort $`Y`$. Then we can write $`X`$ as a disjoint union of definable sets $`U_i`$, where each $`U_i`$ is a Zariski open subset of some sort (complex space) $`X_i`$ such that for each $`i`$, the restriction of $`f`$ to $`U_i`$ is holomorphic and is the restriction to $`U_i`$ of a meromorphic mapping from $`X_i`$ to $`Y`$. We say, with possibly some abuse of language, that definable functions are piecewise meromorphic.
Zilber proved finiteness of Morley rank. Elimination of imaginaries was observed in .
$`\mathrm{}_1`$-compactness means that any countable family of definable subsets of some sort $`X`$ has nonempty intersection as long a every finite subfamily does.
With Fact 2.9 there is a remarkable parallel between complex-analytic and model-theoretic structural and classification results. We refer the reader to , , for more discussion. Note that on the face of it $`𝒜`$ is not $`\mathrm{}_1`$-saturated, as each element of each sort is essentially named by a constant. One can ask whether there is some sublanguage $`L_0`$ of the full language $`L`$ described above such that every relation in $`L`$ is definable possibly with parameters in the language $`L_0`$ and such that the reduct $`𝒜|L_0`$ is $`\mathrm{}_1`$-saturated. This is not true, as, for example, a general generalized Hopf surface has continuum many holomorphic automorphisms but our Proposition 5.2 shows that it has trivial generic type and hence cannot have a non-trivial definable family of automorphisms. On the other hand, $`𝒞`$ can be considered as a reduct of $`𝒜`$ (fewer sorts but the full structure on each sort), which has quantifier-elimination and elimination of imaginaries in its own right. Fujiki’s results on the Douady spaces of manifolds in $`𝒞`$ imply that the structure $`𝒞`$ is $`\mathrm{}_1`$-saturated in a suitable sublanguage. (See .)
$`𝒜^{}`$ will denote a very saturated elementary extension of $`𝒜`$. For any $`X𝒜`$, $`X^{}`$ denotes its canonical extension in $`𝒜^{}`$. We will often work model-theoretically in $`𝒜^{}`$. For example, a definable property holds generically on $`X`$ iff it holds of a generic point of $`X^{}`$.
A definable group $`G`$ in $`𝒜`$ will be called connected if $`G`$ has no definable subgroups of finite index. Any meromorphic group is clearly a definable group (using elimination of imaginaries). Methods from the algebraic case due to Hrushovski and van den Dries (see as well as Pillay’s article in ) adapt to yield the important:
###### Fact 2.10
Any definable connected group $`G`$ in $`𝒜`$ is definably isomorphic to a connected meromorphic group $`H`$ (unique up to meromorphic isomorphism).
This fact gives a natural equivalence between the category of definable groups and meromorphic groups. In particular any definable homomorphism between meromorphic groups will be meromorphic (in particular holomorphic). From here on we will use “definable” interchangeably with “meromorphic” when talking about groups and homomorphisms.
A definable set $`X`$ (in $`𝒜`$) is said to be strongly minimal if $`X`$ is infinite and has no infinite coïnfinite definable subsets. A definable connected group $`A`$ is said to be modular if every definable subset of $`A^n`$ is a Boolean combination of translates of definable subgroups. In it was proved that the results of apply to the category $`𝒜`$. This yields.
###### Fact 2.11
Suppose $`G`$ is a definable connected group in $`𝒜`$ which has no infinite normal definable subgroups. Then either $`G`$ is strongly minimal and modular or $`G`$ is definably isomorphic to a (complex) algebraic group.
It follows that if $`T`$ is a nonalgebraic simple complex torus then $`T`$ is modular. (A direct proof, avoiding was given in .)
## 3 Compactifications
We will prove:
###### Theorem 3.1
Let $`G`$ be a connected commutative meromorphic group, which is either strongly minimal, or an extension of a connected $`1`$-dimensional linear algebraic group by a simple complex torus. Then $`G`$ is Fujiki-meromorphic.
A consequence is:
###### Corollary 3.2
* Let $`G`$ be a strongly minimal meromorphic group. Then $`G`$ is meromorphically isomorphic to either a $`1`$-dimensional algebraic group or a simple modular complex torus.
* Let $`G`$ be a commutative meromorphic group which is an extension of a $`1`$-dimensional linear algebraic group by a simple complex torus. Then $`G`$ meromorphically splits.
Proof: (i) $`G`$ is commutative, so by Theorem 3.1, Fujiki meromorphic, thus by Fact 2.8 (ii), meromorphically an extension of a complex torus $`T`$ by a linear algebraic group $`L`$. As $`G`$ is strongly minimal, $`G`$ is either $`T`$ or $`L`$. If $`G=L`$, then $`dim(L)=1`$. If $`G=T`$, then $`T`$ has no proper infinite analytic subsets so is either an elliptic curve or simple and modular (by 2.11).
(ii) Immediate, by Fact 2.8 (ii). $``$
To prove Theorem 3.1 we will find a meromorphic compactification $`G^{}`$ of $`G`$ and then show it to be a Fujiki-compactification. The following general result concerning compactifications of commutative meromorphic groups will be useful.
###### Lemma 3.3
Suppose that the connected commutative meromorphic group $`(G,\mu )`$ has meromorphic compactification $`(G^{},\mu ^{})`$. Suppose $`S=G^{}G`$ is nonempty. Then
* Every component of $`S`$ has co-dimension $`1`$ in $`G^{}`$.
* $`\mu ^{}|(G^{}\times S)`$ is a meromorphic mapping from $`G^{}\times S`$ to $`S`$.
* For each $`gG`$, and component $`C`$ of $`S`$, $`\mu _g^{}=\mu ^{}(g,):CC`$ is biholomorphic on a dense Zariski open subset of $`C`$, and for $`g,hG`$, $`\mu _g^{}.\mu _h^{}=\mu _{g.h}^{}`$ on a dense Zariski open subset of $`C`$.
Proof: (i) Let $`n=dim(G^{})`$ (=$`dim(G)`$). Suppose for the sake of contradiction that there is $`xS`$ such that $`dim_x(S)<n1`$. Let $`\mathrm{\Delta }_n`$ be the open unit disc in $`^n`$ and $`f:\mathrm{\Delta }_nU`$ be a coordinate function for any open neighborhood $`U`$ of $`x`$ in $`G^{}`$ where $`U`$ is chosen such that $`US`$ has dimension $`<n1`$. So if $`A=f^1(US)`$, then $`A`$ is an analytic subset of $`\mathrm{\Delta }_n`$ of codimension at least $`2`$, and $`f_1`$ = $`f|(\mathrm{\Delta }_nA)`$ is a holomorphic embedding into $`G`$. As $`G`$ is a connected commutative Lie group its universal cover is $`\pi :^nG`$. As $`A`$ has co-dimension at least two in $`\mathrm{\Delta }_n`$, $`\mathrm{\Delta }_nA`$ is simply connected, and so $`f_1`$ lifts to a holomorphic map $`f_2:\mathrm{\Delta }_n^n`$ (see ). Let $`g=\pi f_2`$. Then $`g`$ is a holomorphic map from $`\mathrm{\Delta }`$ into $`G^{}`$ which agrees with $`f`$ off the thin analytic subset $`A`$. But then by Riemann’s removable singularity theorem () $`f=g`$, contradicting the fact that $`xG`$. (i) is proved.
(ii) Let $`\mathrm{\Gamma }`$ be the graph of $`\mu ^{}`$. We will first show that for all $`(g,x)`$ in some dense Zariski open subset $`V`$ of $`G^{}\times S`$, $`\{y:(g,x,y)\mathrm{\Gamma }\}`$ is finite. If not, then for a Zariski open subset $`V`$ of $`G^{}\times S`$, the above set of $`y`$’s has positive dimension. It follows from (i) that $`dim(\mathrm{\Gamma }(G^{}\times S\times G^{})2n`$, contradicting irreducibility of $`\mathrm{\Gamma }`$. It now follows by the implicit function theorem that
(a) $`\mu ^{}`$ is holomorphic on $`V`$.
For $`gG`$ let $`\mu _g^{}`$ be $`\mu ^{}(g,)`$, a meromorphic mapping from $`G^{}`$ to $`G^{}`$. Note that if $`\mu _g^{}`$ is defined (single valued) at $`x`$ and $`\mu _h^{}`$ is defined at $`\mu _g^{}(x)`$ then $`\mu _{hg}^{}`$ is defined at $`x`$ and equals $`\mu _h^{}(\mu _g^{}(x))`$. It follows that if $`(g,x)V`$ and $`gG`$, then $`\mu (g,x)S`$. As $`G`$ is Zariski-dense in $`G^{}`$ it follows that
(b) $`\mu ^{}|(G^{}\times S)`$ is a meromorphic mapping into $`S`$, yielding (ii).
(iii) The same argument as above shows that for any $`gG`$, $`\mu _g^{}|S`$ is a meromorphic mapping from $`S`$ to $`S`$. Let $`C_1,..,C_s`$ be the irreducible components of $`S`$. Note that the image of the meromorphic mapping $`\mu _g^{}`$ from $`G^{}`$ to $`G^{}`$ (i.e. projection of its graph on second component) is all of $`G^{}`$. But for $`xG`$, $`\mu _g^{}(x)G`$. Thus the image of the meromorphic mapping $`\mu _g^{}|S`$ is all of $`S`$. We work model-theoretically. Fix $`C_i`$. Let $`x`$ be a generic point of $`C_i^{}`$ over $`𝒜`$. So $`y=\mu _g^{}(x)C_j^{}`$ for some $`j=f(i)`$. It follows that $`\mu _g^{}|C_i`$ is a meromorphic mapping from $`C_i`$ into $`C_{f(i)}`$.
(c) Thus $`f=f_g`$ must be a permutation of $`\{1,\mathrm{},s\}`$.
If for some $`i`$, and $`x`$ as above, $`\mu _g^{}(x)`$ is not a generic point of $`C_{f(i)}^{}`$ over $`𝒜`$, then there is a proper analytic subset $`D_{f(i)}`$ of $`C_{f(i)}`$ such that $`\mu _g^{}|C_i`$ has image contained in $`D_{f(i)}`$. By (c), we contradict the fact that $`\mu _g^{}|S`$ has image all of $`S`$.
Thus for $`xC_i^{}`$ generic, $`\mu _g^{}(x)`$ is generic in $`C_{f(i)}^{}`$ over $`𝒜`$. It follows that $`gf_g`$ gives a definable action of $`G`$ on $`\{1,\mathrm{},s\}`$. As $`G`$ is connected this has to be trivial. This gives (iii). $``$
###### Remark 3.4
(iii) above is interpreted model-theoretically by saying that $`G`$ acts generically on $`C`$: let $`p=p_C`$ be the generic type of the component $`C`$ of $`S`$. Then for $`gG^{}`$ and $`x`$ realizing $`p`$ independent of $`g`$ (over $`𝒜`$), $`\mu ^{}(g,x)`$ is defined, realizes $`p`$ and is independent from $`g`$. Moreover, if $`g,hG^{}`$ and $`x`$ realizes $`p`$ independent of $`g,h`$ then $`\mu ^{}(hg,x)=\mu ^{}(h,\mu ^{}(g,x))`$.
We can now obtain the strongly minimal case of Theorem 3.1.
###### Lemma 3.5
Let $`G`$ be a strongly minimal meromorphic group. Then $`G`$ is Fujiki-meromorphic.
Proof:
*Step 1.* Finding a meromorphic compactification.
By assumption on $`G`$ some open nonempty definable subset $`U`$ of $`G`$ is already a Zariski open subset of a compact complex space $`X`$, which we may assume by resolution of singularities to be a manifold. Moreover by strong minimality of $`G`$, $`GU`$ is finite, say $`\{g_1,\mathrm{},g_n\}`$. For $`i=1,\mathrm{},n`$ let $`V_i`$ be a coordinate neighborhood of $`g_i`$ in $`G`$ such that the closures $`\overline{V}_i`$ of the $`V_i`$ in $`G`$ are disjoint. Note that $`K_i=\overline{V}_i\{g_i\}`$ is contained in $`U`$, so in $`X`$, but is not compact, so not closed in $`X`$. Let $`D_i`$ be the boundary of $`K_i`$ in $`X`$, namely $`\overline{K}_iK_i`$. Let $`\pi :XX^{}`$ be the quotient map which collapses each $`D_i`$ to a point $`c_i`$. Then $`X^{}`$ is compact, $`\pi `$ is holomorphic (in fact is a modification), and is biholomorphic outside the union of the $`D_i`$’s. Let $`f:GX^{}`$ be defined by $`f(x)=\pi (x)`$ for $`xU`$ and $`f(g_i)=c_i`$. Then $`f`$ is a definable, holomorphic embedding.
*Step 2.* $`G^{}`$ is a Fujiki-compactification of $`G`$.
If $`G=G^{}`$ there is nothing to do. Otherwise, (as $`G`$ is commutative) Lemma 3.3 applies. Let $`S`$ be as there. We will show that $`G`$ is holomorphic on $`S`$, and in fact acts as the identity. Note that the generic type of $`G`$ is orthogonal to any set of dimension less than that of $`G`$ ($`G`$ being strongly minimal). In particular $`G`$ is orthogonal to $`S`$. Fix a component $`C`$ of $`S`$. Lemma 3.3 (iii) gives us a generic action of $`G`$ on $`C`$. Let $`g,hG^{}`$ be generic independent elements of $`G`$ and let $`x`$ be generic in $`C^{}`$ over $`\{g,h\}`$. Then by the orthogonality mentioned above, each of $`g`$ and $`h`$ is independent from $`\{x,\mu ^{}(g,x)\}`$. It follows that $`\mu ^{}(g,x)=\mu ^{}(h,x)`$, and thus $`\mu ^{}(h^1.g,x)=x`$. But $`h^1.g`$ is generic in $`G^{}`$ and independent from $`x`$. It follows that $`G`$ acts generically trivially on $`C`$. So the holomorphic map from $`G^{}\times C`$ to $`C`$ taking $`(g,x)`$ to $`x`$ agrees generically with the meromorphic mapping $`\mu ^{}|(G^{}\times C):G^{}\times CC`$. By Fact 2.1, these mappings agree. This shows that $`(G^{},\mu ^{})`$ is a Fujiki-compactification of $`G`$. $``$
We now deal with the case when $`G`$ is a commutative extension of the additive group, $`𝔾_a`$, or the multiplicative group, $`𝔾_m`$, by a simple complex torus $`T`$. We let $`H`$ denote $`G/T`$ (so $`H`$ is $`𝔾_a`$ or $`𝔾_m`$). If $`G`$ is meromorphically isomorphic to an algebraic group, then $`G`$ is clearly Fujiki-meromorphic (in fact the Chevalley theorem applies immediately, yielding Theorem 1.2). If $`T`$ has a definable complement in $`G`$ (up to finite), then again we get the required conclusion. So for the rest of this section we make:
*Assumption.*
* $`G`$ is a commutative meromorphic group, which is meromorphically an extension of $`H`$ by a simple complex torus $`T`$, where $`H`$ is $`𝔾_a`$ or $`𝔾_m`$.
* $`G`$ is not meromorphically isomorphic to an algebraic group.
* There is no definable connected infinite subgroup $`L`$ of $`G`$ with $`LT`$ finite.
We will show that $`G`$ is Fujiki-meromorphic (which actually leads to a contradiction).
We will make use of the “socle theory” from .
###### Lemma 3.6
$`T`$ is the maximal almost pluriminimal definable subgroup of $`G`$.
Proof: Note that $`T`$ being simple is almost strongly minimal. So if the lemma, as $`G/T`$ has dimension $`1`$, $`G`$ is semipluriminimal. By , $`G`$ is an almost direct product of pairwise orthogonal semiminimal groups. If $`G`$ is already semiminimal, then as $`G`$ is nonorthogonal to $`^1`$ via $`GH`$, $`G`$ must be algebraic, contradicting Assumption (b) Thus $`G`$ is the semidirect product of $`T`$ and some $`L`$, contradicting Assumption (c). $``$
###### Lemma 3.7
Let $`X`$ be a definable subset of $`G`$. Assume that the Morley rank of $`X`$ is strictly less than the Morley rank of $`G`$ (equivalently $`X`$ is not Zariski-dense in $`G`$). Then $`X`$ is contained in finitely many translates of $`T`$.
Proof: We prove the lemma by induction on $`RM(X)=m`$. It is clearly true for $`m=0`$. We may assume that the Morley degree of $`X`$ is $`1`$. Let $`S`$ be the (model-theoretic) stabilizer of $`X`$. If $`S`$ is finite then by Lemma 3.6, as well as Proposition 4.3 of , $`X`$ is, up to a set of Morley rank $`<m`$, contained in a single translate of $`T`$. By induction hypothesis, $`X`$ is contained in finitely many translates of $`T`$ as desired. So we may assume that $`S`$ is infinite. By Assumption (c), and the fact that $`T`$ is simple, $`S`$ must contain $`T`$. Note that $`RM(T)=RM(G)1RM(X)`$. But it is well-known that the Morley rank of the stabilizer of a Morley degree $`1`$ set $`X`$ is at most the Morley rank of $`X`$, with equality if and only if the stabilizer is connected and $`X`$ is, up to a set of smaller Morley rank, a translate of this stabilizer. Thus $`RM(S)=RM(X)`$, $`S=T`$ and up to a set of smaller Morley rank, $`X`$ is a translate of $`T`$, so we finish again by induction. $``$
###### Lemma 3.8
$`G`$ is Fujiki-meromorphic.
Proof: As in the strongly minimal case we first find a compact complex manifold $`G^{}`$ containing $`G`$ as a Zariski open set, and then show that this gives $`G`$ a Fujiki-meromorphic structure.
*Step I.* Finding the compactification.
Let $`RM(G)=n`$. By definition of $`G`$ being a meromorphic group, let $`U`$ be a definable subset (with Morley rank $`n`$) of $`G`$ which is a dense Zariski-open subset of a compact complex manifold $`\overline{U}`$. Let $`\pi :GH`$ be the canonical surjective homomorphism. Then $`\pi `$ takes $`U`$ onto a cofinite subset $`\pi (U)`$ of $`H`$.
*Claim 1.* We may assume that for any $`x\pi (U)`$, $`\pi ^1(x)U`$ = $`\pi ^1(x)`$ (a translate of $`T`$).
Proof: $`Y`$ = $`\pi ^1(\pi (U))U`$ is a definable subset of $`G`$ of Morley rank $`<n=RM(G)`$. By Lemma 3.7, $`Y`$ is contained in finitely many translates of $`T`$, namely finitely many fibers of $`\pi `$. Remove these fibers from $`U`$. $``$
Let $`\pi ^{}`$ denote $`\pi |U`$. $`\pi ^{}`$ extends to a meromorphic function $`\overline{\pi }`$ from $`\overline{U}`$ to $`^1`$. Further restricting $`U`$ we may assume:
*Claim 2.* For all $`x\pi ^{}(U)`$, $`(\pi ^{})^1(x)=\overline{\pi }^1(x)`$.
Let $`C`$ be the finite set $`H\pi (U)`$. Then we can find $`h\pi (U)`$ such that $`h.C\pi (U)`$. Let $`gU`$ be a preimage of $`h`$. Let $`\tau _g:GG`$ be multiplication by $`g`$. $`\tau _g|U`$ is not defined everywhere but is holomorphic on the open set where it is defined and so extends to a meromorphic map $`\overline{\tau }_g:\overline{U}\overline{U}`$. By a theorem of Remmert (see Theorem 1.9 in Chapter VII of ), there is a modification $`\nu :\stackrel{~}{U}\overline{U}`$ and a holomorphic map $`\stackrel{~}{\tau }:\stackrel{~}{U}\overline{U}`$ such that $`\stackrel{~}{\tau }=\overline{\tau }_g\nu `$. In particular, for $`x\stackrel{~}{U}`$ such that $`\tau _g|U`$ is defined at $`\nu (x)`$, $`\stackrel{~}{\tau }(x)=\tau _g(\nu (x))`$.
*Claim 3.* $`\overline{\pi }\stackrel{~}{\tau }=\tau _h\overline{\pi }\nu `$.
Proof: This holds generically, so holds everywhere. $``$
We will now construct the required compactification $`G^{}`$ of $`G`$ as a holomorphic image of $`\stackrel{~}{U}`$. Let $`S=^1H`$. So $`S=\{\mathrm{}\}`$ or $`\{\mathrm{},0\}`$. As a set $`G^{}`$ will be the disjoint union of $`G`$ with $`\overline{\pi }^1(S)`$. The manifold structure of $`G^{}`$ is as follows: $`G`$ is given its canonical manifold structure. Now let $`x\overline{\pi }^1(S)`$. Let $`y=\overline{\pi }(x)^1`$. Choose an open neighborhood $`V`$ of $`y`$ in $`^1`$ such that $`V\{y\}U`$. Then $`\overline{\pi }^1(V)G^{}`$ is an open neighborhood of $`x`$. The transition maps are clearly holomorphic, yielding a structure of a complex compact manifold on $`G^{}`$ containing $`G`$ as an open (dense) subset.
Now we define a holomorphic surjective map $`f`$ from $`\stackrel{~}{U}`$ to $`G^{}`$. Let $`x\stackrel{~}{U}`$. If $`\overline{\pi }(\stackrel{~}{\tau }(x))C`$ define $`f(x)=\stackrel{~}{\tau }(x)`$ (so $`f(x)\overline{\pi }^1(\pi (U)S)G^{}`$). On the other hand, if $`\overline{\pi }(\stackrel{~}{\tau }(x))C`$, define $`f(x)=g.\nu (x)`$. Note that in this latter case $`\nu (x)UG`$ and so $`g.\nu (x)G`$ and $`\pi (g.\nu (x))=\overline{\pi }\stackrel{~}{\tau }(x)`$.
It is easy to check, given our assumptions, that $`f`$ is holomorphic and surjective. So $`G^{}`$ is a compact complex manifold, containing $`G`$ as a dense Zariski open subset (the embedding of $`G`$ in $`G^{}`$ is definable and holomorphic).
This completes Step I.
*Step II:* The action of $`G`$ on itself extends to a trivial action on the boundary $`G^{}G`$.
Let $`C_1,\mathrm{},C_k`$ be the irreducible components of $`G^{}G`$. Note that for each $`i`$, $`dim(C_i)=dim(T)`$.
*Claim 4.* For each $`i`$ there is a surjective holomorphic map from $`C_i`$ to $`T`$ (so finite-to-one outside a proper Zariski closed subset $`D_i`$ of $`C_i`$).
Proof: By Step I we have a surjective holomorphic map $`\pi :G^{}^1`$ such that $`\pi ^1(H)=G`$ and $`\pi |G`$ is precisely the canonical surjective homomorphism from $`G`$ to $`H`$. So $`G^{}G`$ is $`\pi ^1(S)`$ where $`S=^1H`$. Consider the map $`\mu (g,h)=gh^1`$ from $`G\times _HG`$ to $`T`$. This is definable and holomorphic so extends to a meromorphic map from $`G^{}\times _^1G^{}`$ to $`T`$, which we also call $`\mu `$. By Lemma 3.3 of this map is holomorphic. In particular, for any $`C_i`$ and $`xC_i`$, $`\mu (,x)|C_i`$ is a holomorphic map from $`C_i`$ into $`T`$. We must show that for suitable $`xC_i`$, this is surjective.
For $`gG`$, let $`f_g`$ be the meromorphic map from $`G^{}`$ to $`G^{}`$ whose restriction to $`G`$ is multiplication by $`g`$ (so $`f_g`$ is $`\overline{\tau }_g`$ in previous notation).
Let $`tT`$. The map taking $`xG`$ to $`\mu (tx,x)T`$ is the constant map with value $`t`$. It follows that whenever $`xG^{}`$ and $`f_t`$ is single valued at $`x`$ then $`\mu (f_t(x),x)=t`$. Choose $`x_0`$ generic in $`C_i`$. Then for a dense open set $`V`$ of $`t^{}s`$ in $`T`$, $`f_t(x_0)`$ is defined and in $`C_i`$. So for each $`tV`$, $`\mu (f_t(x_0),x_0)=t`$. Thus $`\mu (,x_0)|C_i:C_iT`$ is generically surjective, so surjective. In any case this map is finite. $``$
By Lemma 3.3, for each $`i`$, $`f_g`$ induces a generic holomorphic action of $`G`$ on $`C_i`$. Let $`K_i`$ be the subgroup of $`G`$ consisting of those $`gG`$ such that for all $`x`$ in some Zariski open subset of $`C_i`$, $`f_g(x)=x`$. For dimension reasons $`K_i`$ is a definable infinite subgroup of $`G`$, so contains $`T`$. Moreover we have an induced generic action of $`G/K_i`$ on $`C_i`$. Let $`D_i`$ be as in Claim 4. Let $`x_0C_i`$ be such that for any $`h`$ in some dense Zariski open subset of $`G/K_i`$, $`h.x_0C_iD_i`$. This gives a meromorphic map from $`G/K_i`$ to $`C_i`$ whose image contains infinitely many points outside $`D_i`$. Composing with the holomorphic map from $`C_i`$ to $`T`$ given by Claim 4 yields a meromorphic nonconstant map from $`^1`$ into $`T`$ which is impossible.
Thus $`K_i=G`$ and the generic action of $`G`$ on $`C_i`$ is trivial. This holds for each $`i`$. So the generic action of $`G`$ on $`G^{}G`$ is trivial. That is, if $`gG`$, then the meromorphic mapping $`f_g:G^{}(G^{}G)`$ agrees with the identity map on a dense Zariski open set. This implies that $`f_g`$ is the identity map. So $`G^{}`$ witnesses $`G`$ being Fujiki meromorphic. The proof of Lemma 3.8 is complete as well as the proofs of Theorem 3.1 and Corollary 3.2 $``$
## 4 Composition series
In this section we will prove
###### Theorem 4.1
Suppose $`G`$ is a connected meromorphic group. Then $`G`$ is regular (in the sense of Definition 2.6).
We first state a consequence of Theorem 3.1
###### Proposition 4.2
Suppose the connected meromorphic group $`G`$ is simple, in the sense that $`G`$ has no nontrivial, connected normal definable subgroup. Then $`G`$ is either
* a (almost simple) noncommutative algebraic group,
* $`𝔾_a`$ or $`𝔾_m`$,
* a simple abelian variety, or
* a strongly minimal modular complex torus.
Proof: Simplicity of $`G`$ together with the the dichotomy theorem from , implies that $`G`$ is either nonorthogonal to $`^1`$ (namely has nontrivial algebraic reduction) or $`G`$ is modular. In the first case, $`G`$ is an algebraic group, so (i), (ii) or (iii) hold. In the second case, every definable subset of $`G`$ is a Boolean combination of cosets of definable subgroups. Simplicity implies that $`G`$ is strongly minimal. By 3.2, $`G`$ is a complex torus. $``$
The following will be crucial. The special case for Fujiki-meromorphic groups in the class $`𝒞`$ was proved in . In any case the “classical” theory of groups of finite Morley rank enters the picture.
###### Lemma 4.3
Let $`1LGH1`$ be an exact sequence of connected meromorphic groups and suppose that $`L`$ and $`H`$ are (meromorphically isomorphic to) linear algebraic groups. Then so is $`G`$.
Proof: Note that if $`G`$ satisfies the hypotheses of the lemma and $`G_1`$ is a connected definable subgroup of $`G`$, or an image of $`G`$ under a meromorphic homomorphism then $`G_1`$ satisfies the hypotheses too (for suitable $`L_1,H_1`$).
We will prove the lemma by induction on $`dim(G)=n`$. We consider various possibilities for $`G`$,
*Case I.* $`G`$ has an an infinite center.
By the hypotheses, $`Z(G)`$ contains an infinite definable linear algebraic group and thus, by the structure of commutative linear algebraic groups, $`Z(G)`$ contains a definable $`1`$-dimensional connected linear algebraic group $`A`$. $`A`$ is normal in $`G`$, so by induction hypothesis $`G/A`$ is (meromorphically isomorphic to) a linear algebraic group. It makes sense to talk about the algebraic dimension $`a(G)`$ of $`G`$. Note that $`dim(G/A)=n1`$ so the map $`GG/A`$ witnesses that $`a(G)n1`$. If $`a(G)=n`$, then $`G`$ is already isomorphic to an algebraic group, so a linear algebraic group. Otherwise $`a(G)=n1`$, and it is well-known (see ) that the general fiber of the algebraic reduction $`\pi :GX`$ is an elliptic curve $`E`$. But the map $`GG/A`$ must meromorphically factor through $`\pi `$. The general fiber of the first map is $`^1`$ and thus we see that $`E`$ is an image of $`^1`$ under a meromorphic (i.e. rational) map, a contradiction. Thus $`G`$ is linear algebraic.
Case II. $`G`$ is solvable.
We may assume, by Case I, that $`Z(G)`$ is finite. But then $`G/Z(G)`$ is centerless and easily $`G`$ is linear algebraic iff $`G/Z(G)`$ is. So we may assume that $`G`$ is centerless. As is shown in Chapter 3 of or Chapter 9 of , the commutator subgroup $`G^{}`$ of $`G`$ is connected, and nilpotent, so $`Z(G^{})`$ is infinite and contains a minimal definable connected $`G`$-normal subgroup $`A`$. $`G/G^{}`$ defines an infinite group of automorphisms of $`A`$. By again results in or , $`A`$ is the additive group of a definable field $`K`$. As $`A`$ is by hypothesis linear algebraic, the field $`K`$ has to be (definably isomorphic to) $``$ and $`dim(A)=1`$. $`G/A`$ is by induction hypothesis algebraic, and as in Case 1 we deduce that $`G`$ is (definably) algebraic.
Case III. $`G`$ is nonsolvable.
Note that $`G`$ is among other things a connected complex Lie group, and as such we have the Levi-Malcev decomposition $`G=R.S`$ where $`R`$ is the maximal normal solvable connected subgroup of $`G`$, $`S`$ is a semisimple (complex) Lie group (unique up to conjugacy in $`G`$), and $`RS`$ is discrete. Finite Morley rank considerations (see 5.38 in ) show that $`R`$ is definable, so linear algebraic by Case II (or induction). It is probably then well-known that $`G`$ must be isomorphic as a complex Lie group to a linear algebraic group. However we want $`G`$ to be definably isomorphic to such a group. So we must do a little more work, although maybe there is a more direct way. We will show
Claim. $`S`$ is a definable subgroup of $`G`$.
Proof: We may assume that $`R`$ is a proper, nontrivial subgroup of $`G`$, definably isomorphic to a linear algebraic group. We will first reduce to the case where $`R`$ is commutative and unipotent. Let $`H`$ be the connected component of the center of the commutator subgroup of $`R`$. $`H`$ is then a nontrivial commutative connected linear algebraic group, normal in $`G`$. So $`H`$ is the direct product $`U.T`$ of a commutative unipotent group $`U`$ and an algebraic torus $`T`$. Note that both $`U`$ and $`T`$ are definable connected normal subgroups of $`G`$. $`T`$ has no infinite definable group of automorphisms, so is central in $`G`$. By Case I we may assume $`T`$ to be trivial. Thus $`H=U`$ is unipotent. By induction hypothesis, $`G/H`$ is linear algebraic. Clearly $`R/H`$ is the solvable radical of $`G/H`$. Thus $`G/H`$ is an almost direct product of $`R/H`$ with a semisimple algebraic group $`G_1/H`$ (where $`G_1`$ is a definable connected subgroup of $`G`$ containing $`H`$). As $`S`$ is unique up to conjugacy we may assume that $`G_1=H.S`$. Note that the homomorphism $`\mu :G_1G_1/H`$ is an isomorphism on $`S`$.
Note that $`G_1/H`$ is linear algebraic, by induction hypothesis among other things. Now $`S`$ (being semisimple) is isomorphic (uniquely) as a complex Lie group to a linear algebraic group, so it makes sense to talk about an element of $`S`$ being unipotent. $`S`$ is an almost direct product of almost simple groups $`S_1,.,S_r`$. Fix a nontrivial unipotent element $`a`$ in some $`S_iH`$. Work now inside the definable group $`G_1`$. Let $`a_1=\mu (a)G_1/H`$. $`a_1`$ is then unipotent, and let $`U_1`$ be a $`1`$-dimensional definable unipotent subgroup of $`G_1/H`$ containing $`a_1`$. Let $`U_2=\mu ^1(U_1)`$. Then $`U_2`$ is an extension of a unipotent linear algebraic group ($`U_1`$) by a linear algebraic unipotent group $`H`$ so is (by induction) linear algebraic unipotent. $`aU_2`$. We can find a definable commutative connected subgroup $`U_3`$ of $`U_2`$ containing $`a`$. $`U_3`$ is definably a vector space over $`𝐂`$. The $`1`$-dimensional subspace $`U_4`$ generated by $`a`$ is a definable subgroup of $`G`$ contained in $`S_i`$. We have found an infinite connected subgroup $`U_4`$ of $`S_i`$ which is definable in $`G`$. The group generated by all the $`U_4^g`$, where $`gS_i`$ is definable and must be equal to $`S_i`$. So $`S_i`$ is definable. As $`i`$ was arbitrary, $`S`$ is definable. The claim is proved. $``$
We want $`S`$ to be definably isomorphic to a linear algebraic group. Recall that $`S`$ as a complex Lie group is the almost direct product of almost simple (discrete centre) groups $`S_1,..,S_r`$. As the centre of $`S`$ is definable, each $`S_i`$ has finite centre. By considering centralizers, each $`S_i`$ is definable. $`S_i`$ being almost simple is almost strongly minimal, hence by the validity of the Zilber trichotomy in $`𝒜`$, modular or nonorthogonal to $`^1`$. But $`S_i`$ is nonabelian. So it must be nonorthogonal to $`^1`$, so algebraic. Thus $`S`$ is definably an algebraic group, so by semismplicity, linear algebraic.
Finally $`G`$, being the almost semidirect product of linear algebraic $`R`$ with linear algebraic $`S`$, must be linear algebraic. Case III is complete, as well as the proof of Lemma 4.3.
$``$
Proof of Theorem 4.1. The proof will by induction on $`dim(G)`$.
We first deal with the case when $`G`$ is commutative. Let $`H`$ be a minimal definable connected subgroup of $`G`$. By 4.2, $`H`$ is either (a) a linear algebraic group or (b) a simple complex torus. If $`G=H`$ we are finished. Otherwise, applying the induction hypotheses, $`G/H`$ is definably an extension of a complex torus $`T`$ by a linear algebraic group $`L/H`$. In case (a) by Lemma 4.3, $`L`$ is linear algebraic. So $`G`$ is definably an extension of $`T`$ by $`L`$, and we finish. So suppose (b) holds. If $`L/H`$ is trivial, $`G`$ is an extension of a complex torus by a complex torus, so also a complex torus. Otherwise let $`L_1/H`$ be a $`1`$-dimensional subgroup of $`L/H`$. Then $`L`$ is definably an extension of $`𝔾_a`$ or $`𝔾_m`$ by the simple complex torus $`T`$. By Corollary 3.2 (ii) $`L`$ splits, yielding a $`1`$-dimensional linear algebraic subgroup of $`G`$. We are now back in case (a). This proves Case I.
We now deal with the general case.
If $`G`$ has no proper normal nontrivial definable connected subgroup, then we are finished by 4.2
Otherwise let $`H`$ be a proper normal nontrivial connected subgroup of $`G`$. The induction hypothesis applies to $`H`$. If the maximal connected linear algebraic normal subgroup $`L`$ of $`H`$ is nontrivial, then $`L`$ is normal in $`G`$ as $`L`$ is characteristic in $`H`$ by Remark 2.7(i), and by applying the induction hypothesis to $`G/L`$ and applying Lemma 4.3 we finish. Otherwise $`H`$ is a complex torus. As a complex torus has no infinite definable group of automorphisms, $`H`$ is central in $`G`$. The induction hypothesis applies to $`G/H`$. If the latter is a complex torus so is $`G`$. Otherwise $`G`$ has a connected normal definable subgroup $`G_1`$ containing $`H`$ such that $`G_1/H`$ is linear algebraic. If $`G_1/H`$ is semisimple (equivalently contains no infinite normal solvable subgroup), then $`G_1`$ is the almost direct product of its commutator subgroup $`G_1^{}`$ and $`H`$. $`G_1^{}`$ is semisimple so (definably) linear algebraic. We conclude by applying the induction hypothesis to $`G/G_1^{}`$ and using Lemma 4.3.
If $`G_1/H`$ is not semisimple, then there is a definable nontrivial connected subgroup $`A`$ of $`G_1`$ containing $`H`$ such that $`A`$ is normal in $`G`$ and $`A/H`$ is commutative (and linear algebraic). $`A/H`$ is (definably) a product of $`𝔾_a`$’s and $`𝔾_m`$’s. Remember that $`H`$ is central in $`A`$, so if $`A`$ is not commutative then the commutator map yields a nonconstant meromorphic map from $`A/H\times A/H`$ into the complex torus $`H`$, which is impossible. So $`A`$ has to be commutative. By induction hypothesis, or by the first part of the proof (if $`A=G`$), $`A`$ has a definable connected (nontrivial) subgroup which is normal in $`G`$ and (definably) linear algebraic. By induction hypothesis and Lemma 4.3 we finish.
This completes the proof of Theorem 4.1. $``$
## 5 Additional remarks and questions
The following was first proved by the second author . We give a quick proof using Theorem 1.1.
###### Proposition 5.1
Let $`X`$ be a strongly minimal compact complex manifold. Then $`X`$ is either a (smooth projective) algebraic curve, a complex torus, or is trivial.
Proof: Suppose that $`X`$ is neither trivial, nor an algebraic curve. Then $`dim(X)>1`$ and by and , there is a strongly minimal group $`G`$ definable in $`X`$. By Theorem 1.1 $`G`$ must be a complex torus, nonorthogonal to $`X`$. Nonorthogonality is witnessed by an analytic subset $`\mathrm{\Gamma }`$ of $`X\times A`$ which projects generically finite-to-one on each of $`X`$ and $`G`$. As both $`X,G`$ are strongly minimal, both projections $`\pi _1:\mathrm{\Gamma }X`$ and $`\pi _2:\mathrm{\Gamma }G`$ are finite-to-one, and $`\mathrm{\Gamma }`$ is strongly minimal with $`dim(\mathrm{\Gamma })=dim(X)=dim(G)>1`$. Replacing $`\mathrm{\Gamma }`$ by its normalization we may assume $`\mathrm{\Gamma }`$ is normal. $`\mathrm{\Gamma }`$ has no proper infinite analytic subsets, in particular no co-dimension $`1`$ analytic subsets. By the purity of branch theorem, $`\pi _2`$ is an unramified covering. Thus $`\mathrm{\Gamma }`$ is a complex torus. As $`\pi _1`$ is also an unramified covering $`X`$ is a complex torus. $``$
We can more generally give natural necessary and sufficient conditions for a type of $`U`$-rank $`1`$ to be trivial. First we call a compact complex manifold (or space) $`X`$ simple if there is no definable family $`\{Y_t:tT\}`$ of positive dimensional proper analytic subvarieties of $`X`$ such that $`\{Y_t:tT\}`$ contains a Zariski open subset of $`X`$. It is easy to see that $`X`$ is simple if and only if its generic type $`p_X`$ has $`U`$-rank $`1`$. A Kummer manifold is a compact complex space which is bimeromorphic with a space of the form $`T/G`$ where $`T`$ is a complex torus and $`G`$ a finite group of (holomorphic) automorphisms of $`T`$.
The following proposition explains the trichotomy within simple compact complex spaces between algebraic curves, nonalgebraic Kummer manifolds, and manifolds of zero algebraic and Kummer dimension in terms of the Zilber trichotomy for Zariski geometries.
###### Proposition 5.2
Suppose $`X`$ is a simple compact complex manifold. Then $`p_X`$ is trivial if and only if
* the algebraic dimension $`a(X)`$ of $`X`$ is $`0`$, and
* there is no surjective meromorphic map from $`X`$ to a Kummer manifold (i.e. $`k(X)=0`$ in the notation of ).
In particular, if $`X𝒞`$ then $`p_X`$ is trivial.
Proof: ($``$): suppose $`p_X`$ is nonorthogonal to $`^1`$. Then clearly $`a(X)>0`$. Suppose $`p_X`$ is nontrivial and modular. Then $`p_X`$ is nonorthogonal to the generic type of a strongly minimal modular torus $`T`$. Let $`a`$ be a generic point of $`X^{}`$ (i.e. a realization of $`p_X`$). Then there is generic $`bT^{}`$ in $`acl(a)`$. Let $`\{b_1,\mathrm{},b_n\}`$ be the finite set of realizations of $`tp(b/a)`$. Then by the modularity of $`T`$, $`(b_1,\mathrm{},b_n)`$ is a generic point of a translate of a (strongly minimal) subtorus $`S`$ of $`T^n`$. After translating we may assume that $`(b_1,\mathrm{},b_n)`$ is a generic point of $`S`$. By elimination of imaginaries the finite set $`\{b_1,\mathrm{},b_n\}`$ is coded by some $`c`$. As $`cdcl(a)`$, we obtain a surjective meromorphic map $`f`$ from $`X`$ to a compact complex manifold $`Y`$ where $`c`$ is a generic point of $`Y`$. But the map taking $`(b_1,\mathrm{},b_n)`$ to $`c`$ extends to a meromorphic map from $`S`$ to $`Y`$. Modularity of $`S`$ implies that this induces a bimeromorphic map between $`S/G`$ and $`Y`$ for some finite group $`G`$ of automorphisms of $`S`$. Thus $`Y`$ is Kummer. So if (i) and (ii) hold, the only possibility left for $`p_X`$ is to be trivial.
($``$): If (i) fails then $`p_X`$ is nonorthogonal to $`^1`$ so is nontrivial. If (ii) fails and there is a surjective meromorphic map to the Kummer manifold $`Y`$ then clearly $`Y`$ is also simple and its generic type is nontrivial. So $`p_X`$ is nontrivial.
The “in particular” clause follows from the observations that any Kummer manifold $`Y`$ is in $`𝒞`$ as well as any compact complex manifold which maps meromorphically and generically finite-to-one on $`Y`$. $``$
Our classification of meromorphic groups together with Fact 2.8 (i) shows that any meromorphic group $`G`$ is of “type $`𝒞`$”, that is already definable in the structure $`𝒞`$. As $`𝒞`$ is saturated in a suitable language, our results also classify definable groups in all models of $`Th(𝒞)`$. What about groups definable in elementary extensions of $`𝒜`$? Here is a conjecture:
###### Conjecture 5.3
Let $`G`$ be a definable group in $`𝒜^{}`$. Then $`G`$ is definably (in $`𝒜^{}`$) isomorphic to a group $`H`$ definable in the reduct $`𝒞^{}`$.
The conjecture can be restated in terms of families of groups in $`𝒜`$.
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# Hierarchy of Critical Exponents on Sierpinski fractal resistor networks
## 1 INTRODUCTION
The study of infinite sets of exponents which originated in the field of turbulence , has recently become the focus of attention in a number of fields involving fractals or scaling objects , ranging from random resistor networks , dynamical systems, diffusion limited aggregates (DLA) , to localization. What is common to all of these different fields is that, one wants to characterize the properties of a ”weight” or ”measure” associated to different parts of a fractal object. Modelization of electrical transport properties for inhomogeneous and composite materials by random resistor networks have been the subject of many recent works, Also other physical phenomena such as diffusion problems can be formulated in terms of electrical problem. Distribution of currents (or voltage drops) on a percolating structure in the scaling region are multifractals, in the sense that different moments scale with different exponents, that is, if we consider a system of length $`L`$, then the $`q`$-moment of the current distribution :
$$M_q=\underset{r}{}I_r^^q$$
(1-1)
scales as $`L^{D_q}`$, where $`D_q`$ is by no means a simple function of $`q`$. Thus each moment scales with its own anomalous dimension. This phenomena is characteristic of multifractals distributions. Actually this set of exponents first appeared in the field of turbulence and has recently become focus of attention in a number different fields such as diffusion limited aggregation, dynamical system and random resistor networks as mentioned above. Here in this paper we study the multifractals structure of current distribution on Sierpinsky fractal, since as Kirkpatric had suggested, the so called back-bone of the percolating random resistor networks could be modeled by a fractal structure and among the fractal objects, the $`n`$-simplex one is simplest to study the various physical problems from random walk to electrical one on it .
Here by using the $`S_3`$-symmetry of Sierpinski fractal resistor networks (see Fig. 1) together with the minimization of the electrical power, we have been able to determine the current distribution in Sierpinsky fractals with decimation numbers $`b=2,3,4`$, and $`5`$. Then, using the independent Shure’s $`S_3`$ invariant polynomials, which is proved that the required number of independent Shure’s $`S_3`$ -invariant polynomials of degree $`q`$ is $`[q/4]+1`$, with \[ \] indicating the greatest integer parts, we have derived the results of reference for $`b=2`$ up to $`q=12`$ and we have calculated $`D_q`$ up to $`q=22`$ for $`b=2,3,4`$ and $`5`$. The organization of the article is as follows:
In section 2, we give a brief description of Sierpinski fractals, then in section 3, using the $`S_3`$-symmetry of Sierpinsky fractal resistor networks and minimization of electrical power we have determined the inward flowing current of subfractals. In section 4 we talk about the independent Shure,s $`S_3`$-symmetry invariant polynomials of input currents. Section 5 is about the moments current distributions and their multifractals spectrum where it contains the main results of this paper. The paper ends with a brief conclusion and 5 appendices.
## 2 Sierpinski Fractal
To obtain Sierpinski fractal with decimation number $`b`$, we choose a triangle and divide its sides into $`b`$ parts and then draw all possible lines through the dividend points parallel to the side of the triangle. Next, having omitted every other inner triangle, we repeat this for the remaining triangles or for the subfractals of the next higher order. This way Sierpinski fractals are constructed. To calculate the fractal dimension, we label subfractals of order $`(l+1)`$ in terms of partition of $`(b1)`$ into $`3`$ positive integers $`\lambda __1`$, $`\lambda __2`$ and $`\lambda __3`$. Each partition represents a subfractal of order $`l`$ and $`\lambda `$ shows the distance of the corresponding subfractal from the sides of triangle. As an illustrating example, we show in Figure 2 the method of labeling a Sierpinski fractal with decimation number $`b=3`$. Obviously, the number of all possible partitions is equal to the distribution of $`(b1)`$ objects among three boxes, which is the same as the Bose-Einstein distribution of $`(b1)`$ identical bosons in 3 quantum states. This is equal to
$$c=\frac{(b+1)!}{(b1)!2!}.$$
According to the following definition, the fractal dimension $`d_F`$ of a self-similar object is
$$(N^r)^{d_F}=1,$$
with $`N`$ as the number of similar objects, up to translation and rotation. For self-similar fractals, $`N`$ is equal to the number of subfractals. Therefore, we have $`N=C^l`$ and $`r=b^l`$. Hence $`d_F=\frac{lnC}{lnb}`$, or
$$d_F=\frac{ln(b+1)!/(b1)!2!}{lnb}.$$
## 3 Determination of inward flowing current of subfractals
We denote the j-th inward flowing current of subfractal which corresponds to the partition $`\lambda _1,\lambda _2`$ and $`\lambda _3`$ by $`I_{\lambda _1,\lambda _2,\lambda _3}(\lambda _1+\delta _{1,j},\lambda _2+\delta _{2,j},\lambda _3+\delta _{3,j})`$. Therefore, $`I_1`$, $`I_2`$ and $`I_3`$ can be denoted by
$$I_{b1,0,0}(b,0,0),I_{0,b1,0}(0,b,0)\text{and}I_{0,0,b1}(0,0,b).$$
In order to determine the inward flowing currents in terms of $`I_j,j=1,2`$ and $`3`$, besides using Kirchhoff’s law at each node, we have to minimize the power of Sierpinski fractal, that is we minimize the following expression:
$$\underset{overpossiblepartitions}{}\underset{j=1}{\overset{3}{}}I_{\lambda _1,\lambda _2\lambda _3}^2(\lambda _1+\delta _{1,j},\lambda _2+\delta _{2,j},\lambda _3+\delta _{3,j})$$
$$+\mu _{\lambda _1,\lambda _2,\lambda _3}(\underset{j=1}{\overset{3}{}}I_{\lambda _1,\lambda _2,\lambda _3}(\lambda _1+\delta _{1,j},\lambda _2+\delta _{2,j},\lambda _3+\delta _{3,j})$$
$$+2\underset{overallnodes}{}\nu _{\eta _1,\eta _2,\eta _3}(\underset{j=1}{\overset{3}{}}I_{\eta _1,\eta _2,\eta _3}(\eta _1\delta _{1,j},\eta _2\delta _{2,j},\eta _3\delta _{3,j}).$$
(3-1)
where $`\mu _{\lambda _1,\lambda _2,\lambda _3}`$ and $`\nu _{eta_1,eta_2,\eta _3}`$ are Lagrange multipliers which are considered because of Kirchhoff,s law on each subfractal, and also each node, respectively. By minimizing the energy given by expression (1-2), we get the following equations for the inner flowing currents:
$$I_{\lambda _1,\lambda _2,\lambda _3}(\lambda _1+\delta _{1,j},\lambda _2+\delta _{2,j},\lambda _3+\delta _{3,j})+\mu _{\lambda _1,\lambda _2,\lambda _3}+nu_{\lambda _1+\delta _{1,j},\lambda _2+\delta _{2,j},\lambda _3+\delta _{3,j}}=0,$$
(3-2)
together with the Kirchhoff’s law for each subfractal and each vertex, respectively:
$$\underset{j=1}{\overset{3}{}}I_{\lambda _1,\lambda _2,\lambda _3}(\lambda _1+\delta _{1,j},\lambda _2+\delta _{2,j},\lambda _3+\delta _{3,j})=0,$$
(3-3)
$$\underset{j=1}{\overset{3}{}}I_{\eta _1,\eta _2,\eta _3}(\eta _1\delta _{1,j},\eta _2\delta _{2,j},\eta _3\delta _{3,j})=0.$$
(3-4)
Now, using the $`S_3`$-permutation symmetry of Sierpinski fractal we propose the following ansatz for the Lagrange multipliers:
$$\mu _{\lambda _1,\lambda _2,\lambda _3}=\underset{k=1}{\overset{3}{}}a_{\lambda _k}I_k,$$
(3-5)
$$\nu _{\lambda _1,\lambda _2,\lambda _3}=\underset{k=1}{\overset{3}{}}b_{\lambda _k}I_k,$$
(3-6)
where $`a_0`$ is assumed to be zero. Substituting the ansatz (3-5) and (3-6) in equation (3-2), then the inward flowing currents can be given in terms of a’s and b,s , respectively.
$$I_{\lambda _1,\lambda _2,\lambda _3}(\eta _1,\eta _2,\eta _3)=\underset{k=1}{\overset{3}{}}(a_{\lambda _k}+b_{\eta _k})I_k.$$
(3-7)
Actually one could write the currents in terms of input ones as in (3-7) simply by using the symmetry of simplex fracal where the minimization of power is not required. Finally the a’s and b’s themselves can be determined through the equations (3-3) and (3-4). Here we determine the currents for b=2,3,4 and 5, respectively. First for b=2 we have
$$I_{1,0,0}(2,0,0)=I_1,I_{0,1,0}(0,2,0)=I_2,I_{0,0,1}(0,0,2)=I_3,$$
$$I_{\delta _{1,j},\delta _{2,j},\delta _{3,j}}(\delta _{1,j}+\delta _{1,k},\delta _{2,j}+\delta _{2,k}\delta _{3,j}+\delta _{3,k})=a_1I_J+b_1I_j+b_1I_k$$
Using equation (3-4) we obtain:
$$a_1+2b_1=0$$
and using equation (3-4) we get:
$$1+2a_1+b_1=0.$$
Solving the above equations we get the following result:
$$I_{\delta _{1,j},\delta _{2,j},\delta _{3,j}}(\delta _{1j}+\delta _{1k},\delta _{2j}+\delta _{2k},\delta _{3j}+\delta _{3k})=\frac{(I_kI_j)}{3}.$$
Similarly for b=3 we have
$$I_{2,0,0}(3,0,0)=I_1,I_{0,2,0}(0,3,0)=I_2,I_{0,0,2}(0,0,3)=I_3,$$
$$I_{2\delta _{1,j},2\delta _{2,j},2\delta _{3,j}}(2\delta _{1,j}+\delta _{1,k},2\delta _{2,j}+\delta _{2,k},2\delta _{3,j}+\delta _{3,k})=a_2I_j+b_2I_j+b_1I_k,$$
$$I_{\delta _{1,j}+\delta _{1,k},\delta _{2,j}+\delta _{2,k},\delta _{3,j}+\delta _{3,k}}(2\delta _{1,j}+\delta _{1,k},2\delta _{2,j}+\delta _{2,k},2\delta _{3,j}+\delta _{3,k})=a_1I_j+a_1I_k+b_2I_j+b_1I_k,$$
$$I_{\delta _{1,j}+\delta _{1,k},\delta _{2,j}+\delta _{2,k},\delta _{3,j}+\delta _{3,k}}(\delta _{1,j}+\delta _{1,k}+\delta _{1,l},\delta _{2,j}+\delta _{2,k}+\delta _{2,l},\delta _{3,j}+\delta _{3,k}+\delta _{3,l})$$
$$=a_1(I_j+I_k)+b_1(I_j+I_k+I_l).$$
Using equation (3-3) in subfractal $`(2\delta _{1,j},2\delta _{2,j},2\delta _{3,j})`$, we get
$$1+2(a_2+b_2)b_1=0.$$
Also using equation (3-3) in subfractal $`(\delta _{1,j}+\delta _{1,K},\delta _{2,j}+\delta _{2,K},\delta _{3,j}+\delta _{3,K})`$ we get
$$3a_1+2b_1+b_2=0.$$
Also using equation (3-4) in the vertices we have
$$a_1+a_2+2b_1=0,$$
$$a_1+2b_1=0,$$
$$2a_1+3b_1=0.$$
After solving the above equations we get the following result for the currents for b=3
$$I_{2\delta _{1,j},2\delta _{2,j},2\delta _{3,j}}(2\delta _{1,j}+\delta _{1,k},2\delta _{2,j}+\delta _{2,k},2\delta _{3,j}+\delta _{3,k})=\frac{9}{21}I_j+\frac{3}{21}I_k,$$
$$I_{\delta _{1,j}+\delta _{1,k},\delta _{2,j}+\delta _{2,k},\delta _{3,j}+\delta _{3,k}}(2\delta _{1,j}+\delta _{1,k},2\delta _{2,j}+\delta _{2,k},2\delta _{3,j}+\delta _{3,k})$$
$$=\frac{9}{21}I_j\frac{3}{21}I_k,$$
$$I_{\delta _{1,j}+\delta _{1,k},\delta _{2,j}+\delta _{2,k},\delta _{3,j}+\delta _{3,k}}(\delta _{1,j}+\delta _{1,k}+\delta _{1,l},\delta _{2,j}+\delta _{2,k}+\delta _{2,l},\delta _{3,j}+\delta _{3,k}+\delta _{3,l})$$
$$=\frac{6}{21}(I_j+I_k)+\frac{4}{21}(I_j+I_k+I_l).$$
By the same procedure explained above, we can calculate the inner inward flowing currents for decimation number $`b=4`$ and $`b=5`$, where we quote only the results below and give the details of calculation in Appendix I and II.
### 3.1 Inner inward flowing currents corresponding to $`b=4`$:
$$I_{3\delta _{1,j},3\delta _{2,j},3\delta _{3,j}}(3\delta _{1,j}+\delta _{1,k},3\delta _{2,j}+\delta _{2,k},3\delta _{3,j}+\delta _{3,k})=\frac{19}{41}I_j+\frac{3}{41}I_k,$$
$$I_{2\delta _{1,j}+\delta _{1,k},2\delta _{2,j}+\delta _{2,k},2\delta _{3,j}+\delta _{3,k}}(3\delta _{1,j}+\delta _{1,k},\delta _{2,j}+\delta _{2,k},3\delta _{3,j}+\delta _{3,k})=\frac{19}{41}I_j\frac{3}{41}I_k,$$
$$I_{2\delta _{1,j}+\delta _{1,k},2\delta _{2,j}+\delta _{2,k},2\delta _{3,j}+\delta _{3,k}}(2\delta _{1,j}+2\delta _{1,k},2\delta _{2,j}+2\delta _{2,k},2\delta _{3,j}+2\delta _{3,k})$$
$$=\frac{9}{41}(I_jI_k),$$
$$I_{2\delta _{1,j}+\delta _{1,k},2\delta _{2,j}+\delta _{2,k},2\delta _{3,j}+\delta _{3,k}}(2\delta _{1,j}+\delta _{1,k}+\delta _{1,l},2\delta _{2,j}+\delta _{2,k}+\delta _{2,l},2\delta _{3,j}+\delta _{3,k}+\delta _{3,l})$$
$$=\frac{184}{1353}I_j\frac{52}{1353}I_k+\frac{146}{1353}I_l,$$
$$I_{\delta _{1,j}+\delta _{1,k}+\delta _{1,l},\delta _{2,j}+\delta _{2,k}+\delta _{2,l},\delta _{3,j}+\delta _{3,k}+\delta _{3,l}}(2\delta _{1,j}+\delta _{1,k}+\delta _{1,l},2\delta _{2,j}+\delta _{2,k}+\delta _{2,l},2\delta _{3,j}+\delta _{3,k}+\delta _{3,l})$$
$$=\frac{368}{1353}I_j\frac{94}{1353}(I_k+I_l).$$
### 3.2 Inner inward flowing currents corresponding to $`b=5`$:
$$I_{4\delta _{1,j},4\delta _{2,j},4\delta _{3,j}}(4\delta _{1,j}+\delta _{1,k},4\delta _{2,j}+\delta _{2,k},4\delta _{3,j}+\delta _{3,k})=\frac{283}{591}I_j\frac{41375}{1015929}I_k,$$
$$I_{3\delta _{1,j}+\delta _{1,k},3\delta _{2,j}+\delta _{2,k},3\delta _{3,j}+\delta _{3,k}}(4\delta _{1,j}+\delta _{1,k},4\delta _{2,j}+\delta _{2,k},4\delta _{3,j}+\delta _{3,k})$$
$$=\frac{283}{591}I_j+\frac{41375}{1015929}I_k,$$
$$I_{3\delta _{1,j}+\delta _{1,k},3\delta _{2,j}+\delta _{2,k},3\delta _{3,j}+\delta _{3,k}}(3\delta _{1,j}+2\delta _{1,k},3\delta _{2,j}+2\delta _{2,k},3\delta _{3,j}+2\delta _{3,k})$$
$$=\frac{51}{197}I_j\frac{25}{197}I_k,$$
$$I_{3\delta _{1,j}+\delta _{1,k},3\delta _{2,j}+\delta _{2,k},3\delta _{3,j}+\delta _{3,k}}(3\delta _{1,j}+\delta _{1,k}+\delta _{1,l},3\delta _{2,j}+\delta _{2,k}+\delta _{2,l},3\delta _{3,j}+\delta _{3,k}+\delta _{3,l})$$
$$=2\frac{17486}{112881}I_j+2\frac{2206}{112881}I_k2\frac{2448}{37627}I_l,$$
$$I_{2\delta _{1,j}+2\delta _{1,k},2\delta _{2,j}+2\delta _{2,k},2\delta _{3,j}+2\delta _{3,k}}(3\delta _{1,j}+2\delta _{1,k},3\delta _{2,j}+2\delta _{2,k},3\delta _{3,j}+2\delta _{3,k})$$
$$=\frac{51}{197}I_j\frac{25}{197}I_k,$$
$$I_{2\delta _{1,j}+2\delta _{1,k},2\delta _{2,j}+2\delta _{2,k},2\delta _{3,j}+2\delta _{3,k}}(2\delta _{1,j}+2\delta _{1,k}+\delta _{1,l},2\delta _{2,j}+2\delta _{2,k}+\delta _{2,l},2\delta _{3,j}+2\delta _{3,k}\delta _{3,l})$$
$$=\frac{9865}{338643}(2I_j+I_k)\frac{12482}{338643}I_l$$
$$I_{2\delta _{1,j}+\delta _{1,k}+\delta _{1,l},2\delta _{2,j}+\delta _{2,k}+\delta _{2,l},2\delta _{3,j}+\delta _{3,k}+\delta _{3,l}}(3\delta _{1,j}+\delta _{1,k}+\delta _{1,l},3\delta _{2,j}+\delta _{2,k}+\delta _{2,l},3\delta _{3,j}+\delta _{3,k}+\delta _{3,l})$$
$$=2\frac{5138}{112881}(I_k+I_l)4\frac{34972}{112881}I_j,$$
$$I_{2\delta _{1,j}+\delta _{1,k}+\delta _{1,l},2\delta _{2,j}+\delta _{2,k}+\delta _{2,l},2\delta _{3,j}+\delta _{3,k}+\delta _{3,l}}(2\delta _{1,j}+2\delta _{1,k}+\delta _{1,l},2\delta _{2,j}+2\delta _{2,k}+\delta _{2,l},2\delta _{3,j}+2\delta _{3,k}+\delta _{3,l})$$
$$=2\frac{18847}{3047787}I_j2\frac{14356}{3047787}I_k+\frac{624}{3047787}I_l.$$
## 4 Shure’s Polynomials of Inward Flowing Currents
Shure’s $`S_3`$-invariant polynomials are homogeneous polynomials of degree 3 of variables $`I_1,I_2`$ and $`I_3`$:
$$s_{\lambda _1,\lambda _2,\lambda _3}=\underset{permutationof(1,2,\text{3})}{}I_1^{\lambda _1}I_2^{\lambda _2}I_3^{\lambda _3}$$
where $`\lambda _1,\lambda _2,\lambda _3`$ are partitions of m into $`3`$ non-negative integers, that is:
$$\lambda _1+\lambda _2+\lambda _3=m.$$
Because of the following equation due to Kirchhoff’s law:
$$S_1=\underset{k=1}{\overset{3}{}}I_k=0,$$
(4-1)
all Schure’s polynomials of degree m, corresponding to all possible partitions of m, are not independent. In calculation of the multifractals critical exponents $`D_q`$, we must use the independent ones. By multiplying both sides of (4-2) by $`S_{\lambda _1,\lambda _2,\lambda _3}`$, we get
$$0=S_1S_{\lambda _1,\lambda _2,\lambda _3}=a_{\mu _1,\mu _2,\mu _3}S_{\mu _1,\mu _2,\mu _3},$$
(4-2)
where $`(\mu _1,\mu _2,\mu _3)`$ and $`(\lambda _1,\lambda _2,\lambda _3)`$ correspond to partition of $`m`$-1 and $`m`$ respectively. From the formula (2-10) it follows that there are $`P_3(m+1)`$ constraint over $`P_3(m)`$ shure polynomials of degree m, where $`p_3(m)`$ takes all possible partitions of $`m`$ into 3 non-negative integers. Therefore, the number of invariant polynomials of degree m is:
$$P_3(m)P_3(m1).$$
(4-3)
For example for $`m=2`$ we have
$$0=S_1S_1=S_2+2S_{1,1}$$
therefore, using the above equation we can write $`S_{1,1}`$ in terms of $`S_2`$ as:
$$S_{1,1}=\frac{S_2}{2}.$$
(4-4)
Thus we have only one invariant polynomial for $`q=2`$. Also in the case of $`q=4`$ we have
$$S_1S_3=S_4+S_{3,1}=0,$$
$$S_1S_{2,1}=S_{3,1}+2S_{2,2}+S_{2,1,1}=0,$$
$$S_1S_{1,1,1}=S_{2,1,1}=0,$$
hence there is only one independent polynomial such as $`S_4`$ and the others can be written
in terms of $`S_4`$ as follows:
$$S_{3,1}=S_4,S_{2,2}=\frac{S_4}{2},$$
$$S_{2,1,1}=0.$$
(4-5)
For $`q=6`$ we have:
$$S_1S_5=S_6+S_{5,1}=0,$$
$$S_1S_{4,1}=S_{5,1}+S_{4,2}+S_{4,1,1}=0,$$
$$S_1S_{3,2}=S_{4,2}+2S_{3,3}+S_{3,2,1}=0,$$
$$S_1S_{3,1,1}=S_{4,1,1}+S_{3,2,1}=0.$$
Therefore, there are only two independent invariant polynomials such as $`S_6`$, $`S_{3,3}`$ and the other dependent one can be written in terms of them as follows:
$$S_{5,1}=S_6,S_{4,2}=\frac{S_6}{2}S_{3,3},S_{321}=\frac{S_6}{2}S_{3,3},S_{4,1,1}=\frac{S_6}{2}+S_{3,3}$$
In Appendix III, we have proved that the number of independent Schure’s invariant polynomials of degree $`q`$ is equal to:
$$[q/4]+1$$
(4-6)
where \[ \] means the greatest integer part.
Below we give some of the constraints over Schure’s invariant polynomials of degrees $`q=8`$ and $`10`$ which are occurring through imposing the Kirschhof’s law over Schure’s polynomials of order eight:
$$S_8+S_{7,1}=0,S_{7,1}+S_{6,2}+2S_{6,1,1}=0,$$
$$S_{6,2}+S_{5,3}+S_{5,2,1}=0,S_{6,1,1}+S_{5,2,1}=0,$$
$$S_{5,3}+2S_{4,4}+S_{4,3,1}=0,S_{5,2,1}+S_{4,3,1}+2S_{4,2,2}=0,$$
where, $`S_8`$ and $`S_{4,2,2}`$ are considered as the invariant polynomials and other dependent invariant can be expressed in terms of them as follows:
$$S_{7,1}=S_8,S_{3,3,2}=S_{4,2,2},$$
$$S_{4,3,1}=S_{4,2,2},S_{5,2,1}=5S_{4,2,2},$$
$$S_{6,1,1}=5S_{4,2,2},S_{6,2}=S_8+10S_{4,2,2},$$
$$S_{5,3}=15S_{4,2,2}S_8,S_{4,4}=\frac{S_8+16S_{4,2,2}}{2}.$$
Constrains over Schure’s polynomials of order ten are:
$$S_{10}+S_{9,1}=0,S_{8,2}+S_{9,1}+2S_{8,1,1}=0,$$
$$S_{8,2}+S_{7,3}+S_{7,2,1}=0,S_{8,1,1}+S_{7,2,1}=0,$$
$$S_{7,3}+S_{6,4}+S_{6,3,1}=0,S_{7,2,1}+S_{6,3,1}+2S_{6,2,2},$$
$$S_{6,4}+2S_{5,5}+S_{5,4,1}=0S_{6,3,1}+S_{5,4,1}+S_{5,3,2}=0,$$
$$S_{6,2,2}+S_{5,3,2}=0$$
where $`S_{10}`$ and $`S_{4,4,2}`$ are considered as the invariant polynomials and other dependent invariants can be expressed in terms of them as follows:
$$S_{4,3,3}=0,S_{5,3,2}=2S_{4,4,2},$$
$$S_{6,2,2}=2S_{4,4,2},S_{9,1}=S_{10},$$
$$S_{5,4,1}=S_{4,4,2},S_{6,3,1}=3S_{4,4,2},$$
$$S_{7,2,1}=7S_{4,4,2},S_{8,1,1}=7S_{4.4,2},$$
$$S_{8,2}=S_{10}14S_{4,4,2},S_{7,3}=21S_{4,4,2}S_{10},$$
$$S_{6,4}=S_{10}24S_{4,4,2},S_{5,5}=\frac{25S_{4,4,2}S_{10}}{2},$$
In Appendix IV we use the constraints concerned with the invariant polynomials of order up to $`22`$ to express the dependent invariant polynomials in terms of the independent ones.
## 5 Moments of Current Distribution and Multifractal Spectrum
In order to study the multifractals behaviour of current distribution we consider their $`q`$-moments defined as:
$$M_q(n)=\underset{r}{}I_r(n)^q$$
where $`I_r`$ is the current in the $`r`$-th bond of subfractals of generation level n. From the $`S_3`$ symmetry of Sierpinsky fractal, it is clear that the $`q`$-moments depend only on the independent Schure’s $`S_3`$ invariant polynomials of degree q of input currents defined in section IV, that is
$$M_q(n+1)=\underset{partitionscorrespondingtoindependentpolynomials}{}A_{\lambda _1,\lambda _2,\lambda _3}(n+1)S_{\lambda _1,\lambda _2,\lambda _3}(n+1),$$
(5-1)
where $`A_{\lambda _1,\lambda _2,\lambda _3}`$,s are some constants.
On the other hand, $`M_q(n+1)`$ can be written in terms of the invariant polynomials of their level $`n`$ subfractals, that is
$$M_q(n+1)=\underset{partitionscorrespondingtoinvariantpolynomials}{}A_{\lambda _1,\lambda _2,\lambda _3}(n)S_{\lambda _1,\lambda _2,\lambda _3}(n).$$
(5-2)
By comparing the expressions (5-1) and (5-2) we obtain the recursion relations between $`A_{\lambda _1,\lambda _2,\lambda _3}(n)`$ and $`A_{\lambda _1,\lambda _2,\lambda _3}(n+1)`$. Then the scaling factor is defined as:
$$\lambda (q)=\underset{n\mathrm{}}{lim}\frac{M_q(n+1)}{M_q(n)}.$$
(5-3)
Obviously $`\lambda (q)`$ is the maximum eigenvalue of the matrix connecting A(n) and A(n+1). Then $`D(q)`$, the multifractals scaling exponents, are defined as:
$$D(q)=\frac{\mathrm{ln}(\lambda (q))}{\mathrm{ln}(b)},$$
(5-4)
since the $`M_q(n)`$ scale as:
$$Lim_n\mathrm{}M_q(n)=L_n^{D(q)},$$
where $`L_n=b^n`$.
Now, as an example, we obtain $`D_2`$, the power scaling exponent of Sierpinsky fractal with decimation numbers $`b=2,3,4,4`$ and $`5`$.
According to formula (4-4) for $`q=2`$ we have only one independent invariant polynomial, where we can consider $`S_2`$ as the independent invariant polynomial. Therefore the total power is proportional to $`S_2`$, that is:
$$P(n+1)=A_2(n+1)S_2(n+1).$$
(5-5)
It is straightforward to show that:
$$\frac{A_2(n+1)}{A_2(n)}=\frac{5}{3}forb=2$$
$$\frac{A_2(n+1)}{A_2(n)}=\frac{45}{21}forb=3$$
$$\frac{A_2(n+1)}{A_2(n)}=\frac{3399}{1353}forb=4$$
$$\frac{A_2(n+1)}{A_2(n)}=\frac{8576091}{3047787}forb=5$$
Therefore, we have:
$$D(2)=.7369655945forb=2$$
$$D(2)=.6937297714forb=3$$
$$D(2)=.6644742613forb=4$$
$$D(2)=.6428097998forb=5$$
.
In case of $`D(4)`$,we have to consider the $`S_4`$,$`S_{3,1}`$,$`S_{2,2}`$ ,$`S_{2,1,1}`$, and due to relations (4-5) only one of them, say $`S_4`$, is independent and the others can be written in terms of it. Again by computing the fourth moments of currents of fractals which are proportional to $`S_4`$:
$$M_4(n)=A_4(n)S_4(n)M_4(n+1)=A_4(n+1)S_4(n+1)$$
(5-6)
one can easily show that:
$$\frac{A_{4(n+1)}}{A_{4(n)}}=1.222222222forb=2$$
$$\frac{A_{4(n+1)}}{A_{4(n)}}=1.288213244forb=3$$
$$\frac{A_{4(n+1)}}{A_{4(n)}}=1.323683604forb=4$$
$$\frac{A_{4(n+1)}}{A_{4(n)}}=1.231193828forb=5.$$
Hence using the formula (5-4) we get
$$D(4)=.2895066169forb=2$$
$$D(4)=.2305237058forb=3$$
$$D(4)=.2022791602forb=4$$
$$D(4)=.1292279056forb=5$$
With the above explained prescription, we can calculate the higher moments and consequently higher multifractals exponents, where we quote only the multifractals exponents below in the remaining part of this section and give the other information such as the recursion relations in appendix V:
Using the above results, the best fit we can get for various multifractals exponents are:
$$D(q,b=2)=1+4\times 2^q,$$
$$D(q,b=3)=1+51.47353178\times 3^q,$$
$$D(q,b=4)=1+291.7913871\times 4^q,$$
$$D(q,b=5)=1+650.6706017\times 5^q,$$
where the first formula is the same as the formula of reference. The above formulas show the scaling behaviour of the multifractals spectra. Appendix I: Calculation of currents of $`b=4`$.
Here in this Appendix we give the detail of the calculation of inner inward flowing currents corresponding to decimation number b=4. Following the procedure of section III, for b=4 we have:
$$I_{3\delta _{1,j},3\delta _{2,j},3\delta _{3,j}}(4\delta _{1,j},4\delta _{2,j},4\delta _{3,j})=I_j,$$
$$I_{3\delta _{1,j},3\delta _{2,j},3\delta _{3,j}}(3\delta _{1,j}+\delta _{1,k},3\delta _{2,j}+\delta _{2,k},3\delta _{3,j}+\delta _{3,k})=a_3(3)I_j+b_{31}(3)I_j+b_{31}(1)I_k,$$
$$I_{2\delta _{1,j}+\delta _{1,k},2\delta _{2,j}+\delta _{2,k},2\delta _{3,j}+\delta _{3,k}}(3\delta _{1,j}+\delta _{1,k},3\delta _{2,j}+\delta _{2,k},3\delta _{3,j}+\delta _{3,k})=a_{21}(2)I_j+a_{21}(1)I_k+b_{31}(3)I_j+b_{31}(1)I_k,$$
$$I_{2\delta _{1,j}+\delta _{1,k},2\delta _{2,j}+\delta _{2,k},2\delta _{3,j}+\delta _{3,k}}(2\delta _{1,j}+2\delta _{1,k},2\delta _{2,j}+2\delta _{2,k},2\delta _{3,j}+2\delta _{3,k})=a_{21}(2)I_j+a_{21}(1)I_k+b_{22}(2)(I_j+I_k),$$
$$I_{2\delta _{1,j}+\delta _{1,k},2\delta _{2,j}+\delta _{2,k},2\delta _{3,j}+\delta _{3,k}}(2\delta _{1,j}+\delta _{1,k}+\delta _{1,l},2\delta _{2,j}+\delta _{2,k}+\delta _{2,l},2\delta _{3,j}+\delta _{3,k}+\delta _{3,l})$$
$$=a_{21}(2)I_j+a_{21}(1)+b_{211}(2)I_j+b_{211}(1)(I_k+I_l),$$
$$I_{\delta _{1,j}+\delta _{1,k}+\delta _{1,l},\delta _{2,j}+\delta _{2,k}+\delta _{2,l},\delta _{3,j}+\delta _{3,k}+\delta _{3,l}}(2\delta _{1,j}+\delta _{1,k}+\delta _{1,l},2\delta _{2,j}+\delta _{2,k}+\delta _{2,l},2\delta _{3,j}+\delta _{3,k}+\delta _{3,l})$$
$$=47a_{111}(1)(I_j+I_k+I_l)+b_{211}(2)I_j+b_{211}(1)(I_k+I_l).$$
Now, imposing Kirchhoff’s law on subfractals and on vertices, we get the following equations for $`𝐚`$ and $`𝐛`$
$`1+2a_3(3)+b_{31}(3)b_{31}(1)=0,`$
$`3a_{21}(2)+b_{31}(3)+b_{22}(2)+2b_{211}(2)b_{211}(1)=0,`$
$`3a_{21}(2)+b_{31}(1)+b_{22}(2)+b_{211}(1)=0,`$
$`3a_{111}(2)+b_{211}(2)=0`$
$`a_{21}(1)+2b_{31}(1)=0,`$
$`a_3(3)+b_{21}(2)+2b_{31}(3)=0,`$
$`a_{21}(2)+a_{21}(1)+2b_{22}(2)=0,`$
$`2a_{21}(2)+a_{111}(1)+3b_{211}(2)=0,`$
$`a_{21}(1)+a_{111}(1)+3b_{211}(1)=0,`$
$`3a_{111}(1)+4b_{1111}(1)=0.`$
By solving the above equations we can determine inner inward flowing currents corresponding to decimation number b=4 which is given in subsection 3.1.
Appendix II: Calculation of currents of b=5.
Here in this Appendix we give the detail of the calculation of inner inward flowing currents corresponding to decimation number b=5. Following the procedure of section III, for b=5 we have:
$$I_{4\delta _{1,j},4\delta _{2,j},4\delta _{3,j}}(5\delta _{1,j},5\delta _{2,j},5\delta _{3,j})=I_j,$$
$$I_{4\delta _{1,j},4\delta _{2,j},4\delta _{3,j}}(4\delta _{1,j}+\delta _{1,k},4\delta _{2,j}+\delta _{2,k},4\delta _{3,j}+\delta _{3,k})=a_4(4)I_j+b_{41}(4)I_j+b_{41}(1)I_k,$$
$$I_{3\delta _{1,j}+\delta _{1,k},3\delta _{2,j}+\delta _{2,k},3\delta _{3,j}+\delta _{3,k}}(4\delta _{1,j}+\delta _{1,k},4\delta _{2,j}+\delta _{2,k},4\delta _{3,j}+\delta _{3,k})=a_{31}(3)I_j+a_{31}(1)I_k+b_{41}(4)I_j+b_{41}(1)I_k,$$
$$I_{3\delta _{1,j}+\delta _{1,k},3\delta _{2,j}+\delta _{2,k},3\delta _{3,j}+\delta _{3,k}}(3\delta _{1,j}+2\delta _{1,k},3\delta _{2,j}+2\delta _{2,k},3\delta _{3,j}+2\delta _{3,k})$$
$$=a_{31}(3)I_j+a_{31}(1)I_k+b_{32}(3)I_j+b_{32}(3)I_k,$$
$$I_{3\delta _{1,j}+\delta _{1,k},3\delta _{2,j}+\delta _{2,k},3\delta _{3,j}+\delta _{3,k}}(3\delta _{1,j}+\delta _{1,k}+\delta _{1,l},3\delta _{2,j}+\delta _{2,k}+\delta _{2,l},3\delta _{3,j}+\delta _{3,k}+\delta _{3,l})$$
$$=a_{31}(3)I_j+a_{31}(1)+b_{311}(3)I_j+b_{311}(1)(I_k+I_l),$$
$$I_{2\delta _{1,j}+2\delta _{1,k},2\delta _{2,j}+2\delta _{2,k},2\delta _{3,j}+2\delta _{3,k}}(3\delta _{1,j}+2\delta _{1,k},3\delta _{2,j}+2\delta _{2,k},3\delta _{3,j}+2\delta _{3,k})=a_{22}(2)(I_j+I_k)+b_{32}(3)I_j+b_{32}(2)I_k.$$
Again imposing Kirchhoff’s law on subfractals and on vertices, we get the following equations for $`𝐚`$ and $`𝐛`$:
$`1+2a_4(4)+b_{41}(3)b_{41}(1)=0,`$
$`3a_{31}(3)+b_{41}(4)+b_{32}(3)+b_{311}(3)b_{311}(1)=0,`$
$`3a_{31}(1)+b_{41}(1)+b_{32}(2)=0,`$
$`3a_{22}(1)+b_{32}(3)+b_{32}(2)+b_{221}(2)b_{221}(1)=0,`$
$`3a_{211}(2)+b_{311}(3)+2b_{221}(2)b_{2111}(1)=0,`$
$`3a_{211}(2)+b_{311}(1)+b_{221}(2)+b_{221}(1)=0,`$
$`3a_{1111}(2)+b_{2111}(2)+3b_{2111}(1)=0,`$
$`a_4(4)+a_{31}(3)+2b_{41}(4)=0,`$
$`a_{31}(1)+2b_{41}(1)=0,`$
$`a_{31}(3)+a_{22}(2)+2b_{32}(3)=0,`$
$`a_{31}(1)+a_{22}(2)+2b_{32}(2)=0,`$
$`2a_{31}(3)+a_{211}(2)+3b_{311}(3)=0,`$
$`a_{31}(1)+a_{211}(1)+3b_{311}(1)=0,`$
$`a_{22}(2)+a_{211}(2)+a_{211}(1)+3b_{221}(2)=0,`$
$`2a_{211}(1)+3b_{221}(1)=0,`$
$`3a_{211}(2)+a_{1111}(1)+4b_{2111}(2)=0,`$
$`2a_{211}(1)+a_{1111}(1)+4b_{2111}(1)=0,`$
$`4a_{1111}(1)+5b_{11111}(1)=0.`$
By solving the above equations we can determine inner inward flowing currents corresponding to decimation number b=5 which is given in subsection 3.2.
Appendix III: Proof of the formula (4-6):
Here we give the proof of the formula (4-6). The number of independent Shure’s invariant polynomials of degree $`2k`$ of 3 variables $`I_1`$,$`I_2`$,$`I_3`$ with the constraint:
$$I_1+I_2+I_3=0,$$
(III-1)
is equal to
$$P_3(2K)P_3(2k1),$$
(III-2)
where $`P_3(m)`$ is the number of partition of $`m`$ into at most three independent non-negative integers. If we define $`M_k(n)`$,the number of partitions of $`n`$ into exactly $`k`$ non-negative integers, then we have
$$P_3(n)=\underset{k=1}{\overset{3}{}}M_k(n).$$
(III-3)
Obviously
$$M_1(2k)=M_1(2k1).$$
(III-4)
and
$$M_2(2k)=M_2(2k1)+1,M_2(2k)=M_2(2k+1).$$
(III-5)
If we denote the partition of $`2k`$, $`2k1`$ and $`2k+1`$ into two non-negative integers respectively by: ($`l_1,l_2`$), ($`m_1,m_2`$) and ($`n_1,n_2`$) then in the case of $`l_2=m_2=n_2`$ we will have
$$l_1=m_11=n_1+1.$$
(III-6)
Therefore, for all values of $`l_1>k`$, there is a one to one correspondence between the $`M_2(2k)`$,$`M_2(2k1)`$ and $`M_2(2k+1)`$. Only for $`l_1=k`$, $`n_1`$ can be equal to $`k+1`$, but $`m_1`$ cannot be equal to $`k1`$, thereof the relations (III-4) and (III-5) follows. Now, we are ready to prove that
$$M_3(2k)=M_3(2k1)+[k/3].$$
(III-7)
If we denote the partition of $`2k`$ and $`2k1`$ into three non-negative integers by $`(l_1,l_2,l_3)`$ and $`(m_1,m_2,m_3)`$ respectively, then using the relations (III-4) and (III-5), we can prove $`M_3(2k)=M_3(2k1)`$for $`m_1=l_1=odd`$; and for $`l_1=m_1=even`$, we would have $`M_3(2k)=M_3(2k1)`$. Since $`l_1`$ takes values between $`1`$ and $`[2k/3]`$, where $`[[2k/3]/2]=[k/3]`$ of them correspond to even values of $`l_1`$, the relation (III-7) follows and the proof is complete. APPENDIX IV:Solution of Constraints Over Schure’s invariants polynomials:
Here in this appendix by solving the constraints over Schure’s polynomials of degree $`12`$, $`14`$,$`16`$,$`18`$,$`20`$ and $`22`$, we have expressed the dependent invariant polynomials in terms of the independent invariant polynomials.
1) Solution of Constraints of degree $`12`$:
The invariant polynomials $`S_{12}`$, $`S_{8,2,2}`$ and $`S_{6,3,3}`$ are considered to be independent and the other dependent invariant polynomials can be written in terms of them as follows:
$$S_{5,4,3}=S_{6,3,3},S_{4,4,4}=\frac{S_{6,3,3}}{3}$$
$$S_{7,3,2}=S_{8,2,2},S_{6,4,2}=S_{8,2,2}2S_{6,3,3},$$
$$S_{5,5,2}=\frac{3S_{6,3,3}S_{8,2,2}}{2},S_{6,5,1}=\frac{S_{8,2,2}3S_{5,3,3}}{2}$$
$$S_{7,4,1}=\frac{7S_{6,3,3}3S_{8,2,2}}{2},S_{8,3,1}=\frac{5S_{8,2,2}7S_{6,3,3}}{2},$$
$$S_{9,2,1}=\frac{7S_{6,3,3}9S_{8,2,2}}{2},S_{10,1,1}=\frac{9S_{8,2,2}7S_{6,3,3}}{2},$$
$$S_{11,1}=S_{12},$$
$$S_{10,2}=S_{12}+7S_{6,3,3}9S_{8,2,2},$$
$$S_{9,3}=\frac{27S_{8,2,2}21S_{6,3,3}2S_{12}}{2},$$
$$S_{8,4}=\frac{28S_{6,3,3}32S_{8,2,2}+2S_{12}}{2},$$
$$S_{7,5}=\frac{35S_{8,2,2}2S_{12}35S_{6,3,3}}{2},$$
$$S_{6,6}=\frac{2S_{12}+38S_{6,3,3}36S_{8,2,2}}{4}.$$
2) Solution of Constraints of degree 14:
The invariant polynomials $`S_{14}`$, $`S_{6,6,2}`$ and $`S_{5,5,4}`$ are considered to be independent one and the other dependent invariant polynomials can be written in terms of them as follows:
$$S_{10,3,1}=7S_{6,2,2}23S_{5,5,4},$$
$$S_{8,5,1}=3S_{6,6,2}S_{5,5,4},$$
$$S_{9,4,1}=5S_{6,6,2}+7S_{5,5,4},$$
$$S_{8,3,3}=5S_{5,5,4},S_{7,4,3}=5S_{5,5,4},$$
$$S_{6,5,3}=S_{5,5,4},S_{6,4,4}=2S_{5,5,4},$$
$$S_{9,5}=S_{14}+45S_{6,6,2}195S_{5,5,4},$$
$$S_{10,4}=S_{14}40S_{6,6,2}+188S_{5,5,4},$$
$$S_{11,3}=S_{14}+33S_{6,6,2}165S_{5,5,4},$$
$$S_{7,7}=\frac{S_{14}+49S_{6,6,2}196S_{5,5,4}}{2},$$
$$S_{8,6}=S_{14}48S_{6,6,2}+196S_{5,5,4},$$
$$S_{13,1}=S_{14},$$
$$S_{8,4,2}=2S_{6,62}6S_{5,5,4},$$
$$S_{12,2}=S_{14}22S_{6,6,2}+110S_{5,5,4},$$
$$S_{9,3,2}=2S_{6,6,2}+16S_{5,5,4},$$
$$S_{7,5,2}=2S_{6,6,2}+S_{5,5,4},$$
$$S_{10,2,2}=2S_{6,6,2}16S_{5,5,4},$$
$$S_{7,6,1}=S_{6,6,2},$$
$$S_{11,2,1}=11S_{6,6,2}+55S_{5,5,4},$$
$$S_{12,1,1}=11S_{6,6,2}55S_{5,5,4}.$$
3) Solution of Constraints of degree 16: The invariant polynomials $`S_{16}`$, $`S_{7,7,2}`$, $`S_{6,6,4}`$ are considered to be independent and the other dependent invariant polynomials can be written in terms of them as follows:
$$S_{6,5,5}=0$$
$$S_{7,5,4}=2S_{6,6,4},S_{7,6,3}=S_{6,6,4},$$
$$S_{8,8,4}=2S_{6,6,4},S_{8,5,3}=3S_{6,6,4},$$
$$S_{8,7,1}=S_{7,7,2},S_{8,6,2}=2S_{7,7,2}+S_{6,6,4},$$
$$S_{9,4,3}=7S_{6,6,4},S_{9,5,2}=2S_{7,7,2}4S_{6,6,4},$$
$$S_{9,6,1}=3S_{7,7,2}S_{6,6,4},S_{10,3,3}=7S_{6,6,4},$$
$$S_{10,4,2}=2S_{7,7,2}+11S_{6,6,4},$$
$$S_{10,5,1}=5S_{7,7,2}+5S_{6,6,4},$$
$$S_{11,3,2}=2S_{7,7,2}25S_{6,6,4},$$
$$S_{11,4,1}=7S_{7,7,2}16S_{6,6,4},$$
$$S_{12,2,2}=2S_{7,7,2}+25S_{6,6,4},$$
$$S_{12,3,1}=9S_{7,7,2}+41S_{6,6,4},$$
$$S_{13,2,1}=13S_{7,7,2}91S_{6,6,4},$$
$$S_{14,1,1}=13S_{7,7,2}+91S_{6,6,4},$$
$$S_{15,1}=S_{16},$$
$$S_{14,2}=S_{16}+26S_{7,7,2}182S_{6,6,4},$$
$$S_{13,3}=S_{16}39S_{7,7,2}+273S_{6,6,4},$$
$$S_{12,4}=S_{16}+48S_{7,7,2}314S_{6,6,4},$$
$$S_{11,5}=S_{16}55S_{7,7,2}+330S_{6,6,4},$$
$$S_{10,6}=S_{16}+60S_{7,7,2}335S_{6,6,4},$$
$$S_{9,7}=S_{16}63S_{7,7,2}+336S_{6,6,4},$$
$$S_{8,8}=\frac{S_{16}+64S_{7,7,2}336S_{6,6,4}}{2}.$$
4) Solution of Constraints of degree 18:
The invariant polynomials $`S_{18}`$, $`S_{14,2,2}`$, $`S_{7,7,4}`$, $`S_{8,5,5}`$ are considered to be independent and the other dependent invariant polynomials can be written in terms of them as follows:
$$S_{7,6,5}=S_{8,5,5},S_{6,6,6}=\frac{S_{8,5,5}}{3},$$
$$S_{8,7,3}=S_{7,7,4},S_{8,6,4}=S_{8,5,5}2S_{7,7,4},$$
$$S_{9,5,4}=3S_{8,5,5}+2S_{7,7,4},S_{9,6,3}=3S_{7,7,4}S_{8,5,5},$$
$$S_{10,5,3}=4S_{8,5,5}5S_{7,7,4},S_{10,4,4}=3S_{8,5,5}2S_{7,7,4},$$
$$S_{11,4,3}=9S_{7,7,4}10S_{8,5,5},S_{12,3,3}=10S_{8,5,5}9S_{7,7,4},$$
$$S_{13,3,2}=S_{14,2,2},$$
$$S_{12,4,2}=S_{14,2,2}20S_{8,5,5}+18S_{7,7,4},$$
$$S_{11,5,2}=S_{14,2,2}+30S_{8,5,5}27S_{7,7,4},$$
$$S_{10,6,2}=32S_{7,7,4}34S_{8,5,5}+S_{14,2,2},$$
$$S_{9,7,2}=35S_{8,5,5}35S_{7,7,4}S_{14,2,2},$$
$$S_{8,8,2}=\frac{36S_{7,7,4}+S_{14,2,2}35S_{8,5,5}}{2},$$
$$S_{9,8,1}=\frac{35S_{8,5,5}S_{14,2,2}36S_{7,7,4}}{2},$$
$$S_{10,7,1}=\frac{106S_{7,7,4}105S_{8,5,5}+3S_{14,2,2}}{2},$$
$$S_{11,6,1}=\frac{173S_{8,5,5}170S_{7,7,4}5S_{14,2,2}}{2},$$
$$S_{12,5,1}=\frac{224S_{7,7,4}233S_{8,5,5}+7S_{14,2,2}}{2},$$
$$S_{13,4,1}=\frac{273S_{8,5,5}260S_{7,7,4}9S_{14,2,2}}{2},$$
$$S_{14,3,1}=260S_{7,7,4}273S_{8,5,5}+11S_{14,2,2},$$
$$S_{15,2,1}=\frac{273S_{8,5,5}260S_{7,7,4}15S_{14,2,2}}{2},$$
$$S_{16,1,1}=\frac{260S_{7,7,4}+15S_{14,2,2}273S_{8,5,5}}{2},$$
$$S_{17,1}=S_{18},$$
$$S_{16,2}=S_{18}260S_{7,7,4}15S_{14,2,2}+273S_{8,5,5},$$
$$S_{15,3}=\frac{780S_{7,7,4}2S_{18}+45S_{14,2,2}819S_{8,5,5}}{2},$$
$$S_{14,4}=\frac{1092S_{8,5,5}56S_{14,2,2}+2S_{18}1040S_{7,7,4}}{2},$$
$$S_{13,5}=\frac{65S_{14,2,2}2S_{18}1365S_{8,5,5}+1300S_{7,7,4}}{2},$$
$$S_{12,6}=\frac{2S_{18}+1598S_{8,5,5}1524S_{7,7,4}72S_{14,2,2}}{2},$$
$$S_{11,7}=\frac{1694S_{7,7,4}+77S_{14,2,2}2S_{18}1771S_{8,5,5}}{2},$$
$$S_{10,8}=\frac{2S_{18}+1876S_{8,5,5}80S_{14,2,2}1800S_{7,7,4}}{2},$$
$$S_{9,9}=\frac{81S_{14,2,2}+1836S_{7,7,4}2S_{18}1911S_{8,5,5}}{4}.$$
5) Solution of Constraints of degree 20:
The invariant polynomials $`S_{20}`$, $`S_{16,2,2}`$, $`S_{8,8,4}`$ and $`S_{7,7,6}`$ are considered to be independent and the other dependent invariant polynomials can be written in terms of them as follows:
$$S_{8,8,6}=2S_{7,7,6},S_{8,7,5}=S_{7,7,6},$$
$$S_{9,6,5}=5S_{7,7,6},S_{9,8,3}=S_{8,8,4},$$
$$S_{9,7,4}=S_{7,7,6}2S_{8,8,4},S_{10,5,5}=5S_{7,7,6},$$
$$S_{10,6,4}=6S_{7,7,6}+2S_{8,8,4},S_{10,7,3}=3S_{8,8,4}S_{7,7,6},$$
$$S_{11,6,3}=5S_{8,8,4}+7S_{7,7,6},S_{11,5,4}=+16S_{7,7,6}2S_{8,8,4},$$
$$S_{12,4,4}=(16S_{7,7,6}+2S_{8,8,4}),S_{12,5,3}=23S_{7,7,6}+7S_{8,8,4},$$
$$S_{13,4,3}=55S_{7,7,6}11S_{8,8,4},S_{14,3,3}=55S_{7,7,6}+11S_{8,8,4},$$
$$S_{15,3,2}=S_{16,2,2},$$
$$S_{14,4,2}=22S_{8,8,4}+S_{16,2,2}+110S_{7,7,6},$$
$$S_{13,5,2}=165S_{7,7,6}S_{16,2,2}+33S_{8,8,4},$$
$$S_{12,6,2}=40S_{8,8,4}+188S_{7,7,6}+S_{16,2,2},$$
$$S_{11,7,2}=195S_{7,7,6}+45S_{8,8,4}S_{16,2,2},$$
$$S_{10,8,2}=48S_{8,8,4}+196S_{7,7,6}+S_{16,2,2},$$
$$S_{9,9,2}=\frac{196S_{7,7,6}+49S_{8,8,4}S_{16,2,2}}{2},$$
$$S_{10,9,1}=\frac{S_{16,2,2}+196S_{7,7,6}49S_{8,8,4}}{2},$$
$$S_{11,8,1}=\frac{588S_{7,7,6}+145S_{8,8,4}3S_{16,2,2}}{2},$$
$$S_{12,7,1}=\frac{+978S_{7,7,6}235S_{8,8,4}+5S_{16,2,2}}{2},$$
$$S_{13,6,1}=\frac{1354S_{7,7,6}+315S_{8,8,4}7S_{16,2,2}}{2},$$
$$S_{14,5,1}=\frac{9S_{16,2,2}381S_{8,8,4}+1684S_{7,7,6}}{2},$$
$$S_{15,4,1}=\frac{425S_{8,8,4}1904S_{7,7,6}11S_{16,2,2}}{2},$$
$$S_{16,3,1}=\frac{13S_{16,2,2}+1904S_{7,7,6}425S_{8,8,4}}{2},$$
$$S_{17,2,1}=\frac{1904S_{7,7,6}+425S_{8,8,4}17S_{16,2,2}}{2},$$
$$S_{18,1,1}=\frac{17S_{16,2,2}425S_{8,8,4}+1904S_{7,7,6}}{2},$$
$$S_{19,1}=S_{20},$$
$$S_{18,2}=S_{20}+425S_{8,8,4}1904S_{7,7,6}17S_{16,2,2},$$
$$S_{17,3}=\frac{51S_{16,2,2}1275S_{8,8,4}+5712S_{7,7,6}2S_{20}}{2},$$
$$S_{16,4}=\frac{2S_{20}64S_{16,2,2}+1700S_{8,8,4}7616S_{7,7,6}}{2},$$
$$S_{15,5}=\frac{75S_{16,2,2}2125S_{8,8,4}2S_{20}+9520S_{7,7,6}}{2},$$
$$S_{14,6}=\frac{2S_{20}+2506S_{8,8,4}11204S_{7,7,6}84S_{16,2,2}}{2},$$
$$S_{13,7}=\frac{91S_{16,2,2}+12558S_{7,7,6}2821S_{8,8,4}2S_{20}}{2},$$
$$S_{12,8}=\frac{2S_{20}96S_{16,2,2}+3056S_{8,8,4}13536S_{7,7,6}}{2},$$
$$S_{11,9}=\frac{99S_{16,2,2}2S_{20}3201S_{8,8,4}+14124S_{7,7,6}}{2},$$
$$S_{10,10}=\frac{2S_{20}100S_{16,2,2}+3250S_{8,8,4}14320S_{7,7,6}}{2}.$$
6) Solution of Constraints of degree 22:
The invariant polynomials $`S_{22}`$, $`S_{18,2,2}`$, $`S_{9,9,4}`$ and $`S_{8,8,6}`$ are considered to be independent and the other dependent invariant polynomials can be written in terms of them as follows:
$$S_{8,7,7}=0,S_{9,7,6}=2S_{8,8,6},$$
$$S_{9,8,5}=S_{8,8,6},S_{10,6,6}=2S_{8,8,6},$$
$$S_{10,7,5}=3S_{8,8,6},S_{10,9,3}=S_{9,9,4},$$
$$S_{10,8,4}=S_{8,8,6}2S_{9,9,4},S_{11,6,5}=7S_{8,8,6},$$
$$S_{11,7,4}=2S_{9,9,4}4S_{8,8,6},$$
$$S_{11,8,3}=3S_{9,9,4}S_{8,8,6},S_{12,5,5}=7S_{8,8,6},$$
$$S_{12,6,4}=11S_{8,8,6}+2S_{9,9,4},S_{12,7,3}=5S_{8,8,6}5S_{9,9,4},$$
$$S_{13,5,4}=2S_{9,9,4}25S_{8,8,6},S_{13,6,3}=7S_{9,9,4}16S_{8,8,6},$$
$$S_{14,4,4}=25S_{8,8,6}2S_{9,9,4},S_{14,5,3}=41S_{8,8,6}9S_{9,9,4},$$
$$S_{15,4,3}=13S_{9,9,4}91S_{8,8,6},S_{16,3,3}=91S_{8,8,6}13S_{9,9,4},$$
$$S_{17,3,2}=S_{18,2,2},$$
$$S_{16,4,2}=26S_{9,9,4}182S_{8,8,6}+S_{18,2,2},$$
$$S_{15,5,2}=273S_{8,8,6}39S_{9,9,4}S_{18,2,2},$$
$$S_{14,6,2}=48S_{9,9,4}314S_{8,8,6}+S_{18,2,2},$$
$$S_{13,7,2}=330S_{8,8,6}55S_{9,9,4}S_{18,2,2},$$
$$S_{12,8,2}=60S_{9,9,4}+S_{18,2,2}335S_{8,8,6},$$
$$S_{11,9,2}=336S_{8,8,6}63S_{9,9,4}S_{18,2,2},$$
$$S_{10,10,2}=\frac{S_{18,2,2}+64S_{9,9,4}336S_{8,8,6}}{2},$$
$$S_{11,10,1}=\frac{64S_{9,9,4}+336S_{8,8,6}S_{18,2,2}}{2},$$
$$S_{12,9,1}=\frac{190S_{9,9,4}+3S_{18,2,2}1008S_{8,8,6}}{2},$$
$$S_{13,8,1}=\frac{310S_{9,9,4}5S_{18,2,2}+1678S_{8,8,6}}{2},$$
$$S_{14,7,1}=\frac{2338S_{8,8,6}+7S_{18,2,2}+420S_{9,9,4}}{2},$$
$$S_{15,6,1}=\frac{2966S_{8,8,6}9S_{18,2,2}516S_{9,9,4}}{2},$$
$$S_{16,5,1}=\frac{11S_{18,2,2}+594S_{9,9,4}3512S_{8,8,6}}{2},$$
$$S_{17,4,1}=\frac{3876S_{8,8,6}13S_{18,2,2}646S_{9,9,4}}{2},$$
$$S_{18,3,1}=\frac{3876S_{8,8,6}+646S_{9,9,4}+15S_{18,2,2}}{2},$$
$$S_{19,2,1}=\frac{646S_{9,9,4}19S_{18,2,2}+3876S_{8,8,6}}{2},$$
$$S_{20,1,1}=\frac{3876S_{8,8,6}+19S_{18,2,2}+646S_{9,9,4}}{2},$$
$$S_{21,1}=S_{22},$$
$$S_{20,2}=S_{22}19S_{18,2,2}646S_{9,9,4}+3876S_{8,8,6},$$
$$S_{19,3}=\frac{57S_{18,2,2}2S_{22}+1938S_{9,9,4}11628S_{8,8,6}}{2},$$
$$S_{18,4}=\frac{2S_{22}72S_{18,2,2}2584S_{9,9,4}+15504S_{8,8,6}}{2},$$
$$S_{17,5}=\frac{3230S_{9,9,4}+85S_{18,2,2}19380S_{8,8,6}2S_{22}}{2},$$
$$S_{16,6}=\frac{2S_{22}96S_{18,2,2}3824S_{9,9,4}+22892S_{8,8,6}}{2},$$
$$S_{15,7}=\frac{105S_{18,2,2}+4340S_{9,9,4}25858S_{8,8,6}2S_{22}}{2},$$
$$S_{14,8}=\frac{28196S_{8,8,6}112S_{18,2,2}+2S_{22}4760S_{9,9,4}}{2},$$
$$S_{13,9}=\frac{5070S_{9,9,4}+117S_{18,2,2}29874S_{8,8,6}2S_{22}}{2},$$
$$S_{12,10}=\frac{30882S_{8,8,6}+2S_{22}120S_{18,2,2}5260S_{9,9,4}}{2},$$
$$S_{11,11}=\frac{121S_{18,2,2}+4678S_{9,9,4}2S_{22}31218S_{8,8,6}}{4}.$$
APPENDIX IV:the recursion relation and $`\lambda _{max}`$ for $`q6`$ a)b=2 q=6
$$\left(\begin{array}{c}A_{6(n+1)}\\ A_{4,1,1(n+1)}\end{array}\right)=\left(\begin{array}{cc}1.090534979& .1646090535\\ .1481481481& .1851851852\end{array}\right)\left(\begin{array}{c}A_{6(n)}\\ A_{4,1,1(n)}\end{array}\right),$$
$$\lambda _{max}=1.062745991.$$
q=8
$$\left(\begin{array}{c}A_{8(n+1)}\\ A_{3,3,2(n+1)}\end{array}\right)=\left(\begin{array}{cc}1.039323274& .5121170553\\ .03840877915& .1742112483\end{array}\right)\left(\begin{array}{c}A_{8(n)}\\ A_{3,3,2(n)}\end{array}\right),$$
$$\lambda _{max}=1.015955376.$$
q=10
$$\left(\begin{array}{c}A_{10(n+1)}\\ A_{4,4,2(n+1)}\end{array}\right)=\left(\begin{array}{cc}1.017375400& .3840877915\\ .02987349489& .1486053955\end{array}\right)\left(\begin{array}{c}A_{10(n)}\\ A_{4,4,2(n)}\end{array}\right),$$
$$\lambda _{max}=1.003961045.$$
q=12
$$\left(\begin{array}{c}A_{12(n+1)}\\ A_{8,2,2(n+1)}\\ A_{6,3,3(n+1)}\end{array}\right)=\left(\begin{array}{ccc}1.007711110& .1266744568& .1217068310\\ .05378583888& .1908753747& .1417805551\\ .007315957933& .07681755830& .04709647920\end{array}\right)\left(\begin{array}{c}A_{12(n)}\\ A_{8,2,2(n)}\\ A_{6,3,3(n+1)}\end{array}\right),$$
$$\lambda _{max}=1.000983621.$$
q=14
$$\left(\begin{array}{c}A_{14(n+1)}\\ A_{6,6,2(n+1)}\\ A_{5,5,4(n+1)}\end{array}\right)=\left(\begin{array}{ccc}1.003425906& .1557074696& .6579896295\\ .01048470103& .09541040304& .5668362057\\ .001896729835& .04335382479& .2332808346\end{array}\right)\left(\begin{array}{c}A_{14(n)}\\ A_{6,6,2(n)}\\ A_{5,5,4(n+1)}\end{array}\right),$$
$$\lambda _{max}=1.000244974.$$
q=16
$$\left(\begin{array}{c}A_{16(n+1)}\\ A_{7,7,2(n+1)}\\ A_{6,6,4(n+1)}\end{array}\right)=\left(\begin{array}{ccc}1.001522485& .09132402907& .5053262942\\ .006679533152& .1334905160& 1.056414123\\ .001475234316& .03567658498& .2661499583\end{array}\right)\left(\begin{array}{c}A_{16(n)}\\ A_{7,7,2(n)}\\ A_{6,6,4(n+1)}\end{array}\right),$$
$$\lambda _{max}=1.000061131.$$
q=18
$$\left(\begin{array}{c}A_{18\left(n+1\right)}\\ A_{14,2,2\left(n+1\right)}\\ A_{7,7,4\left(n+1\right)}\\ A_{8,5,5\left(n+1\right)}\end{array}\right)=\left(\begin{array}{cccc}1.000676645& .02587988061& .5694995445& .5946011699\\ .04976341868& .3702685508& 6.165437409& 6.474053436\\ .0007861225946& .004847198466& .1026962541& .1087868226\\ .0003612818732& .009483649173& .1320936849& .1378233896\\ 1.000015270& .1122123554& .01717707276& .001020854222\end{array}\right)\left(\begin{array}{c}A_{18\left(n\right)}\\ A_{14,2,2\left(n\right)}\\ A_{7,7,4\left(n\right)}\\ A_{8,5,5\left(n\right)}\end{array}\right),$$
$$\lambda _{max}=1.000015270.$$
q=20
$$\left(\begin{array}{c}A_{20\left(n+1\right)}\\ A_{16,2,2\left(n+1\right)}\\ A_{8,8,4\left(n+1\right)}\\ A_{7,7,6\left(n+1\right)}\end{array}\right)=\left(\begin{array}{cccc}1.000300729& .01428439338& .4535545013& 2.007507524\\ .04955188625& .4249966010& 10.27981253& 46.05861448\\ .0005177630138& .0009212316079& .05329922806& .2455943625\\ .00009366567084& .002585841556& .05473587640& .2437584270\end{array}\right)\left(\begin{array}{c}A_{20\left(n\right)}\\ A_{16,2,2\left(n\right)}\\ A_{8,8,4\left(n\right)}\\ A_{7,7,6\left(n\right)}\end{array}\right),$$
$$\lambda _{max}=1.000003815.$$
q=22
$$\left(\begin{array}{c}A_{22\left(n+1\right)}\\ A_{18,2,2\left(n+1\right)}\\ A_{9,9,4\left(n+1\right)}\\ A_{8,8,6\left(n+1\right)}\end{array}\right)=\left(\begin{array}{cccc}1.000133657& .007718672824& .3331935728& 1.963553834\\ .04945789936& .4775535736& 15.78748967& 94.73506749\\ .0003298534890& .001784479705& .02202297308& .1217001716\\ .00007285107732& .002090751581& .06241497465& .3728127167\end{array}\right)\left(\begin{array}{c}A_{22\left(n\right)}\\ A_{18,2,2\left(n\right)}\\ A_{9,9,4\left(n\right)}\\ A_{8,8,6\left(n\right)}\end{array}\right),$$
$$\lambda _{max}=1.000000954.$$
b)b=3 q=6
$$\left(\begin{array}{c}A_{6(n+1)}\\ A_{4,1,1(n+1)}\end{array}\right)=\left(\begin{array}{cc}1.082584637& .07343878826\\ .1808430161& .04755671531\end{array}\right)\left(\begin{array}{c}A_{6(n)}\\ A_{4,1,1(n)}\end{array}\right),$$
$$\lambda _{max}=1.069590059.$$
q=8
$$\left(\begin{array}{c}A_{8(n+1)}\\ A_{3,3,2(n+1)}\end{array}\right)=\left(\begin{array}{cc}1.025056713& .1818484281\\ .04533096632& .03647844913\end{array}\right)\left(\begin{array}{c}A_{8(n)}\\ A_{3,3,2(n)}\end{array}\right),$$
$$\lambda _{max}=1.016646559.$$
q=10
$$\left(\begin{array}{c}A_{10(n+1)}\\ A_{4,4,2(n+1)}\end{array}\right)=\left(\begin{array}{cc}1.007845912& .1085523780\\ .03420274178& .02815706873\end{array}\right)\left(\begin{array}{c}A_{10(n)}\\ A_{4,4,2(n)}\end{array}\right),$$
$$\lambda _{max}=1.004041374.$$
q=12
$$\left(\begin{array}{c}A_{12(n+1)}\\ A_{8,2,2(n+1)}\\ A_{6,3,3(n+1)}\end{array}\right)=\left(\begin{array}{ccc}1.002501314& .02774717658& .02388825224\\ .06022405661& .05660671065& .04390015579\\ .01079114607& .01664662055& .01252511118\end{array}\right)\left(\begin{array}{c}A_{12(n)}\\ A_{8,2,2(n)}\\ A_{6,3,3(n+1)}\end{array}\right),$$
$$\lambda _{max}=1.001026135.$$
q=14
$$\left(\begin{array}{c}A_{14(n+1)}\\ A_{6,6,2(n+1)}\\ A_{5,5,4(n+1)}\end{array}\right)=\left(\begin{array}{ccc}1.000805713& .02587192638& .1198110012\\ .008633787717& .02491545222& .1279091010\\ .002716116832& .009331353906& .04723947173\end{array}\right)\left(\begin{array}{c}A_{14(n)}\\ A_{6,6,2(n)}\\ A_{5,5,4(n+1)}\end{array}\right),$$
$$\lambda _{max}=1.000246359.$$
q=16
$$\left(\begin{array}{c}A_{16(n+1)}\\ A_{7,7,2(n+1)}\\ A_{6,6,4(n+1)}\end{array}\right)=\left(\begin{array}{ccc}1.000261069& .01135511506& .07104143816\\ .004468408282& .03190697902& .2276301973\\ .002050937173& .007546589748& .05339121346\end{array}\right)\left(\begin{array}{c}A_{16(n)}\\ A_{7,7,2(n)}\\ A_{6,6,4(n+1)}\end{array}\right),$$
$$\lambda _{max}=1.000061317.$$
q=18
$$\left(\begin{array}{c}A_{18\left(n+1\right)}\\ A_{14,2,2\left(n+1\right)}\\ A_{7,7,4\left(n+1\right)}\\ A_{8,5,5\left(n+1\right)}\end{array}\right)=\left(\begin{array}{cccc}1.000084877& .002388883753& .04740112573& .04969629320\\ .06005002155& .07399773075& 1.274170117& 1.337882220\\ .0009014749149& .0004616094519& .008788472030& .009233751643\\ .0006471907605& .002494396769& .04199894504& .04409378714\end{array}\right)\left(\begin{array}{c}A_{18\left(n\right)}\\ A_{1422\left(n\right)}\\ A_{7,7,4\left(n\right)}\\ A_{8,5,5\left(n\right)}\end{array}\right),$$
$$\lambda _{max}=1.000015297.$$
q=20
$$\left(\begin{array}{c}A_{20\left(n+1\right)}\\ A_{16,2,2\left(n+1\right)}\\ A_{8,8,4\left(n+1\right)}\\ A_{7,7,6\left(n+1\right)}\end{array}\right)=\left(\begin{array}{cccc}1.000027647& .0009744218073& .02806977909& .1253211808\\ .05999917318& .08477443355& 2.107893057& 9.443400292\\ .0005178070963& .0004282216637& .009448647802& .04228755067\\ .0001628983371& .0006696634259& .01635654929& .07326845888\end{array}\right)\left(\begin{array}{c}A_{20\left(n\right)}\\ A_{16,2,2\left(n\right)}\\ A_{8,8,4\left(n\right)}\\ A_{7,7,6\left(n\right)}\end{array}\right),$$
$$\lambda _{max}=1.000003819.$$
q=22
$$\left(\begin{array}{c}A_{22\left(n+1\right)}\\ A_{18,2,2\left(n+1\right)}\\ A_{9,9,4\left(n+1\right)}\\ A_{8,8,6\left(n+1\right)}\end{array}\right)=\left(\begin{array}{cccc}1.000006761& .0002279253840& .009315026611& .05551458536\\ .05998283956& .09512343412& 3.219304290& 19.31589634\\ .0002679923170& .0009873502040& .03199666409& .1919144288\\ .0001230048643& .0005306869789& .01770755556& .1062348858\end{array}\right)\left(\begin{array}{c}A_{22\left(n\right)}\\ A_{18,2,2\left(n\right)}\\ A_{9,9,4\left(n\right)}\\ A_{8,8,6\left(n\right)}\end{array}\right),$$
$$\lambda _{max}=1.000002379.$$
c)b=4 q=6
$$\left(\begin{array}{c}A_{6(n+1)}\\ A_{4,1,1(n+1)}\end{array}\right)=\left(\begin{array}{cc}1.076310660& .03365228885\\ .1889456076& .01616947736\end{array}\right)\left(\begin{array}{c}A_{6(n)}\\ A_{4,1,1(n)}\end{array}\right),$$
$$\lambda _{max}=1.070278597.$$
q=8
$$\left(\begin{array}{c}A_{8(n+1)}\\ A_{3,3,2(n+1)}\end{array}\right)=\left(\begin{array}{cc}1.019524120& .06535850011\\ .05434716964& 371.8425102\end{array}\right)\left(\begin{array}{c}A_{8(n)}\\ A_{3,3,2(n)}\end{array}\right),$$
$$\lambda _{max}=1.0195338.$$
q=10
$$\left(\begin{array}{c}A_{10(n+1)}\\ A_{4,4,2(n+1)}\end{array}\right)=\left(\begin{array}{cc}1.005151073& .03282144763\\ .03495279823& .007384841623\end{array}\right)\left(\begin{array}{c}A_{10(n)}\\ A_{4,4,2(n)}\end{array}\right),$$
$$\lambda _{max}=1.003999975.$$
q=12
$$\left(\begin{array}{c}A_{12(n+1)}\\ A_{8,2,2(n+1)}\\ A_{6,3,3(n+1)}\end{array}\right)=\left(\begin{array}{ccc}1.001388320& .007332723661& .006026327075\\ .06394351082& .009657753638& .007508055989\\ .01147176558& .004507477340& .003470225654\end{array}\right)\left(\begin{array}{c}A_{12(n)}\\ A_{8,2,2(n)}\\ A_{6,3,3(n+1)}\end{array}\right),$$
$$\lambda _{max}=1.000986557.$$
q=14
$$\left(\begin{array}{c}A_{14(n+1)}\\ A_{6,6,2(n+1)}\\ A_{5,5,4(n+1)}\end{array}\right)=\left(\begin{array}{ccc}1.000380234& .006079338806& .02938006638\\ .008215782561& .006916822481& .03482160520\\ .002872694778& .002513285925& .01260695484\end{array}\right)\left(\begin{array}{c}A_{14(n)}\\ A_{6,6,2(n)}\\ A_{5,5,4(n+1)}\end{array}\right),$$
$$\lambda _{max}=1.000245203.$$
q=16
$$\left(\begin{array}{c}A_{16(n+1)}\\ A_{7,7,2(n+1)}\\ A_{6,6,4(n+1)}\end{array}\right)=\left(\begin{array}{ccc}1.000105428& .002387138955& .01597611501\\ .004013610557& .008712843003& .06128427304\\ .002158295483& .002024493055& .01421013434\end{array}\right)\left(\begin{array}{c}A_{16(n)}\\ A_{7,7,2(n)}\\ A_{6,6,4(n+1)}\end{array}\right),$$
$$\lambda _{max}=1.000061148.$$
q=18
$$\left(\begin{array}{c}A_{18\left(n+1\right)}\\ A_{14,2,2\left(n+1\right)}\\ A_{7,7,4\left(n+1\right)}\\ A_{8,5,5\left(n+1\right)}\end{array}\right)=\left(\begin{array}{cccc}1.000029509& .0004497502350& .008243519261& .008652383249\\ .06187422841& .01967111602& .3403871980& .3574066079\\ .0009123885980& .00009719487958& .001741086261& .001828250148\\ .0007092075724& .0006921244620& .01190983853& .01250523720\end{array}\right)\left(\begin{array}{c}A_{18\left(n\right)}\\ A_{14,2,2\left(n\right)}\\ A_{7,7,4\left(n\right)}\\ A_{8,5,5\left(n\right)}\end{array}\right),$$
$$\lambda _{max}=1.000015271.$$
q=20
$$\left(\begin{array}{c}A_{20\left(n+1\right)}\\ A_{16,2,2\left(n+1\right)}\\ A_{8,8,4\left(n+1\right)}\\ A_{7,7,6\left(n+1\right)}\end{array}\right)=\left(\begin{array}{cccc}1.000008319& .0001645061367& .004348107103& .01946561802\\ .06184409730& .02252802695& .5624149410& 2.519619600\\ .0005078921879& .0001439100811& .003508323524& .01571652660\\ .0001776181030& .0001853563027& .004606804705& .02063832629\end{array}\right)\left(\begin{array}{c}A_{20\left(n\right)}\\ A_{16,2,2\left(n\right)}\\ A_{8,8,4\left(n\right)}\\ A_{7,7,6\left(n\right)}\end{array}\right),$$
$$\lambda _{max}=1.000003788.$$
q=22
$$\left(\begin{array}{c}A_{22\left(n+1\right)}\\ A_{18,2,2\left(n+1\right)}\\ A_{9,9,4\left(n+1\right)}\\ A_{8,8,6\left(n+1\right)}\end{array}\right)=\left(\begin{array}{cccc}1.000002358& .00005834048617& .002108427732& .01263785041\\ .06183577589& .02527262136& .8582479458& 5.149488875\\ .0002481471736& .0002926869789& .009839879848& .05903810198\\ .0001334512767& .0001464026161& .004954070521& .02972423554\end{array}\right)\left(\begin{array}{c}A_{22\left(n\right)}\\ A_{18,2,2\left(n\right)}\\ A_{9,9,4\left(n\right)}\\ A_{8,8,6\left(n\right)}\end{array}\right),$$
$$\lambda _{max}=1.000000954.$$
d)b=5 q=6
$$\left(\begin{array}{c}A_{6(n+1)}\\ A_{4,1,1(n+1)}\end{array}\right)=\left(\begin{array}{cc}1.044360741& .01006023599\\ .3156061304& .003157528970\end{array}\right)\left(\begin{array}{c}A_{6(n)}\\ A_{4,1,1(n)}\end{array}\right),$$
$$\lambda _{max}=1.041302331.$$
q=8
$$\left(\begin{array}{c}A_{8(n+1)}\\ A_{3,3,2(n+1)}\end{array}\right)=\left(\begin{array}{cc}1.009037038& .01421998195\\ .07691145227& .005139564359\end{array}\right)\left(\begin{array}{c}A_{8(n)}\\ A_{3,3,2(n)}\end{array}\right),$$
$$\lambda _{max}=1.010114286.$$
q=10
$$\left(\begin{array}{c}A_{10(n+1)}\\ A_{4,4,2(n+1)}\end{array}\right)=\left(\begin{array}{cc}1.001897065& .005511338227\\ .03479355808& .002253185745\end{array}\right)\left(\begin{array}{c}A_{10(n)}\\ A_{4,4,2(n)}\end{array}\right),$$
$$\lambda _{max}=1.001705201.$$
q=12
$$\left(\begin{array}{c}A_{12(n+1)}\\ A_{8,2,2(n+1)}\\ A_{6,3,3(n+1)}\end{array}\right)=\left(\begin{array}{ccc}1.000405841& .0009777920845& .0007888947248\\ .06389106764& .002988237303& .002324086713\\ .01925523989& .001402362796& .001089147707\end{array}\right)\left(\begin{array}{c}A_{12(n)}\\ A_{8,2,2(n)}\\ A_{6,3,3(n+1)}\end{array}\right),$$
$$\lambda _{max}=1.000358489.$$
q=14
$$\left(\begin{array}{c}A_{14(n+1)}\\ A_{6,6,2(n+1)}\\ A_{5,5,4(n+1)}\end{array}\right)=\left(\begin{array}{ccc}1.000087998& .0006533768126& .003203625771\\ .008004412267& .002154118594& .01079323997\\ .004778768070& .001472733806& .007372115949\end{array}\right)\left(\begin{array}{c}A_{14(n)}\\ A_{6,6,2(n)}\\ A_{5,5,4(n+1)}\end{array}\right),$$
$$\lambda _{max}=1.000097981.$$
q=16
$$\left(\begin{array}{c}A_{16(n+1)}\\ A_{7,7,2(n+1)}\\ A_{6,6,4(n+1)}\end{array}\right)=\left(\begin{array}{ccc}1.000019275& .0002082580899& .001424787973\\ .003961084604& .002614693987& .01832747999\\ .002150223007& .0006222811611& .004359625253\end{array}\right)\left(\begin{array}{c}A_{16(n)}\\ A_{7,7,2(n)}\\ A_{6,6,4(n+1)}\end{array}\right),$$
$$\lambda _{max}=1.000017019.$$
q=18
$$\left(\begin{array}{c}A_{18\left(n+1\right)}\\ A_{1422\left(n+1\right)}\\ A_{774\left(n+1\right)}\\ A_{855\left(n+1\right)}\end{array}\right)=\left(\begin{array}{cccc}1.000004256& .00003193265577& .0005685434366& .0005968699293\\ .06193405435& .006096146153& .1056112005& .1108917618\\ .001017169377& .00006930842110& .001204750315& .001264989081\\ .001191712555& .0002627970450& .004548661710& .004776093729\end{array}\right)\left(\begin{array}{c}A_{18\left(n\right)}\\ A_{14,2,2\left(n\right)}\\ A_{7,7,4\left(n\right)}\\ A_{855\left(n\right)}\end{array}\right),$$
$$\lambda _{max}=1.000002408.$$
q=20
$$\left(\begin{array}{c}A_{20\left(n+1\right)}\\ A_{16,2,2\left(n+1\right)}\\ A_{8,8,4\left(n+1\right)}\\ A_{7,7,6\left(n+1\right)}\end{array}\right)=\left(\begin{array}{cccc}1.000000946& .000009482285099& .0002434259750& .001090262082\\ .06190933065& .006979673785& .1744166676& .7813866910\\ .0004953876882& .00004636993771& .001150699865& .005155112958\\ .0002958221549& .0001041331110& .002597984312& .01163895822\end{array}\right)\left(\begin{array}{c}A_{20\left(n\right)}\\ A_{16,2,2\left(n\right)}\\ A_{8,8,4\left(n\right)}\\ A_{7,7,6\left(n\right)}\end{array}\right),$$
$$\lambda _{max}=1.000000141.$$
q=22
$$\left(\begin{array}{c}A_{22\left(n+1\right)}\\ A_{18,2,2\left(n+1\right)}\\ A_{9,9,4\left(n+1\right)}\\ A_{8,8,6\left(n+1\right)}\end{array}\right)=\left(\begin{array}{cccc}1.000000211& .000002740626906& .00009573710813& .0005742469550\\ .06190294246& .007829018982& .2660888394& 1.596533072\\ .0002452104252& .00008750002217& .2660888394& .01778872116\\ .0001330977107& .00004502739760& .001528690467& .009172137322\end{array}\right)\left(\begin{array}{c}A_{22\left(n\right)}\\ A_{18,2,2\left(n\right)}\\ A_{9,9,4\left(n\right)}\\ A_{8,8,6\left(n\right)}\end{array}\right),$$
$$\lambda _{max}=1.000000086.$$
ACKNOWLEDGEMENT
We wish to thank Dr. S. K. A. Seyed Yagoobi for his careful reading the article and for his constructive comments.
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# SEMI-INCLUSIVE HADRON-HADRON TRANSVERSE SPIN ASYMMETRIES AND THEIR IMPLICATION FOR POLARIZED DIS
## 1 Internal structure of the nucleon at the parton level
For each quark there are three kinds of number densities:
The usual $`q(x)`$
q(x) is the number density of quarks with momentum fraction in the range $`xp/Px+\mathrm{\Delta }x`$. This is mostly measured in DIS.
The longitudinal polarized density $`\mathrm{\Delta }q(x)`$
$`q_\pm (x)`$ is the number density at $`x`$ with spin $`𝐬`$ along (+) or opposite (-) to the spin $`𝐒`$ ($``$) of the nucleon. The new density is
$$\mathrm{\Delta }q(x)=q_+(x)q_{}(x).$$
(1)
It is measured in DIS using a longitudinally polarized nucleon target.
The transverse polarized density $`\mathrm{\Delta }_Tq(x)`$
$`q_{}(x)`$ are the number densities at $`x`$ with transverse spin $`𝐬`$ along ($``$) or opposite ($``$) to the transverse nucleon spin $`𝐒`$ ($``$). The new density is
$$\mathrm{\Delta }_Tq(x)=q_{}(x)q_{}(x).$$
(2)
Note that $`\mathrm{\Delta }_Tq(x)`$ cannot be measured in DIS with a transversely polarized target; $`g_2(x)`$ does not tell us anything about $`\mathrm{\Delta }_Tq(x)`$.
In summary there are 3 independent functions, all equally fundamental, describing the internal structure of the nucleon: $`q(x)`$, $`\mathrm{\Delta }q(x)`$ and $`\mathrm{\Delta }_Tq(x)`$.
How can we measure $`\mathrm{\Delta }_Tq(x)`$ ? The ideal reaction would be Drell-Yan using transversely polarized beam and target, but this has never been done. It is one of the prime aims at RHIC. Can one use semi-inclusive hadron-hadron reactions with a transversely polarized target ? At first sight, yes. At second sight, no. And finally, yes, but one has to introduce a new theoretical idea and thereby it seems possible to resolve the ancient puzzle of the large transverse spin asymmetries.
## 2 The transverse spin asymmetries
There is a mass of data on reactions of the type $`A^{}+BC+X`$ for which the asymmetry $`A_N`$ under the reversal of the transverse spin is measured:
$$A_N=\frac{d\sigma ^{}d\sigma ^{}}{d\sigma ^{}+d\sigma ^{}}.$$
(3)
Some examples are shown in Fig. 1 for $`p^{}p\pi X`$ and $`\overline{p}^{}p\pi X`$ respectively. From looking at many reactions one concludes that:
– the asymmetries are large !
– they increase with $`p_T`$
– they increase with $`x_F`$
– they seem independent of energy
– they occur in a variety of reactions.
For decades there has been no serious theoretical explanation and, as we shall see, the standard approach via perturbative QCD gives $`A_N=0`$.
## 3 Why it is difficult to explain the asymmetries
The standard parton model picture for $`A^{}+BC+X`$ at large momentum transfer is
The hadronic $`A_N`$ depends upon the asymmetry $`\widehat{a}_N`$ at the parton level, i.e. the asymmetry in
$$q_a^{}+q_bq_c+q_d.$$
(4)
But
$$\widehat{a}_N\mathrm{Im}\{(\mathrm{Helicity}\mathrm{Non}\mathrm{flip})^{}(\mathrm{Single}\mathrm{flip})\}.$$
(5)
In lowest order this is doubly zero: there is no helicity flip and the amplitudes are real. Going to higher order doesn’t help. One finds, if one takes $`m_q0`$,
$$\widehat{a}_N=\alpha _s\frac{m_q}{\sqrt{\widehat{s}}}f(\theta ^{})$$
(6)
which gives asymmetries of less than 1%.
## 4 New soft mechanism
Consider, for concreteness, the reaction $`p^{}p\pi ^\pm X`$. Let us concentrate only on the partons in the polarized proton and follow them through the partonic diagram. We assume that the $`\pi `$’s come mainly from the fragmentation of quarks. The notation is the following: $`f_{q/p}`$ is the number density of $`q`$ in $`p`$ and $`D^{\pi /q}`$ the number density of $`\pi `$ in the fragmentation of $`q`$.
Proceeding blindly to sum over all possible spins of the quarks leads to
$`d\sigma ^{}d\sigma ^{}`$ $`=`$ $`[f_{q_a/p^{}}f_{q_a/p^{}}]\widehat{\sigma }D^{\pi /q}+`$ (7)
$`+`$ $`[f_{q_a^{}/p^{}}f_{q_a^{}/p^{}}]\mathrm{\Delta }\widehat{\sigma }[D^{\pi /q_c}(𝐬_𝐜)D^{\pi /q_c}(𝐬_𝐜)]`$
where $`D^{\pi /q}`$ is the usual, unpolarized fragmentation function, and $`\mathrm{\Delta }\widehat{\sigma }`$ will be defined presently; $`𝐬_𝐜`$ is the polarization vector of quark $`c`$. The key question is: which, if any, of these terms are non-zero ?
With usual collinear kinematics
$$f_{q/p^{}}f_{q/p^{}}=0$$
(8)
Reason ?
There are only two independent vectors, $`𝐏`$ and the pseudovector $`𝐒`$. We cannot construct a scalar which depends on $`𝐒`$. Similarly
$$D^{\pi /q}(𝐬)D^{\pi /q}(𝐬)=0$$
(9)
Again, we cannot construct a scalar from the vector $`𝐏`$ and the pseudovector $`𝐬`$.
Thus both terms in Eq. (7) vanish in the collinear kinematics and $`A_N=0`$.
With intrinsic transverse momentum
Now, apparently, we could have
$$f_{q/p(s)}(x,k_T)=f(x,k_T)+\stackrel{~}{f}(x,k_T)𝐒(𝐏\times 𝐤_T)$$
(10)
implying
$$f_{q/p^{}}f_{q/p^{}}0$$
(11)
This mechanism was proposed by Sivers and further studied in . However, it violates time-reversal invariance, so we shall take the first term in Eq. (7) to be zero. Strangely, the analogous mechanism for the fragmentation
$$D^{\pi /q}(𝐬)D^{\pi /q}(𝐬)0$$
(12)
does not violate time-reversal invariance. This is the Collins mechanism . Hence Eq. (7) becomes
$`d\sigma ^{}d\sigma ^{}`$ $`=`$ $`\left[f_{q_a^{}/p^{}}f_{q_a^{}/p^{}}\right]\mathrm{\Delta }\widehat{\sigma }\left[D^{\pi /q_c}(𝐬_𝐜)D^{\pi /q_c}(𝐬_𝐜)\right]`$ (13)
$`=`$ $`[\mathrm{\Delta }_Tq_a]\left[{\displaystyle \frac{d\widehat{\sigma }}{d\widehat{t}}}(a^{}bc^{}d){\displaystyle \frac{d\widehat{\sigma }}{d\widehat{t}}}(a^{}bc^{}d)\right][\mathrm{\Delta }_ND_{\pi /q_c}].`$
In full detail
$`d\sigma ^{}d\sigma ^{}`$ $``$ $`{\displaystyle }dx_adx_bd^2𝐤_T^\pi q(x_b)\mathrm{\Delta }_Tq(x_a)\times `$ (14)
$`\times `$ $`\left[{\displaystyle \frac{d\widehat{\sigma }}{d\widehat{t}}}(a^{}bc^{}d){\displaystyle \frac{d\widehat{\sigma }}{d\widehat{t}}}(a^{}bc^{}d)\right]\mathrm{\Delta }_ND_{\pi /q_c}(z,𝐤_T^\pi ),`$
where the term $`[\frac{d\widehat{\sigma }}{d\widehat{t}}(a^{}bc^{}d)\frac{d\widehat{\sigma }}{d\widehat{t}}(a^{}bc^{}d)]`$ is calculated in PQCD. The result depends on two unknown functions: $`\mathrm{\Delta }_Tq(x)`$ and $`\mathrm{\Delta }_ND_{\pi /q_c}`$, which we can measure by trying to fit the data.
Now, as stressed earlier, the asymmetries are large, so will demand large values of $`\mathrm{\Delta }_Tq(x)`$ and $`\mathrm{\Delta }_ND_{\pi /q_c}`$. However, positivity requires that
$$|\mathrm{\Delta }_ND_{\pi /q_c}|2D_{\pi /q_c}$$
(15)
and the Soffer bound restricts the magnitude of $`\mathrm{\Delta }_Tq(x)`$ :
$$|\mathrm{\Delta }_Tq(x)|\frac{1}{2}\left[q(x)+\mathrm{\Delta }q(x)\right]$$
(16)
where $`\mathrm{\Delta }q(x)`$ is the usual longitudinal polarized quark density.
How important is the Soffer bound ? In Fig. 2 we show a typical picture of $`\mathrm{\Delta }u(x)`$ and $`\mathrm{\Delta }d(x)`$. We see that:
$`\mathrm{\Delta }u(x)`$ is positive everywhere, so that
* $`u(x)+\mathrm{\Delta }u(x)`$ is big
* RHS of Soffer bound is large
* not very restrictive on $`\mathrm{\Delta }_Tu(x)`$
$`\mathrm{\Delta }u(x)`$ is (usually) negative everywhere, so that
* $`d(x)+\mathrm{\Delta }d(x)`$ is small
* RHS of Soffer bound is small
* highly restrictive on $`\mathrm{\Delta }_Td(x)`$
But the measured asymmetries are such that $`A_N^{\pi ^+}A_N^\pi ^{}`$, so that if the $`\pi ^+`$ come mainly from the $`u`$-quarks and the $`\pi ^{}`$ from $`d`$-quarks we expect trouble in getting a large enough asymmetry for $`\pi ^{}`$. Indeed, if we use the Gehrmann-Stirling (GS) $`\mathrm{\Delta }u(x)`$ and $`\mathrm{\Delta }d(x)`$ to bound $`\mathrm{\Delta }_Tu(x)`$ and $`\mathrm{\Delta }_Td(x)`$ we obtain a catastrophic fit to the data (Fig. 3) with $`\chi _{D.O.F}^225`$ !
Can we escape this dilemma ? There is a surprising escape route !
There is an old PQCD argument that requires for quarks, antiquarks and gluons
$$\frac{\mathrm{\Delta }q(x)}{q(x)}1\mathrm{as}x1$$
(17)
which implies that all $`\mathrm{\Delta }q(x)`$ must become positive as $`x1`$. But almost all fits to polarized DIS ignore this condition on the grounds that (i) Eq. (17) is incompatible with DGLAP evolution and that (ii) the data demand a negative $`\mathrm{\Delta }d(x)`$. In fact, these arguments are spurious because (i) DGLAP is not valid as $`x1`$ where one approaches the exclusive region and (ii) the data do not really extend to large $`x`$.
So let us try to impose $`\mathrm{\Delta }q(x)/q(x)1`$ as $`x1`$ in the fits to polarized DIS. In fact, this was done by Brodsky, Burkhardt and Schmidt (BBS) , but the treatment was rough and evolution was not included. This was improved upon by Leader, Sidorov and Stamenov (LSS)<sub>BBS</sub> so as to include evolution and a reasonably good fit to the polarized DIS data was achieved. In Fig. 4 we compare the GS and BBS $`\mathrm{\Delta }d(x)`$. It is clear that the Soffer bound on $`\mathrm{\Delta }_Td(x)`$ will be much less restrictive at large $`x`$ for the BBS case. Indeed, using the BBS $`\mathrm{\Delta }q(x)`$ to bound the $`\mathrm{\Delta }_Tq(x)`$ has a dramatic effect upon our attempts to fit the $`\pi ^\pm `$ asymmetries as can be seen in Fig. 5 where $`\chi _{D.O.F}^2=1.45`$.
In carrying out the fit we made the following simplifications:
The asymmetry is largest at large $`x_F`$ $``$ large $`x`$ is important. Therefore we used only $`u`$ and $`d`$ quarks.
Large $`x_F`$ $``$ large $`z`$ in the fragmentation. Hence we assumed $`u\pi ^+`$, $`d\pi ^{}`$ only.
The unknown functions were parameterized so that the bounds in Eqs. (15) and (16) are automatically satisfied. Thus we took
$$\mathrm{\Delta }_Tq(x)=N_q\left[\frac{x^a(1x)^b}{\frac{a^ab^b}{(a+b)^{a+b}}}\right]\left\{\frac{1}{2}[q(x)+\mathrm{\Delta }q(x)]\right\}$$
(18)
and
$$\mathrm{\Delta }_ND(z)=N_F\left[\frac{z^\alpha (1z)^\beta }{\frac{\alpha ^\alpha \beta ^\beta }{(\alpha +\beta )^{\alpha +\beta }}}\right]\{2D(z)\},$$
(19)
where $`N_{q,F}`$ are real constants with $`|N_{q,F}|1`$, and the functions in square brackets have modulus $`1`$. The fit to the asymmetry data then determines a range of possible $`\mathrm{\Delta }_Tu(x)`$ and $`\mathrm{\Delta }_Td(x)`$ as shown in Fig. 6.
## 5 Implications and conclusions
It seems that the soft Collins mechanism can explain the semi-inclusive transverse spin asymmetries if $`\mathrm{\Delta }_Tu(x)`$ and $`\mathrm{\Delta }_Td(x)`$ are large enough in magnitude.
This, via the Soffer bound, seems to require $`\mathrm{\Delta }q(x)/q(x)1`$ as $`x1`$.
For the $`d`$-quark this implies that $`\mathrm{\Delta }d(x)`$ must change sign and become positive at large $`x`$.
This, in turn, has a significant effect upon the shape of $`g_1^n(x)`$ at large $`x`$. Fig. 7 compares the behavior of $`g_1^n(x)`$ for the “best” usual fit to the polarized parton densities with that from fits satisfying $`\mathrm{\Delta }d(x)/d(x)1`$. The exciting link between transverse asymmetries and polarized DIS emphasizes the importance of extending that polarized DIS measurements to larger x.
Some notes of caution:
(i) The Collins mechanism does not seem able to produce large enough $`A_N`$ at the largest $`x_F`$ measured. However, there does exist another kind of mechanism, outside the framework of the usual parton model, based on correlated quark-gluon densities in a hadron, which can also produce a transverse spin asymmetry. It may be that a superposition of the two mechanisms is needed.
(ii) For either of these mechanisms $`A_N`$ must decrease when $`𝐩_T^\pi `$ becomes much greater than the intrinsic $`𝐤_T^\pi `$. So far there is no sign of such a decrease in the data.
Finally, we wish to re-emphasize the beautiful interplay between, at first sight, quite unrelated aspects of particle physics.
## 6 Acknowledgements
E. L. is grateful to G. Bellettini and M. Greco for their hospitality. This research project was supported by the Foundation for Fundamental Research on Matter (FOM) and the Dutch Organization for Scientific Research (NWO).
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# Groups Quasi-isometric to ℍ²×ℝ
## Introduction
The most powerful geometric tools are those of differential geometry, but to apply such techniques to finitely generated groups seems hopeless at first glance since the natural metric on a finitely generated group is discrete. However Gromov recognized that a group can metrically resemble a manifold in such a way that geometric results about that manifold carry over to the group . This resemblance is formalized in the concept of a “quasi-isometry.” This paper contributes to an ongoing program to understand which groups are quasi-isometric to which simply connected, homogeneous, Riemannian manifolds by proving that any group quasi-isometric to $`^2\times `$ is a finite extension of a cocompact lattice in $`Isom(^2\times )`$ or $`Isom(\stackrel{~}{SL}(2,))`$.
###### Theorem
For any group $`\mathrm{\Gamma }`$ quasi-isometric to the hyperbolic plane cross the real line, there is an exact sequence
$$0\mathrm{@}>>>A\mathrm{@}>>>\mathrm{\Gamma }\mathrm{@}>>>G\mathrm{@}>>>0$$
where $`A`$ is virtually infinite cyclic and $`G`$ is a finite extension of a cocompact Fuchsian group.
With this result, the question of which finitely generated groups are quasi-isometric to each of Thurston’s eight geometries remains open only for Sol. Tukia for $`n>2`$ and Gabai , or alternatively Casson–Jungreis , and Tukia in dimension 2 show that any group quasi-isometric to $`^n`$ can be realized as a finite extension of a discrete cocompact subgroup of the isometries of $`^n`$. The analogous result for $`^2\times `$ is not true, because $`^2\times `$ is quasi-isometric to $`\stackrel{~}{SL}(2,)`$, the universal cover of $`SL(2,)`$. Gromov’s work on groups of polynomial growth implies that any group quasi-isometric to $`^n`$ is a finite extension of a discrete, cocompact subgroup of the isometry group of $`^n`$. Similarly for Nil. Kapovich and Leeb show that quasi-isometries preserve the geometric decompostion of Haken manifolds. Their paper, which was written after the present paper, contains another proof, along quite different lines, of the main result in this paper. In their 1996 preprint , Kleiner and Leeb generalized the main results of the present paper to products of simply connected nilpotent Lie group with a symmetric space of non-compact type with no Euclidean de Rham factors. Chow has shown that every group quasi-isometric to the complex hyperbolic plane is a finite extension of a discrete, cocompact subgroup of the isometry group of the complex hyperbolic plane. Pansu has shown that in quaternionic and Cayley hyperbolic space every quasi-isometry is within bounded distance of an isometry, so that in fact every group quasi-isometric to one of these spaces is a finite extension of a group naturally isomorphic to a discrete cocompact subgroup of the isometry group. A series of papers by Pansu , Schwartz , Farb-Schwartz , Kleiner-Leeb , Eskin-Farb , and Eskin succeeded in classifying lattices in semi-simple Lie groups up to quasi-isometry. See for a survey. Of course, the question may be asked for metric spaces other than symmetric spaces or even manifolds. Papers by Gromov and Stallings give the result for trees . Farb and Mosher prove qusi-isometric rigidity for the solvable Baumslag-Solitar groups.
The heart of the paper is contained is Sections 3 and 5. Section 1 gives some background and definitions and sets the notation for the paper. Section 2 describes a “quasi-action” of a group $`\mathrm{\Gamma }`$ quasi-isometric to $`^2\times `$ on $`^2\times `$. Section 3 uses geometric arguments to obtain an action of $`\mathrm{\Gamma }`$ on $`^2`$ by quasisymmetric maps. Section 4 describes the application of results of Hinkkanen , Gabai , Casson–Jungreis , and Tukia to show that this action is conjugate by a quasisymmetric map to an action by a Möbius group. Section 5 uses largely algebraic methods to show that the resulting action on $`^2`$ by isometries must have been properly discontinuous, and thus there is a map $`\mathrm{\Phi }:\mathrm{\Gamma }G`$ where $`G`$ is a discrete, cocompact subgroup of the isometries of $`^2`$. Section 6 uses geometric arguments from Section 3 to show that $`\mathrm{ker}\mathrm{\Phi }`$ must be quasi-isometric to $``$.
The aim of Section 3 is to show that the image of a horizontal hyperbolic plane in $`^2\times `$ under any quasi-isometry induced by $`\mathrm{\Gamma }`$ must have sufficient horizontal expanse there is a natural map from $`^2`$ to itself. The intuition behind the proof is that slicing $`^2\times `$ vertically gives Euclidean planes, which have much less area than hyperbolic planes, so the image of a horizontal hyperbolic plane under a quasi-isometry of $`^2\times `$ cannot be contained in a vertical slice. The proof formalizes this idea by showing that if the image of vertical geodesics, from a maximal family of geodesics all within a vertical cylinder all at least a certain distance apart, were to quasicross a vertical cross-section, then the number of quasicrossing points, all at least a certain distance apart, exceeds the number possible in a Euclidean disk.
Key to Section 5 is the notion of semilocal growth of a finitely generated subgroup of a Lie group. In particular, Theorem 5.19 states that any finitely generated, non-elementary, non-discrete subgroup of $`PSL(2,)`$ has superpolynomial semilocal growth. The semilocal growth of a subgroup is strictly larger than its local growth, a notion due to Carrière . In fact, all of the results of Section 5 can be obtained replacing semilocal growth with local growth, but the notion of semilocal growth seems more natural in this setting. It perhaps shows the importance of local growth and Theorem 5.19, which had been previously obtained by Carrière and Ghys , that they arose independently in different contexts.
I would like to thank my advisor, Geoffrey Mess, for his guidance throughout this project and for his enthusiasm for the subject. Thanks also go to Mladen Bestvina, Francis Bonahon and Bob Edwards for many helpful conversations, and to Etienne Ghys for suggesting the proof of Lemma 5.8. I am grateful to the superb anonymous referee whose careful reading and detailed comments led to a significantly more elegant and succinct exposition. Finally warmest thanks to my friends and family for their generous support, encouragement, and interest throughout this endeavor.
## 1. Background and Definitions
Two metric spaces are quasi-isometric if there is a relation between them which, except locally, does not increase or decrease distances too much and has the property that every point in each space is close to a point which is related to a point in the other space.
###### Definition \newitem
A map $`\psi :XY`$ is a $`(\lambda ,ϵ,\delta )`$–quasi-isometry if it satisfies
There is some discrepancy in the literature as to whether quasi-isometries are required to be almost surjective or not. Throughout this paper, when we omit the $`\delta `$ in the definition of a quasi-isometry, we will mean that we are not requiring the almost surjectivity condition, or in other words, the map is a quasi-isometry in the sense given above only onto the range of the map.
Two metric spaces are quasi-isometric if there exist quasi-isometries between them that are almost inverses of each other. This idea is formalized in the following definition.
###### Definition \newitem
Two metric spaces $`X`$ and $`Y`$ are $`(\lambda ,ϵ,\delta )`$–quasi-isometric if there exist $`(\lambda ,ϵ,\delta )`$–quasi-isometries $`\psi :XY`$ and $`\omega :YX`$ such that for some $`\kappa `$
Any group $`\mathrm{\Gamma }`$ with generating set $`S`$, where $`S`$ contains the inverse of any element in $`S`$, has a natural metric, the word metric $`d_S`$, given by
$$d_S(g,h)=\mathrm{min}\{n|ga_1a_2\mathrm{}a_n=h\text{ with }a_iS\}.$$
Word metrics coming from two different finite generating sets $`S`$ and $`S^{}`$ are equivalent in the sense that there is a constant $`\lambda `$ such that
$$\frac{1}{\lambda }d_S(x,y)d_S^{}(x,y)\lambda d_S(x,y).$$
So up to quasi-isometry there is a natural metric on a finitely generated group.
###### Note \newitem
Whenever we say “a group is quasi-isometric” implicitly we will be talking about a finitely generated group endowed with one of the equivalent word metrics.
## 2. A quasiaction of $`\mathrm{\Gamma }`$ on $`^2\times `$
Throughout this paper $`\mathrm{\Gamma }`$ will denote a group $`(\lambda ,ϵ,\delta )`$–quasi-isometric to $`X`$ via the word metric on $`\mathrm{\Gamma }`$ coming from a finite generating set.
###### Definition \newitem
A $`(\lambda _0,ϵ_0,\delta _0,\kappa )`$–quasiaction of a group $`\mathrm{\Gamma }`$ on a metric space $`X`$ consists of a family of $`(\lambda _0,ϵ_0,\delta _0)`$–quasi-isometries $`\varphi _u:XX`$ such that $`d(\varphi _u\varphi _u^{}(x),\varphi _{uu^{}}(x))\kappa `$ for all $`u,u^{}\mathrm{\Gamma }`$ and all $`xX`$.
Note that were $`\kappa `$ zero, we would have an action of $`\mathrm{\Gamma }`$ on $`X.`$ Also it’s worth pointing out that the $`\varphi _u`$’s are not necessarily homeomorphisms.
###### Observation\newitem
A group $`\mathrm{\Gamma }`$ quasi-isometric to a $`^2\times `$ has a natural quasiaction on $`^2\times `$, with $`\kappa =\lambda \delta +ϵ`$, given by $`\varphi _u=\psi u\omega `$, where $`\psi `$ and $`\omega `$ are the $`(\lambda ,ϵ,\delta )`$–quasi-isometries from the definition of $`^2\times `$ and $`\mathrm{\Gamma }`$ being quasi-isometric, and $`u\mathrm{\Gamma }`$ acts on $`\mathrm{\Gamma }`$ by left multiplication.
## 3. Constructing an action of $`\mathrm{\Gamma }`$ on $`^2`$.
Throughout this section “vertical” will refer to the $``$ coordinates and “horizontal” to the $`^2`$ coordinates. Let $`\pi `$ be the projection of $`^2\times `$ onto $`^2\times \{0\}`$. In order to obtain an action of $`\mathrm{\Gamma }`$ on $`^2`$ of $`^2`$ we want to show that each $`\pi \varphi _u`$ restricted to the horizontal plane $`^2\times \{0\}`$ is a quasi-isometry. The idea is to show that each $`\varphi _u`$ must preserve horizontal and vertical in some rough sense. The proof exploits the fact that there is a lot more room in a disk in the hyperbolic plane than in a disk in the Euclidean plane. We begin with two estimates exhibiting this difference.
###### Estimate \newitem
An upper bound for the number of disks of radius $`r`$ in a disk of radius $`R`$ all at least $`2s`$ apart in Euclidean space is given by
$$\left(\frac{(R+s)}{(r+s)}\right)^2.$$
Recall that the area of a hyperbolic disk of radius $`R`$ is given by $`2\pi (\mathrm{cosh}R1).`$
###### Estimate \newitem
A lower bound for a maximal number of disks of radius $`r`$ in a disk of radius $`R`$ all at least $`2s`$ apart in hyperbolic space is given by
$$\frac{e^R2}{2(\mathrm{cosh}(2(r+s))1)}.$$
By “the horizontal distance between points $`x`$ and $`y`$,” I will mean the distance in $`^2\times \{0\}`$ between $`\pi (x)`$ and $`\pi (y)`$, and by “vertical distance,” the distance between the $``$–coordinates.
###### Proposition \newitem
Let $`\varphi :^2\times ^2\times `$ be a $`(\lambda ,ϵ)`$–quasi-isometry. Let $`C_0`$ be a vertical solid cylinder of radius $`r`$, and let $`p_0`$ and $`q_0`$ be points in $`C_0`$ such that $`\varphi (p_0)`$ and $`\varphi (q_0)`$ have vertical distance no greater than $`h_0`$. Then $`p_0`$ and $`q_0`$ must be no farther apart than
$$ch_0+c$$
where $`c`$ is a constant depending only on $`\lambda `$, $`ϵ`$ and $`r`$.
###### Demonstration Proof
Let $`d(p_0,q_0)=L`$. Let $`w`$ denote the horizontal distance between $`\varphi (p_0)`$ and $`\varphi (q_0)`$ and $`h`$ the vertical. Now
$`L/\lambda ϵ`$ $`d(\varphi (p_0),\varphi (q_0))`$
$`=\sqrt{w^2+h^2}`$
$`w+h_0.`$
So
$$wL/\lambda ϵh_0.$$
Project $`\varphi (p_0)`$ and $`\varphi (q_0)`$ onto $`^2\times \{0\}`$. Let $`\alpha `$ be the geodesic which intersects the geodesic between the projected images of $`\varphi (p_0)`$ and $`\varphi (q_0)`$ perpendicularly at the point halfway between them.
For every vertical solid cylinder $`C`$ of radius $`r`$ whose central axis is distance $`D`$ from the central axis of $`C_0`$, there are points $`p`$ and $`q`$ that are the translation of $`p_0`$ and $`q_0`$ under the translation that takes the central axis of $`C_0`$ to the central axis of $`C`$ and preserves the vertical height of each point. Let $`\gamma `$ be the geodesic between $`p`$ and $`q`$. The points $`\varphi (p)`$ and $`\varphi (q)`$ are on opposite sides of $`\{\alpha \}\times `$, if
$$d(\varphi (p_0),\varphi (p))1/2w$$
and
$$d(\varphi (q_0),\varphi (q))1/2w.$$
Let $`R`$ be the maximal distance $`D`$ for which these inequalities hold:
$$R=\frac{1}{2\lambda }(L/\lambda 3ϵh_0).$$
Let $`𝒞`$ be a maximal set of cylinders of radius $`r`$ all $`s`$ apart within the cylinder of radius $`R`$ about $`C_0`$, where $`s=\lambda (1+2ϵ)`$. For each cylinder $`C`$ in $`𝒞`$, let $`x`$ be a point in $`\{\alpha \}\times `$ that is as close as possible to $`\varphi (\gamma )`$. Such points are called quasicrossing-points of $`\varphi (\gamma )`$ with $`\{\alpha \}\times `$ since $`d(x,\varphi (\gamma ))ϵ/2`$. Two quasicrossing-points $`x`$ and $`x^{}`$ associated with cylinders $`C`$ and $`C^{}`$ in $`𝒞`$ must be at least $`1`$ apart as
$`d(x,x^{})`$ $`d(y,y^{})d(x,y)d(x^{},y^{})`$
$`s/\lambda ϵϵ/2ϵ/2=1,`$
where $`y`$ and $`y^{}`$ are points on $`\varphi (\gamma )`$ and $`\varphi (\gamma ^{})`$ within $`ϵ/2`$ of $`x`$ and $`x^{}`$ respectively.
To see that
$$d(x,x^{})\lambda (L+R)+2ϵ,$$
let $`z`$ and $`z^{}`$ be points on $`\gamma `$ such that $`\varphi (z)=y`$ and $`\varphi (z^{})=y^{}`$. Since all points on $`\gamma `$ and $`\gamma ^{}`$ lie within a cylinder of radius $`R`$ and height $`L`$, $`d(z,z^{})L+R`$. So
$`d(x,x^{})`$ $`d(y,y^{})+d(x,y)+d(x^{},y^{})`$
$`\lambda d(z,z^{})+ϵ+ϵ`$
$`\lambda (L+R)+2ϵ.`$
By Estimate 3.0 there are at least $`\frac{e^R2}{2(\mathrm{cosh}(2(r+s/2))1)}`$ vertical cylinders of radius $`r`$ in $`𝒞`$. Let $`X`$ be a set of quasicrossing-points, one for each cylinder $`C`$ in $`𝒞`$. By estimate 1 there are at most
$$(2\lambda (L+R)+4ϵ+1)^2$$
such points. Thus
$$\frac{e^R2}{2(\mathrm{cosh}(2(r+s/2))1)}(2\lambda (L+R)+4ϵ+1)^2.$$
Solving for $`L`$ in the definition of $`R`$ gives
$$L=2\lambda ^2R+3\lambda ϵ+\lambda h_0.$$
Substituting this expression for $`L`$ in the above inequality shows that $`e^R`$ must be at most some quadratic function of $`\lambda `$, $`ϵ`$, $`r`$ and $`h_0`$. Furthermore
$$e^Rb_2R^2+b_1R+b_0,$$
where the coefficients $`b_0`$, $`b_1`$ and $`b_2`$ are polynomial functions of $`h_0`$ of degree no more than 2. The first four terms of the Taylor expansion of $`e^R`$ tell us that $`R`$ must be no more than $`b_2+2b_1+6b_0`$. So $`R`$ is no more than some quadratic function of $`h_0`$. The definition of $`R`$ together with this bound give a similar bound on $`L`$,
$$Lc_2h_0^2+c_1h_0+c_0,$$
where the coefficients depend only on $`\lambda `$, $`ϵ`$, and $`r`$. We can get a linear bound on $`L`$ with respect to $`h_0`$ as follows. Let $`c=c_2+c_1+c_0`$. A geodesic between $`\varphi (p_0)`$ and $`\varphi (q_0)`$ of vertical distance $`hh_0`$ can be split up into no more than $`h_0+1`$ pieces of vertical distance no greater than 1. Using the above bound on each piece and adding them together, we know that $`p_0`$ and $`q_0`$ can be no farther apart than $`c(h_0+1)`$. ∎
###### Remark \newitem
For future reference note that $`c=aS`$ where $`S=2(\mathrm{cosh}(2r+\lambda (1+2ϵ))1)`$ and $`a`$ is some constant depending only on $`\lambda `$ and $`ϵ`$.
###### Corollary \newitem
Let $`\varphi _u:^2\times ^2\times `$ be a ($`\lambda ,ϵ`$)–quasi-isometry coming from the quasiaction of $`\mathrm{\Gamma }`$ on $`^2\times `$. Let $`x`$ and $`y`$ be points $`D`$ apart in some horizontal $`^2`$. Then the horizontal distance $`l`$ between $`\varphi _u(x)`$ and $`\varphi _u(y)`$ must be at least
$$l\mathrm{ln}((\frac{\frac{D}{\lambda }ϵ}{(a(2\kappa +2\delta +1)}1)\lambda (2ϵ+1)+1$$
where $`\kappa `$ is the quasiaction constant of Observation 2.0, and $`a`$ depends only on $`\lambda `$ and $`ϵ`$.
###### Demonstration Proof
Since $`d(\varphi _{u^1}\varphi _u(x),x)d(\varphi _{u^1}\varphi _u(x),\varphi _e(x))+d(\varphi _e(x),x)\kappa +\delta `$ and similarly for $`y`$, the vertical distance between $`\varphi _{u^1}(\varphi _u(x))`$ and $`\varphi _{u^1}(\varphi _u(y))`$ must be less than $`2(\kappa +\delta )`$. Hence by Proposition 3.0, $`d(\varphi _u(x),\varphi _u(y))2c(\kappa +\delta )+c`$. By Remark 3.0,
$$d(\varphi _u(x),\varphi _u(y))2a\mathrm{cosh}(2r+\lambda (2ϵ+1))1)(2(\kappa +\delta )+1),$$
where $`r`$ is the radius of a vertical cylinder containing $`\varphi _u(x)`$ and $`\varphi _u(y)`$, so can be taken to be half the horizontal distance $`l`$ between $`\varphi _u(x)`$ and $`\varphi _u(y)`$. It is $`l`$ we are trying to bound from below. So
$$\frac{D}{\lambda }ϵd(\varphi _u(x),\varphi _u(y))(2a\mathrm{cosh}(l+\lambda (2ϵ+1)1)(2(\kappa +\delta )+1),$$
from which the result follows. ∎
###### Proposition \newitem
$`\pi \varphi :^2\times \{0\}^2\times \{0\}`$ is a ($`\lambda ^{},ϵ^{}`$)–quasi-isometry where $`\lambda ^{}=\mathrm{max}\{\lambda ,D^{}\}`$ and $`ϵ^{}=\mathrm{max}\{ϵ,1\}`$ taking $`D^{}=\lambda a(2\kappa +2\delta +1)(e^{\lambda (2ϵ+1)}+1)+ϵ).`$
###### Demonstration Proof
The projection $`\pi `$ is distance nonincreasing, so for any two points $`x`$ and $`y`$ in $`^2\times \{0\}`$ we know
$$d(\pi \varphi (x),\pi \varphi (y))\lambda d(x,y)+ϵ.$$
We need to show that $`\pi \varphi `$ does not decrease distances too much. By Corollary3.0 the images under $`\pi \varphi `$ of any two points at least distance $`D^{}`$ apart cannot be closer than 1 unit apart. Let $`L=d(\pi \varphi (x),\pi \varphi (y))`$, and let $`\gamma `$ be the geodesic between $`\pi \varphi (x)`$ and $`\pi \varphi (y)`$. Let the $`x_i`$ be $`NL+1`$ points along $`\gamma `$ such that successive $`x_i`$ are within distance 1 of each other and $`x_0=\pi \varphi (x)`$ and $`x_N=\pi \varphi (y)`$. So successive quasi-preimages of the $`x_i`$ must be no more than $`D^{}`$ apart. Hence
$$d(x,y)D^{}(L+1)=D^{}d(\pi \varphi (x),\pi \varphi (y))+D^{}.$$
Therefore,
$$d(\pi \varphi (x),\pi \varphi (y))\frac{1}{D^{}}d(x,y)1.$$
A quasisymmetric map of $`S^1`$ viewed as the boundary of the Poincaré model for the hyperbolic plane is a map that extends to a quasiconformal map of $`^2`$. If we view $`^2`$ as the upper half plane of $``$, Beurling and Ahlfors showed that $`f`$ is quasisymmetric and fixes the point at infinity exactly when there exists a constant $`ϵ`$ such that
$$1/ϵ\frac{f(x+t)f(x)}{f(x)f(xt)}ϵ.$$
It is well-known (see for example ) that any quasi-isometry of $`^2`$ extends to a map on $`^2`$ that is quasisymmetric and the quasisymmetry constant only depends on the quasi-isometry constants. The central idea is that the image of any geodesic under a quasi-isometry lies in the $`B`$ neighborhood of a geodesic, where the constant $`B`$ depends only on $`\lambda `$ and $`ϵ`$. Although we only had a quasiaction of $`\mathrm{\Gamma }`$ on $`^2\times `$ and also on $`^2\times \{0\}`$, we get a true action on $`^2`$ for the following reasons. Since the images of a geodesic in $`^2\times \{0\}`$ under $`\pi \varphi _{uu^{}}`$ and $`(\pi \varphi _u)(\pi \varphi _u^{})`$ are within bounded distance of each other, they must be within bounded distance of the same geodesic. To see where a point $`x`$ in $`^2`$ goes under a quasi-isometry $`\varphi `$, take a geodesic ray $`\gamma (t)`$ with $`x`$ as its endpoint. The subset $`\varphi (\gamma )`$ is within a bounded distance of some geodesic ray. Map $`x`$ to the appropriate endpoint of this geodesic ray. Since this geodesic ray is the same for both maps, both maps must give the same action on the boundary.
###### Corollary \newitem
There is a canonical surjective homomorphism $`\mathrm{\Xi }:\mathrm{\Gamma }F`$ where $`F`$ is the uniformly quasisymmetric group consisting of the boundary values of the $`\pi \varphi _u`$’s.
## 4. The action of $`\mathrm{\Gamma }`$ on $`^2`$ is conjugate by a quasisymmetric map to an action by a Möbius group.
Let $`F^+`$ be the group of orientation preserving elements of $`F`$. This section shows that $`F^+`$ may be conjugated by a quasisymmetric map to a Möbius group $`G`$. The results of this section were previously known. For example, Lemma 4.3 is a special case of Theorem 9 in , and Lemma 4.4 follows trivially from results in Section C of .
For $`F^+`$ not discrete in the topology of pointwise convergence (a case we will rule out in section 5), we may use the following theorem of Hinkkanen.
###### Theorem \newitem(Hinkkanen)
Let $`G`$ be a uniformly quasisymmetric group containing a sequence of distinct elements that tend to the identity pointwise. Then $`G`$ is a quasisymmetric conjugate of a Möbius group.
Throughout the rest of the section we will assume $`F^+`$ is discrete. Any discrete uniformly quasisymmetric group of orientation preserving homeomorphisms is a convergence group, so the following theorem holds for $`F^+`$.
###### Theorem \newitem(Casson–Jungreis , Gabai , Tukia )
$`G`$ is a convergence group if and only if $`G`$ is conjugate in Homeo($`S^1`$) to the restriction of a Fuchsian group.
This theorem was proved by Gabai , and independently by Casson–Jungreis , building on work of Tukia . The rest of this section is devoted to showing that the map conjugating $`F^+`$ to a Möbius group may be taken to be quasisymmetric.
###### Lemma \newitem
Say $`\mathrm{\Theta }:\mathrm{\Gamma }^2`$ and $`\mathrm{\Psi }:\mathrm{\Gamma }^2`$ are two maps which induce a quasi-isometry between $`\mathrm{\Gamma }`$ and $`^2`$. Then the quasisymmetric maps they induce on $`^2`$ are conjugate by a quasisymmetric map.
###### Demonstration Proof
Define $`H:\mathrm{\Theta }(\mathrm{\Gamma })\mathrm{\Psi }(\mathrm{\Gamma })`$ to be the composition of $`\mathrm{\Psi }`$ with a quasi-inverse of $`\mathrm{\Theta }`$. As we saw in the previous section $`H`$ can be extended to a quasi-isometry of $`^2`$ to itself, which induces a quasisymmetric map $`h`$ on $`^2`$. By examining the construction, we see that $`h`$ conjugates the map induced on $`^2`$ by $`\mathrm{\Psi }`$ to the one induced by $`\mathrm{\Theta }`$. ∎
###### Lemma \newitem
For any $`(\lambda ,ϵ)`$–quasi-isometry $`f:^2^2`$ that has the same boundary values as some isometry $`g:^2^2`$, there is a constant $`L`$ such that $`d(f(x),g(x))L`$ for all $`x^2`$, where $`L`$ is dependent only on $`\lambda `$ and $`ϵ`$.
###### Demonstration Proof
Let $`y_0,y_1,y_2`$ be vertices of an ideal triangle such that $`x`$ is within the same distance, say $`P`$, of each of the geodesics connecting the three points. Denote by $`Y_i`$ the geodesic connecting $`y_i`$ and $`y_{(i+1)mod3}`$. As was explained in the paragraph following the proof of Lemma 3.0 , $`f(Y_i)`$ is a curve that remains within distance $`D`$ (depending only on $`\lambda `$ and $`ϵ`$) of $`g(Y_i)`$. The point $`f(x)`$ is within distance $`\lambda P+ϵ`$ of each $`f(Y_i)`$, so must be within $`\lambda P+ϵ+D`$ of each $`g(Y_i)`$. The intersection of these regions has bounded diameter, say $`L`$. Since $`g(x)`$ is certainly within this region, $`f(x)`$ must be within $`L`$ of $`g(x)`$. ∎
###### Lemma \newitem
Any finitely generated, discrete, uniformly $`(\lambda ,ϵ)`$-quasisymmetric group $`G`$ acting on $`S^1`$, that induces a cocompact action on the space of triples $`T`$, is quasi-isometric to $`^2`$.
###### Demonstration Proof
To any ordered triple of distinct points in $`S^1`$ we associate a point in hyperbolic space as follows. Connect the first two points by a geodesic and drop a perpendicular from the third. The intersection will be the point associated to the triple. Choose a triple $`t_0`$ whose associated point is $`x_0`$. Let $`f:G^2`$ be the map sending an element $`aG`$ to the point associated to the triple $`a(t_0)`$. We will show that $`f`$ is a quasi–isometry. For any element $`aG`$, choose a $`(\lambda ,ϵ)`$–quasi-isometry $`\eta _a:^2^2`$ with boundary values $`a`$. Note $`\eta _a(f(a^{}))`$ must be within $`2B`$ of $`f(aa^{})`$ (where $`B`$ was defined in the paragraph following Proposition 3.0), since $`\eta _a(f(a^{}))`$ must be within $`B`$ of each of the geodesics used to construct $`f(aa^{})`$.
Since $`G`$ is cocompact on the space of triples, there is a constant $`E`$ such that every point in $`^2`$ is within $`E`$ of some point in the image of $`f`$. Let $`J^{}`$ be a finite generating set for $`G`$. Enlarge $`J^{}`$ to $`J`$ by including all elements $`a`$ such that $`d(\eta _a(x_0),x_0)3\lambda E+ϵ+6B`$. By the discreteness of $`G`$, there are only finitely many such elements. Let $`M=\mathrm{max}_{bJ}d(f(b),x_0)`$. We are now ready to check that $`f`$ is a quasi-isometry. If $`d(a,a^{})=m`$, then for some $`b_iJ`$, $`a=a^{}b_1b_2\mathrm{}b_m`$. So,
$`d(f(a),f(a^{}))`$ $`d(f(a^{}b_1\mathrm{}b_m),f(a^{}b_1\mathrm{}b_{m1}))`$
$`+d(f(a^{}b_1\mathrm{}b_{m1}),f(a^{}b_1\mathrm{}b_{m2}))+\mathrm{}+d(f(a^{}b_1),f(a^{}))`$
$`d(\eta _{a^{}b_1\mathrm{}b_{m1}}f(b_m),\eta _{a^{}b_1\mathrm{}b_{m1}}f(e))+\mathrm{}+d(\eta _a^{}f(b_1),\eta _a^{}f(e))+4mB`$
$`\lambda d(f(b_m),f(e))+ϵ+\mathrm{}+\lambda d(f(b_1),f(e))+ϵ+4mB`$
$`(\lambda M+ϵ+4B)m`$
$`=(\lambda M+ϵ+4B)d(a,a^{}).`$
To get the lower bound, we compute as follows. Let $`d(f(a),f(a^{}))=L`$. Divide the geodesic between $`f(a)`$ and $`f(a^{})`$ into $`[L/E]+1`$ segments of length no more than $`E`$, with endpoints $`x_0,x_1,\mathrm{},x_N`$, where $`f(a)=x_0`$ and $`f(a^{})=x_N`$. Each $`x_i`$ is within $`E`$ of some point $`f(a_i)`$. By construction, for every $`i`$, $`d(f(a_{i1}),f(a_i))3E`$. We wish to show that there is a generator taking $`a_i`$ to $`a_{i+1}`$ for all $`i`$. Recall that $`d(\eta _a(f(a^{}),f(aa^{}))2B`$ and $`x_0=f(e)`$. Note that
$`d(x_0,f(a_{i1}^1a_i))`$ $`d(x_0,\eta _{a_{i1}^1}(f(a_{i1})))+d(\eta _{a_{i1}^1}(f(a_{i1})),\eta _{a_{i1}^1}(f(a_i)))`$
$`+d(\eta _{a_{i1}^1}(f(a_i)),f(a_{i1}^1a_i))`$
$`2B+3\lambda E+ϵ+2B.`$
Furthermore,
$`d(\eta _{a_{i1}^1a_i}(x_0),x_0)`$ $`d(x_0,f(a_{i1}^1a_i)))+d(f(a_{i1}^1a_i),\eta _{a_{i1}^1a_i}(x_0)`$
$`3\lambda E+ϵ+6B.`$
So for all $`i`$, $`a_{i1}^1a_i=b_i`$ is in the generating set for $`G`$.
Since
$$a^{}=a_N=a_{N1}b_N=\mathrm{}=ab_1b_2\mathrm{}b_N,$$
we have
$`d(a,a^{})`$ $`N`$
$`=[L/E]+1`$
$`{\displaystyle \frac{1}{E}}d(f(a),f(a^{}))+1.`$
###### Observation \newitem
Section 2 together with the last paragraph of section 3 shows how a quasi-isometry of a group $`G`$ with the hyperbolic plane induces an action on $`^2`$. It is easy to check that the action of $`F^+`$ on $`^2`$ that we get from the quasi-isometry of $`F^+`$ with $`^2`$ given by the Lemma 4.0 is the same action of $`F^+`$ on $`^2`$ that we started with.
###### Proposition \newitem
$`F^+`$ is conjugate to a Möbius group by a quasisymmetric homeomorphism $`h:S^1S^1`$.
###### Demonstration Proof
For $`F^+`$ not discrete this is a result of Hinkkanen .
Gabai \[G\], or Casson–Jungreis , together with Tukia show that any discrete convergence group (in particular, any quasisymmetric group) is conjugate to a finite extension of a Fuchsian group by some homeomorphism $`h_1`$ of $`S^1`$. Our $`F^+`$ acts cocompactly on the space of triples, so $`G^+=h_1F^+h_1^1`$ must also. Thus $`G^+`$ is a Fuchsian group acting discretely and cocompactly on the space of triples, so when we extend the action to $`^2`$ we get a discrete cocompact group of hyperbolic isometries. This gives us a quasi-isometry of $`G^+`$ with $`^2`$. We get a different quasi-isometry $`\mathrm{\Psi }:G\mathrm{\Psi }(G)`$ by using the isometry of $`G^+`$ with $`F^+`$ and applying Lemma 4.0 to $`F^+`$. Applying Lemma 4.0 to these two quasi-isometries, which have the same boundary values by Observation 4.0 , yields the desired result.
Let $`\mathrm{\Phi }:\mathrm{\Gamma }G`$ be the homomorphism we have obtained where $`G`$ is either the Möbius group $`G^+`$ or a $`/2`$ extension of $`G^+`$, depending on whether $`F`$ contained orientation reversing elements or not.
## 5. The discreteness of the image of $`\mathrm{\Phi }`$
Our goal now is to show that $`G`$ is discrete. The idea is that any non-elementary Möbius group which is not discrete has many more small elements than $`G`$. More precisely, choose a left-invariant metric on $`PSL(2,)`$. For any $`ϵ>0`$, denote by $`N_ϵ`$ the $`ϵ`$-neighborhood of the identity in $`PSL(2,)`$. For any finitely generated subgroup $`H`$ of $`PSL(2,)`$ , let $`H_ϵ^n`$ denote the set of $`h_uN_ϵ`$ such that $`u`$ has word length less than or equal to $`n`$. Our claim is that for some $`ϵ`$ (specified later), $`|G_ϵ^n|`$, the number of elements in $`G_ϵ^n`$, grows linearly with $`n`$, while for any finitely generated nonelementary group $`H`$ that is not discrete $`|H_ϵ^n|`$ grows exponentially in $`n`$.
###### Definition \newitem
Let $`G`$ be a finitely generated group with word metric $`d`$ imbedded in another group $`L`$ with metric $`\rho `$. Then the semilocal growth of $`G`$ in $`L`$ is defined to be the growth rate of the number of elements in
$$G_ϵ^n=\{gG|d(g,e)n,\rho (g,e)ϵ\}.$$
###### Note \newitem
This notion of local growth is strictly bigger than Carrière’s notion of local growth , which counts only elements $`g`$ such that the subwords of increasing length making up $`g`$ are all in the $`ϵ`$-neighborhood of the identity. More, precisely it counts only $`g=a_{i_1}a_{i_2}\mathrm{}a_{i_m}`$ where $`g_j=a_{i_1}a_{i_2}\mathrm{}a_{i_j}N_ϵ`$ for all $`jm`$, where the $`a_k`$ are generators for $`G`$. All of the results hold equally well for local growth as for semilocal growth, but as local growth imposed an extra, unnecessary condition on the elements we are counting, we will prove the results for semilocal growth.
### $`G`$ has linear semilocal growth in $`PSL(2,)`$
Recall that $`\psi :\mathrm{\Gamma }^2\times `$ is the quasi-isometry that came from the definition of $`\mathrm{\Gamma }`$ and $`^2\times `$ being quasi-isometric.
###### Observation \newitem
The number of $`u\mathrm{\Gamma }`$ such that the images of $`z_0=\psi (e)`$ under $`\varphi _u`$ lies in a vertical cylinder $`K`$ of height 1 and radius $`R`$ centered about the vertical geodesic through $`z_0`$ is bounded by some constant $`N`$ depending only on $`\lambda `$, $`ϵ`$ and $`R`$.
###### Observation \newitem
The number of $`u\mathrm{\Gamma }`$ of word length $`n`$ such that $`\varphi _u(z_0)`$ lies in $`C_R`$, the vertical cylinder of radius $`R`$ centered about $`z_0`$, is bounded by $`An+B`$ where $`A=2\lambda N`$ and $`B=(2ϵ+1)N`$.
###### Theorem \newitem
For any $`ϵ`$, $`|G_ϵ^n|`$ grows linearly with $`n`$. More explicitly, there exist constants $`A`$ and $`B`$ such that the number of elements in $`G_ϵ^n`$ is less than or equal to $`An+B`$.
###### Demonstration Proof
For any $`u\mathrm{\Gamma }`$ let $`g_u=\mathrm{\Phi }(u)`$. There is some constant $`r_ϵ`$ such that any $`g_u`$ within $`ϵ`$ of the identity moves $`x_0`$ no more than $`r_ϵ`$. Let $`f_u`$ be the map whose conjugate under $`h`$ is $`g_u`$. By the Lemma 4.0 , $`f_u`$ must not move $`x_0`$ by more than $`L+r_ϵ`$. The map $`f_u`$ came from the projection of the image of the horizontal $`^2`$ containing $`\psi (e)`$ under $`\varphi _u`$. Since $`\psi (e)`$ projects to $`x_0`$, we conclude that $`\varphi _u`$ must send $`\psi (e)`$ to some point in the vertical cylinder of radius $`L+r_ϵ`$ centered about $`\psi (e)`$. Taking $`R=L+r_ϵ`$ in Observation 5.0 , the number of such $`u`$ of word length less than or equal to $`n`$ is bounded by $`An+B`$. ∎
###### Lemma \newitem
The kernel of $`\mathrm{\Phi }`$ is either finite or contains an element of infinite order.
###### Demonstration Proof
Passing if necessary to an index 2 subgroup, we may assume that all elements $`u`$ of $`\mathrm{ker}\mathrm{\Phi }`$ are end–preserving in the sense that the image under $`\varphi _u`$ of a sequence of points whose $``$ components tend to $`+\mathrm{}`$ (resp. $`\mathrm{}`$) also have $``$ components which tend to $`+\mathrm{}`$ (resp. $`\mathrm{}`$). Our goal will be to show that if an element $`u\mathrm{ker}\mathrm{\Phi }`$ has finite order, then $`\psi (u)`$ and $`\psi (e)`$ are within a vertical distance of $`\lambda c+ϵ`$ of each other, where $`c`$ is the constant of Proposition 3.0. This would imply that there are only finitely many elements of finite order, since the image under $`\mathrm{\Psi }`$ of any element of finite order in the kernel would have to be in the vertical cylinder $`C_0`$ of radius $`L`$ centered at $`z_0`$ bounded above and below by $`\lambda (cϵ/2+c)ϵ+2(\lambda \delta +ϵ)`$. But by Observation 5.0, the number of such points is finite.
Say $`u`$ has finite order $`n`$ and the vertical distance between $`\psi (u)`$ and $`\psi (e)`$ is greater than $`\lambda c+ϵ`$. We assume without loss of generality that $`\psi (u)`$ is above $`\psi (e)`$. Then for some $`k`$, $`\psi (u^k)`$ is above $`\psi (u^{k+1})`$. Draw a geodesic segment between these two points and then continue it vertically upward from $`\psi (u^k)`$ and vertically downward from $`\psi (u^{k+1})`$. Call this curve $`\gamma `$. Under $`\varphi _{u^{nk}}`$, $`\psi (u^k)`$ is sent to $`\psi (e)`$ and $`\psi (u^{k+1})`$ is sent to $`\psi (u)`$. Since we are assuming $`\varphi _u`$ is end–preserving, it follows that there is some point $`p`$ on $`\gamma `$ above $`\psi (u^k)`$, such that $`\varphi _u(p)`$ is within $`ϵ/2`$ of being at the same height as $`\psi (u)`$. From the previous section we know that this means that $`p`$ and $`\psi (u^{k+1})`$ must be no farther than $`cϵ/2+c`$ apart vertically. Thus,
$$d(\psi (u^{k+1}),\psi (u^k))d(\psi (u^{k+1}),p)cϵ/2+c.$$
Therefore,
$`d(\psi (e),\psi (u))`$ $`d(\varphi _{u^{nk}}\psi (u^{k+1}),\varphi _{u^{nk}}\psi (u^k))+2(\lambda \delta +ϵ)`$
$`\lambda d(\psi (u^{k+1}),\psi (u^k))+ϵ+2(\lambda \delta +ϵ)`$
$`\lambda (cϵ/2+c)ϵ+2(\lambda \delta +ϵ).`$
###### Theorem \newitem
The kernel is either finite or quasi-isometric to $``$.
###### Demonstration Proof
Suppose the kernel is infinite. Let $`u\mathrm{ker}\mathrm{\Phi }`$ be an element of infinite order, whose existence is guaranteed by the previous Lemma. We may assume without loss of generality that $`\varphi _u`$ is end–preserving and that $`u`$ was in our generating set for $`\mathrm{\Gamma }`$ to begin with. Since $`u`$ is of infinite order, there are infinitely many $`\psi (u^i)`$. Since all $`\psi (u^i)`$ are in $`C_0`$, the vertical distance between $`\psi (u^n)`$ and $`z_0`$ must be greater than $`c`$, for $`n`$ sufficiently large. Let $`u^n=v`$. Without loss of generality we may assume that $`\psi (v)`$ is above $`z_0`$. By an argument similar to that of the previous Lemma, $`\psi (v^{k+1})`$ is above $`\psi (v^k)`$. Let $`d(\psi (e),\psi (v))=h`$. Then for all $`k`$, $`d(\psi (v^{k+1}),\psi (v^k))\lambda h+ϵ`$. Thus every point in $`C_0`$ is within $`\frac{1}{2}(\lambda h+ϵ)+L`$ of one of the $`\psi (v^k)`$.
Say $`w\mathrm{ker}\mathrm{\Phi }`$. Then $`\psi (w)`$ is within $`\frac{1}{2}(\lambda h+ϵ)+L`$ of some $`\psi (v^k)`$. Thus,
$`d(z_0,\psi (v^kw))`$ $`\lambda d(\varphi _{v^k}(z_0),\varphi _{v^k}(\psi (v^kw))+ϵ`$
$`=\lambda d(\psi (v^k),\psi (w))+ϵ`$
$`\lambda ({\displaystyle \frac{1}{2}}(\lambda +ϵ)+L)+ϵ.`$
There are only finitely many such points so $`v`$ is of finite index in $`\mathrm{ker}\mathrm{\Phi }`$. Thus $`\mathrm{ker}\mathrm{\Phi }`$ is quasi-isometric to $``$ which in turn is quasi-isometric to $``$. ∎
###### Lemma \newitem
The group $`G`$ is a non-elementary subgroup of $`PSL(2,)`$.
###### Demonstration Proof
By construction $`G`$ does not fix a point or preserve an axis in $`^2`$. Say $`G`$ fixes a point on $`^2`$. Then $`G`$ would be solvable. Since $`\mathrm{ker}\mathrm{\Phi }`$ is either finite or a finite extension of $``$, $`G`$ solvable would imply that $`\mathrm{\Gamma }`$ was amenable. But $`\mathrm{\Gamma }`$ cannot be amenable since it is quasi-isometric to $`^2\times `$. ∎
###### Theorem \newitem
The group $`G`$ is a non-elementary subgroup of $`PSL(2,)`$ with the property that $`|G_ϵ^n|`$ grows linearly.
### Non-discrete, finitely generated, non-elementary subgoups of $`PSL(2,)`$ have exponential semilocal growth
In order to show that $`|H_ϵ^n|`$ grows quickly we need a tool for constructing small elements and another to make sure we can construct enough. The Zassenhaus Lemma will perform the first task while Tits’s Theorem will do the second. We need the following Lemmas before we can apply the results.
###### Lemma \newitem
$`G`$ is either discrete or its closure is all of $`PSL(2,)`$.
###### Demonstration Proof
By construction $`G`$ does not fix any point in $`^2`$ or its boundary, and also does not preserve any axis. It is a well-known fact (see for example , Theorem 4.4.7) that the only closed subgroups of $`PSL(2,)`$ either fix one of the above, are discrete, or are all of $`PSL(2,)`$. ∎
Let $`H`$ be a subgroup of $`PSL(2,)`$ that is not discrete and whose closure is all of $`PSL(2,)`$.
###### Lemma \newitem
$`H`$ is not virtually solvable.
###### Demonstration Proof
The Lemma follows from the fact that a virtually solvable subgroup has virtually solvable closure. ∎
###### Lemma \newitem
(Zassenhaus ) There is a constant $`ϵ_0`$ such that for any $`rϵ_0`$, if $`f`$ and $`g`$ are within distance $`r`$ from the identity, then $`[f,g]`$ is also.
A proof of this Lemma may be found in Raghunathan’s book .
We will use this Lemma to construct more small elements from a few small elements. But we need a way to check that we are indeed constructing new and different elements by conjugating. Tits’s alternative says that any subgroup of a linear group is either virtually solvable or contains a free group on two generators. We need a little more; we need a free group generated by two small generators in the Zassenhaus sense.
###### Note \newitem
Unknown to me at the time, Carrière and Ghys had previously proved the existence of such a free group along similar lines.
We will construct these elements using the following fact.
###### Proposition \newitem
Any finitely generated nonelementary subgroup $`H`$ of the group $`PSL(2,)`$ that is not discrete must contain an elliptic element of infinite order.
###### Demonstration Proof
By Selberg’s Lemma, $`H`$ contains a finite index normal subgroup $`N`$ containing no non-trivial elements of finite order, from which it follows that $`H`$ has no finite order elements of order greater than the index of $`N`$ in $`H`$. But $`H`$ is dense in $`PSL(2,)`$, so there must be elements of $`H`$ which get close to finite order elliptics of higher orders. The only possibilities are elliptics of infinite order, so $`H`$ must contain at least one. ∎
Let $`\alpha `$ be an infinite order elliptic element lying in $`H`$ whose existence is guaranteed by the previous Proposition. The eigenvalues of $`\alpha `$ all have norm 1, but are not roots of unity, so they have infinite order in $`k^{}`$, the subfield of $``$ generated by the matrix entries and eigenvalues of the generators for $`H`$. Thus we may use the following Lemma of Tits \[32, 4.1\], taking $`t`$ to be one of $`\alpha `$’s eigenvalues.
###### Lemma \newitem
Let $`k`$ be a finitely generated field and let $`tk`$ be an element of infinite order. Then there exists a locally compact field $`k^{}`$ endowed with an absolute value $`\omega `$ and a homomorphism $`\sigma :kk^{}`$ such that $`\omega (\sigma (t))1`$.
Tits finds subsets $`S`$ of $`H`$ and $`𝒰`$ of $`SL(2,)\times SL(2,)`$ with the property that for any $`s_1,s_2S`$ with $`s_1s_2`$ and $`(s_1,s_2)𝒰`$ there exists a positive power $`m`$ such that $`s_1^m`$ and $`s_2^m`$ generate a free group. He constructs these sets as follows.
The set $`S`$ consists of elements $`sH`$ with eigenvalues $`\lambda _1`$, $`\lambda _2`$ such that $`\omega (\sigma (\lambda _i))1`$ for $`i=1,2`$. Let $`𝒰`$ be the set of $`(x,y)SL(2,)\times SL(2,)`$ such that $`x`$ and $`y`$ are semisimple with distinct eigenvalues and that, for any eigenvectors $`v`$ and $`w`$ of $`x`$ and $`y`$ respectively, $`w^{}(v)0`$.
Our infinite order elliptic $`\alpha `$ clearly lies in $`S`$. Since the closure of $`H`$ is all of $`PSL(2,)`$, $`H`$ must contain a hyperbolic element. Conjugate $`\alpha `$ by this hyperbolic element to get an elliptic $`\beta `$ with a different fixed point. $`\beta `$ is also semi-simple and has the same eigenvalues as $`\alpha `$, so $`\beta `$ is also in $`S`$. A straightforward calculation shows that the eigenvectors for $`\alpha `$ and $`\beta `$ satisfy the necessary condition, so $`(\alpha ,\beta )𝒰`$. We can now prove:
###### Proposition \newitem
$`H`$ contains two small elements $`a`$ and $`b`$ that generate a free group.
###### Demonstration Proof
Since $`(\alpha ,\beta )`$ lies in $`𝒰`$, there is some integer $`m`$ such that $`\alpha ^m`$ and $`\beta ^m`$ generate a free group. Furthermore since $`\alpha ^m`$ and $`\beta ^m`$ are infinite order elliptics, there exist integers $`k`$ and $`l`$ such that $`\alpha ^{km}`$ and $`\beta ^{lm}`$ are $`ϵ_0`$-close to the identity.
Since a subgroup of a free group is free, $`a=\alpha ^{km}`$ and $`b=\beta ^{lm}`$ are small elements generating a free group. ∎
Our goal now is to use $`a`$ and $`b`$ to find a large number of elements near the identity. Let $`W_i`$ be sets of words in $`a`$ and $`b`$ defined inductively as follows. Let $`W_0=\{a,b,a^1,b^1\}`$. Let $`W_i=\{[x,y]|x,yW_{i1},xy,xy^1\}`$. Let $`W=W_i`$. Let $`C_i`$ be the image of $`W_i`$ in the free group $`a,b`$. Set $`c_i=|C_i|`$.
###### Lemma \newitem
Distinct words in $`W`$ have distinct images in the group $`a,b`$.
###### Demonstration Proof
The idea is to show that any word in $`W`$ can be canonically reconstructed from the reduced word representing the same element in the group $`a,b`$. In order to describe the reconstruction, it is helpful to describe a step-by-step reduction of a word in $`W`$.
This paragraph explains how, given a word $`wW_i`$, words $`w_{i1},w_{i2},\mathrm{},w_0`$ are inductively defined, where $`w_k`$ is a partial reduction of $`w_{k+1}`$. Note the indices decrease as the induction proceeds. View $`w`$ as a word consisting of 4 blocks, each of length $`4^{i1}`$. By construction no two consecutive blocks are inverses of each other. Set $`w_{i1}=w`$. Given $`w_j`$, for $`j1`$, we construct $`w_{j1}`$ by viewing $`w_j`$ as being made up of blocks of length $`4^{j1}`$ and canceling blocks as follows. Scan $`w_j`$ from left to right until a block followed by its inverse appears. Cancel these two blocks and then continue scanning from left to right, starting with the block which followed those just cancelled, until another block is followed by its inverse. Cancel these and continue this process until the end of the word is reached. Call the resulting word $`w_{j1}`$.
Conceivably $`w_{j1}`$ contains consecutive blocks of length $`4^{j1}`$ which are inverses of each other. The next step is to show that there are not any such pairs. We proceed by induction with the indices decreasing. As noted above $`w_{i1}`$ doesn’t contain any consecutive blocks of length $`4^{i1}`$ which are inverses of each other. By induction assume that $`w_j`$ contains no consecutive blocks of length $`4^j`$ which are inverses of each other. We show that $`w_{j1}`$ contains no consecutive blocks of length $`4^{j1}`$ which are inverses of each other. If there were consecutive blocks of length $`4^{j1}`$ which were inverses of each other, they would have to have come from a sequence in $`w_j`$ of the form $`x_ix_jx_j^1x_i^1`$ where $`x_i`$ and $`x_j`$ are in $`W_{j1}`$. Since $`w_j`$ was gotten by canceling only blocks of length $`4^i`$ or longer this sequence must be part of a sequence of the form $`x_i^1x_j^1x_ix_jx_j^1x_i^1x_jx_i`$, but such a sequence cannot occur since by induction we are assuming that $`w_j`$ contains no two consecutive blocks which are inverses of each other. Thus $`w_{j1}`$ contains no consecutive blocks of length $`4^{j1}`$ which are inverses of each other. In particular this argument shows that $`w_0`$ is a reduced word, since it can contain no consecutive blocks of length $`4^0`$ which are inverses of each other.
Let $`R_0`$ be the set of all words gotten by completely reducing some word in $`W`$. Given a word $`w_0`$ in $`R_0`$, we wish to canonically reconstruct the element of $`W`$ it came from. A word $`w_0R_0`$ is either in $`W_0`$ or it came from reducing a word $`v_1`$ with the property that it is made up of blocks of length 4, where each block is an element of $`W_1`$, and no consecutive blocks are inverses of each other. We will say a word $`v`$ has property $`()`$ if it reduces to $`w_0`$ and is made up of blocks of length 4, where each block is an element of $`W_1`$, and no consecutive blocks are inverses of each other.
This paragraph is devoted to showing that there is only one possible word satisfying $`()`$. First we must show that any possible word satisfying $`()`$ must start with the same three letters as $`v_1`$ does. The word $`v_1`$ must start with some sequence of the form $`x_1x_2x_1^1x_2^1`$ for some $`x_iW_0`$. If the second block starts with $`x_2`$ then when reducing to $`w_0`$ the $`x_2^1x_2`$ cancel but no other cancellation takes place between these two blocks since they aren’t allowed to be inverses of each other. Similarly the last and first letter of consecutive blocks may cancel but nothing more, so in particular nothing ever cancels with the first three letters of $`v_1`$. Thus $`w_0`$ and $`v_1`$ must start with the same three letters. Moreover any word satisfying $`()`$ must start with the same four letters as $`v_1`$. If $`w_0`$ and $`v_1`$ begin with the same four letters then let $`w_0^{}=w_0`$. Otherwise let $`w_0^{}`$ be the word gotten by inserting the fourth letter of $`v_1`$ followed by its inverse into $`w_0`$ between its third and fourth letter. We have just shown that the expansion $`w_0^{}`$ of $`w_0`$ is independent of the choice of $`v_1`$. The same reasoning shows that any two elements satisfying $`()`$ and beginning with the same $`4n`$ letters which reduce to $`w_0`$ must agree on the first $`4(n+1)`$ letters. So by induction there is only one $`w_1`$ satisfying $`()`$ which reduces to $`w_0`$.
Now say $`w_i`$ has been reconstructed. Then $`w_i`$ is either in $`W_i`$ or it comes from reducing a word $`w_{i+1}`$ made up of blocks of commutators of elements of $`W_i`$ such that no two consecutive blocks are inverses of each other. By the same reasoning we used in the case $`i=0`$, only one such word exists. Thus by induction we may canonically reconstruct $`wW`$ from $`w_0`$. Thus no two distinct words in $`W`$ can represent the same element of $`a,b`$.
We want to count the elements in $`W`$ and see how the number grows with the word length. Note $`c_i=c_{i1}(c_{i1}2)\frac{1}{2}c_{i1}^2.`$
###### Observation \newitem
$`c_i2^{2^i}.`$
###### Demonstration Proof
We will prove that $`c_i2^{2^i+1}`$, from which the observation follows. For $`n=1`$, we have $`c_1=8=2^{2^1+1}`$. Assume $`c_{i1}2^{2^{i1}+1}`$. Then,
$`c_i`$ $`{\displaystyle \frac{1}{2}}(2^{2^{i1}+1})^2`$
$`={\displaystyle \frac{1}{2}}(2^{2^i+2})`$
$`=2^{2^i+1}.`$
###### Theorem \newitem
For any finitely generated subgroup $`H`$ of $`PSL(2,)`$ which is not discrete or elementary, $`|H_{ϵ_0}^n|`$ grows faster than $`f(n)=2^{\frac{\sqrt{n}}{4}}`$.
###### Demonstration Proof
Elements in $`C_i`$ have word length no more than $`4^i`$. So $`|H_{ϵ_0}^{4^i}|2^{2^i}`$. Given $`n`$, let $`j`$ be such that $`4^jn<4^{j+2}`$.
$$f(n)=2^{\frac{\sqrt{n}}{4}}<2^{2^j}|H_{ϵ_0}^{4^j}||H_{ϵ_0}^n|.$$
###### Theorem \newitem
G must be discrete.
###### Demonstration Proof
Apply Theorem 5.0 together with Theorem 5.0 taking $`ϵ=ϵ_0`$. ∎
## 6. The kernel of $`\mathrm{\Phi }`$ is quasi-isometric to $``$
To complete our proof we need only show that the kernel of $`\mathrm{\Phi }:\mathrm{\Gamma }G`$ is infinite and therefore, by Lemma 5.0, quasi-isometric to $``$.
###### Lemma \newitem
The kernel of $`\mathrm{\Phi }`$ is infinite.
###### Demonstration Proof
Since $`G`$ is discrete there are only finitely many elements $`g_u`$ that move $`x_0`$ less than any given bounded amount. For any $`u\mathrm{\Gamma }`$ such that $`\varphi _u`$ moves $`z_0`$ less than or equal to $`M_1`$ horizontally, the corresponding $`g_u`$ can move $`x_0`$ no more than $`M_1+L`$. But there are infinitely many such $`u`$, since any point in $`^2\times `$ is within $`M_1`$ of some orbit point of $`z_0`$. So some element $`gG`$ must be $`g_u`$ for infinitely many $`u`$, say $`u_1,u_2,u_3,\mathrm{}`$. But then all of the $`g_{u^1u_i}`$ must be the identity, hence $`\mathrm{ker}\mathrm{\Phi }`$ is infinite. ∎
Thus we have shown:
###### Theorem \newitem
Any group $`\mathrm{\Gamma }`$ quasi-isometric to $`^2\times `$ is an extension of a finite extension of a cocompact Fuchsian group by a virtually infinite cyclic group. In other words there is an exact sequence
$$0\mathrm{@}>>>A\mathrm{@}>>>\mathrm{\Gamma }\mathrm{@}>>>G\mathrm{@}>>>0$$
where $`A`$ is virtually infinite cyclic and $`G`$ is a finite extension of a cocompact Fuchsian group.
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# Quasi-normal modes of charged, dilaton black holes
## I INTRODUCTION
The study of the effects produced by the coupling between gravity and scalar fields is a rather hot issue, since a definite prediction of superstring theory is just the existence of a scalar field, namely the *dilaton*. For instance, this coupling could have played an important role during the early phases of the life of the Universe. Indeed, it has been shown that a possible cosmological implication of the dilaton is the so-called “pre-big-bang” scenario , characterized by a dilaton-driven inflationary period, that should have left its peculiar fingerprints on the cosmic gravitational wave background . On the experimental side, recent investigations on the interaction of scalar waves with gravitational detectors have shown that a scalar component of the gravitational radiation should excite the monopole mode of a resonant spherical antenna and also give rise to specific correlations between the signals revealed by a resonant sphere and an interferometer .
Among all possible astrophysical sources of gravitational waves, black holes should have the most typical and in principle recognizable frequency spectrum. According to the Einstein theory of gravity, a black hole which dynamically interacts with its surroundings emits wave bursts having shape and intensity which initially depend on the features of the external perturbations. As a late-time effect, however, it generates gravitational wave trains having characteristic frequencies which are *independent* of the initial perturbations, the so-called *quasi normal modes* . And we know from the theory that these frequencies can only depend on the three parameters characterizing a black hole, namely the mass, the charge and the angular momentum.
The question whether a black hole can support an additional scalar degree of freedom and still generate a regular, asymptotically flat, spacetime free of naked singularities or any other kind of pathological behaviour has been widely investigated in the literature. Having in mind the “no-hair theorem” program, Bekenstein in 1972 showed that black holes cannot support a *minimally coupled* scalar field, if the solution is required to be asymptotically flat. Subsequently, he developed a procedure to generate black hole solutions surrounded by a *conformal* scalar field , but they proved to be *unstable* against radial perturbations . In 1990 Ferrari and Xanthopoulos , using a Kaluza-Klein approach, derived equations describing a gravitational field coupled with a massive scalar field, obtained the conformal structure of the metric, and discussed several kind of possible couplings. Yet, no definite answer was given to the question whether a satisfactory black hole can exist surrounded by a massive scalar field.
The search for regular exact solutions describing black holes endowed with a scalar field got a formidable boost when it became clear that, as long as the curvature is small, all vacuum solutions of general relativity are also approximate solutions of string theory. As a matter of fact, charged black hole solutions coupled to a scalar field, viz a dilaton in this context, can be obtained by the action (we use geometric units throughout)
$$S=d^4x\sqrt{g}[R2(\varphi )^2+e^{2a\varphi }F^2],$$
(1)
which, by variations, gives the equations:
$$_\mu (e^{2a\varphi }F^{\mu \nu })=0,$$
(2)
$$^2\varphi \frac{1}{2}e^{2a\varphi }F^2=0,$$
(3)
$$R_{\mu \nu }=2_\mu \varphi _\nu \varphi 2e^{2a\varphi }F_{\mu \rho }F_\nu ^\rho +\frac{1}{2}g_{\mu \nu }e^{2a\varphi }F^2,$$
(4)
where $`F^2=F_{\mu \nu }F^{\mu \nu }`$ is the first Maxwell invariant and $`a`$ is a non-negative real constant regulating the strength of the coupling between the dilaton and the Maxwell field. This theory has a number of interesting limiting cases. The limit $`a=0`$ corresponds to the ordinary Einstein-Maxwell theory plus a Klein-Gordon scalar field with zero mass. The case $`a=\sqrt{3}`$ is the four-dimensional reduction of Kaluza-Klein theory. Finally, for $`a=1`$ the action (1) describes the tree-level low energy limit of superstring theory in the so-called *Einstein frame*. A large class of black hole solutions of this theory in an arbitrary number of dimensions has been found by Gibbons and Maeda . Their results were next specialized to four dimensions by Garfinkle, Horowitz and Strominger (GHS) who studied a solution describing a spherically symmetric, charged, dilaton black hole. In the GHS solution the electric charge and the dilaton are not independent parameters: when the charge is set equal to zero the dilaton also disappears, and the solution reduces to that of a Schwarzschild black hole. In a sense, this is a consequence of the “no-hair theorems” which limits the number of free parameters of a black hole to three.
The equations satisfied by small perturbations of this solution have been derived and studied by Holzhey and Wilczek in the line of the general treatement given by Chandrasekar . They first showed that it is possible to reduce the perturbed equations to five decoupled wave equations with potential barriers, and, as a by-product of their analysis, they argued that the GHS solution is stable under external small perturbations.
In this paper we take a further step in the study the properties of the coupled emission of electromagnetic, scalar and gravitational radiation by such black holes, computing the *quasi-normal mode frequencies* of the GHS solution. Comparing the spectrum of the latter with those of Schwarzschild and Reissner-Nordström, we show that the presence of the scalar field breaks the relevant feature of the *isospectrality* of the axial and polar perturbations.
The paper is organized as follows. In Section 2 we summarize the features of the GHS solution. In Sections 3 and 4 the equations governing the axial and polar perturbations are discussed separately. We shall not derive the axial equations, since there is nothing to add to the derivation made in Ref. . On the other hand, the polar equations will be discussed in greater detail. Some misprints appearing in Ref. are corrected and the explicit expression of the matrix whose eigenvalues are the potentials governing polar perturbations is given in the Appendix. In Section 5 the quasi-normal frequencies of the GHS black hole obtained with an extended WKB approch are calculated for different values of the parameters. Finally, in Section 6 the results obtained are discussed and the main qualitative differences between the spectra of the dilaton and the Reissner Nordström black holes are pointed out.
## II THE EXACT SOLUTION
The exact solution of equations (2)-(4), with $`a=1`$, describing the charged dilaton black hole we study in this paper is
$$ds^2=\left(1\frac{2M}{r}\right)dt^2\left(1\frac{2M}{r}\right)^1dr^2r\left(r\frac{Q_e^2}{M}\right)\left[d\theta ^2+sin^2\theta d\phi ^2\right]$$
(5)
where $`M`$ and $`Q_e`$ are, respectively, the black hole mass and electric charge. $`F_{tr}=Q_e/r^2`$ is the only non-vanishing component of the electromagnetic tensor, and the scalar field is related to the electric charge by the following equation
$$e^{2\varphi }=\left(1\frac{Q_e^2}{Mr}\right),$$
(6)
which shows that $`\varphi `$ vanishes at radial infinity. Since the dilaton and the electric charge are coupled through Eq. (6), the Reissner Nördstrom solution cannot be obtained as a limiting case of the metric (5). On the other hand, the Schwarzschild solution is recovered from Eq. (5) by setting $`Q_e=0`$.
It should be noted that the usual relation between radius and area of the spheres $`t=const`$, $`r=const`$ is obtained in terms of the modified radial variable
$$\stackrel{~}{r}=\sqrt{r\left(r\frac{Q_e^2}{M}\right)},$$
and not of $`r`$, by the usual relation $`A(r)=\mathrm{\hspace{0.17em}4}\pi \stackrel{~}{r}^2.`$
The metric (5) appears to be singular in $`r=0`$, $`r=Q_e^2/M`$ and $`r=2M`$. The surface $`r=2M`$ is an event horizon, whereas on $`r=Q_e^2/M`$ the curvature scalar
$$R=2g^{rr}(\varphi _{,r})^2=\frac{1}{2}\left(1\frac{2M}{r}\right)\frac{Q_e^2}{r(MrQ_e^2)},$$
diverges, showing that there is a curvature singularity. In $`r=Q_e^2/M`$, the radial coordinate $`\stackrel{~}{r}`$ vanishes and loses its meaning for $`r<Q_e^2/M`$. Thus, a physical observer who crosses the event horizon terminates its jurney on the curvature singularity.
In order to avoid naked singularities, in what follows we shall assume that
$$2M^2>Q_e^2.$$
As shown in Ref. (to be referred to hereafter as MT), the study of the perturbations of the metric (5), can be restricted, without loss of generality, to axisymmetric perturbations only. The appropriate metric in this case is
$$ds^2=e^{2\nu }(dt)^2e^{2\psi }(d\varphi q_2dx^2q_3dx^3\omega dt)^2e^{2\mu _2}(dx^2)^2e^{2\mu _3}(dx^3)^2,$$
(7)
where, in the unperturbed state,
$`e^{2\nu }=`$ $`e^{2\mu _2}`$ $`=\left(1{\displaystyle \frac{2M}{r}}\right),e^{2\psi }=\stackrel{~}{r}^2sin^2\theta `$ (8)
$`e^{2\mu _3}`$ $`=\stackrel{~}{r}^2,\omega =q_2=q_3=0,`$ (9)
and $`(0,1,2,3)`$ stand for $`(t,\phi ,r,\vartheta )`$. We shall assume that, as a consequence of a generic perturbation, the metric functions, the electromagnetic quantities and the scalar field will experience small changes with respect to their unperturbed values
$`\nu \nu +\delta \nu ,`$ $`\mu _2\mu _2+\delta \mu _2,`$ $`\psi \psi +\delta \psi ,`$ (10)
$`\mu _3\mu _3+\delta \mu _3,`$ $`\omega \delta \omega ,`$ $`q_2\delta q_2,`$ (11)
$`q_3\delta q_3,`$ $`F_{\mu \nu }F_{\mu \nu }+\delta F_{\mu \nu },`$ $`\varphi \varphi +\delta \varphi .`$ (12)
Since the perturbation is assumed to be axisymmetric, all perturbed quantities depend on $`t`$, $`r`$ and $`\theta `$ only. As in MT, it is convenient to project Einstein’s and Maxwell’s equations onto an orthonormal tetrad frame, and assume that all perturbed functions have the time dependence $`e^{i\sigma t}`$. The separation of variables is accomplished by expanding all perturbed tensors in tensorial spherical harmonics. These harmonics belong to two different classes, depending on the way they behave under the angular transformation $`\theta \pi \theta `$ and $`\varphi \pi +\varphi .`$ Those that transform like $`(1)^{(\mathrm{}+1)}`$ are termed axial, those that transform like $`(1)^{(\mathrm{})}`$ are termed polar. The perturbed equations split into two decoupled sets corresponding to a different parity.
## III THE AXIAL EQUATIONS
The axial equations for $`\mathrm{}2`$ are obtained by perturbing the {12} and {13}-components of the Einstein equations. The separation of variables, can be accomplished by putting
$`\stackrel{~}{r}^2e^{2\nu }\mathrm{sin}^3\theta \left[q_{2,3}(t,r,\vartheta )q_{3,2}(t,r,\vartheta )\right]=Q_{\mathrm{}}(r,\sigma )C_{\mathrm{}+2}^{3/2}(\vartheta )e^{i\sigma t},`$
and
$`F_{01}(t,r,\vartheta )\mathrm{sin}\vartheta =3B_{\mathrm{}}(r,\sigma )C_{\mathrm{}+1}^{1/2}(\vartheta )e^{i\sigma t},`$
where $`C_{l+1}^{1/2}(\vartheta )`$ are the Gegenbauer polynomials. As shown in Ref. , the axial equations can be cast in the following form
$$\left(\frac{d^2}{dr_{}^{2}}+\sigma ^2\right)\left(\begin{array}{c}H_1\mathrm{}\\ H_2\mathrm{}\end{array}\right)=𝐁\left(\begin{array}{c}H_1\mathrm{}\\ H_2\mathrm{}\end{array}\right)$$
(13)
where
$$𝐁=\frac{e^{2\nu }}{\stackrel{~}{r}^2}\left[\left(\mu ^2+2+\frac{Q_e^2}{r^2}+\frac{3Q_e^4}{4M^2\stackrel{~}{r}^2}e^{2\nu }\right)\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)+\frac{1}{r}\left(\begin{array}{cc}Q_e^2/M& 2\mu Q_e\\ 2\mu Q_e& 6M\end{array}\right)\right],$$
(14)
$`r_{}`$ is the tortoise coordinate defined by the equation
$$dr_{}=e^{2\nu }dr,$$
(15)
and $`\mu ^2=(\mathrm{}1)(\mathrm{}+2)`$. The radial functions $`H_1`$ and $`H_2`$ (from now on we omit the index $`\mathrm{}`$) are related to the perturbed metric and electromagnetic functions by the following relations
$$Q(r,\sigma )=\stackrel{~}{r}H_2(r,\sigma )re^\nu B(r,\sigma )=\frac{H_1(r,\sigma )}{2\mu }.$$
(16)
The right hand side of Eq. (13) can be obviously diagonalized by a linear r-independent transformation, in the form
$`Z_1^a=`$ $`_1H_1+_2H_2,`$ (17)
$`Z_2^a=`$ $`_2H_1_1H_2,`$ (18)
where
$$_1=\frac{Q_e^2}{2M}+3M+\sqrt{\frac{Q_e^4}{4M^2}+9M^2+Q_e^2(3+4\mu ^2)}\text{and}_2=2\mu Q_e.$$
(19)
Then the system (13) decouples in two wave equations:
$$\left(\frac{d^2}{dr_{}^{2}}+\sigma ^2\right)Z_i^a=V_i^aZ_i^a(i=1,2),$$
(20)
where the explicit form of the effective potentials is
$$V_{1,2}^a=\frac{e^{2\nu }}{r\stackrel{~}{r}^2}\left[(\mu ^2+2)r+\frac{Q_e^2}{r}+\frac{3Q_e^4r}{4M^2\stackrel{~}{r}^2}e^{2\nu }+\frac{Q_e^2}{2M}3M\pm \sqrt{\frac{Q_e^4}{4M^2}+9M^2+Q_e^2(3+4\mu ^2)}\right].$$
(21)
Equation (19) shows that when $`Q_e`$ vanishes, $`_2=0`$; in this case $`Z_2^a`$ reduces to the gravitational perturbation and $`V_2^a`$ to the Regge-Wheeler potential, whereas $`Z_1^a`$ reduces to the pure electromagnetic perturbation of a Schwarzschild black hole, which is known to be independent of gravitational contributions. If the electromagnetic charge does not vanish, the gravitational and electromagnetic perturbations are coupled; a gravitational wave incident on the potential barriers induces the emission of electromagnetic radiation and viceversa. On the other hand, the dilaton is not dynamically coupled with the axial perturbations. Its effect is that of shaping the effective potentials together with the electric field. In a similar manner, the energy density and the pressure of the matter composing a perturbed star determine the potential barrier of the axial perturbations without being dynamically coupled to the perturbed gravitational field.
The $`\mathrm{}=2`$ potentials $`V_1^a`$ and $`V_2^a`$ are plotted in Fig. LABEL:fig:1 versus the rescaled tortoise coordinate $`r_{}/M`$ for different values of the electric charge. We see that they always tend to zero at radial infinity and at the black hole horizon, except when the charge assumes its extremal value $`Q_e^2=2M^2.`$ In this case they take the form of a step, reflecting all waves whose square frequency is lower than their limiting value on the horizon.
We also see that for non extremal black holes the maximum of the barrier moves towards the horizon as the electric charge increases.
## IV THE POLAR EQUATIONS
The separation of variables for polar equations is accomplished by requiring that the perturbed functions have the angular dependence deriving from the expansion in tensor spherical harmonics:
$`\delta \nu `$ $`=`$ $`N_{\mathrm{}}(r)P_{\mathrm{}}(\theta )e^{i\sigma t},`$ (22)
$`\delta \mu _2`$ $`=`$ $`L_{\mathrm{}}(r)P_{\mathrm{}}(\theta )e^{i\sigma t},`$ (23)
$`\delta \psi `$ $`=`$ $`\left[T_{\mathrm{}}(r)P_{\mathrm{}}(\theta )+2\mu ^2X_{\mathrm{}}(r)P_{\mathrm{}}(\theta )_{,\theta }\mathrm{cot}\theta \right]e^{i\sigma t},`$ (24)
$`\delta \mu _3`$ $`=`$ $`\left[T_{\mathrm{}}(r)P_{\mathrm{}}(\theta )+2\mu ^2X_{\mathrm{}}(r)P_{\mathrm{}}(\theta )_{,\theta ,\theta }\right]e^{i\sigma t},`$ (25)
$`\delta F_{02}`$ $`=`$ $`{\displaystyle \frac{e^{2\nu }\stackrel{~}{r}^2}{2Q_e}}B_{02\mathrm{}}(r)P_{\mathrm{}}(\theta )e^{i\sigma t}`$ (26)
$`\delta F_{03}`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{r}}{2Q_e}}B_{03\mathrm{}}(r)P_{\mathrm{}}(\theta )_{,\theta }e^{i\sigma t},`$ (27)
$`\delta F_{23}`$ $`=`$ $`{\displaystyle \frac{ie^\nu \sigma \stackrel{~}{r}}{2Q_e}}B_{23\mathrm{}}(r)P_{\mathrm{}}(\theta )_{,\theta }e^{i\sigma t},`$ (28)
$`\delta \varphi `$ $`=`$ $`\mathrm{\Phi }P_{\mathrm{}}(\theta )e^{i\sigma t},`$ (29)
where $`P_{\mathrm{}}(\theta )`$ are the Legendre polinomials. In the following, we shall omit the index $`\mathrm{}`$ in all radial functions. Among the eight radial functions, $`N`$, $`L`$, $`T`$, $`X`$, $`B_{02}`$, $`B_{03}`$, $`B_{23}`$ and $`\mathrm{\Phi }`$, the function $`T`$ can be eliminated from the perturbed equations by making use of the relation that follows from the $`\{03\}`$-component of the perturbed Einstein equations
$$B_{23}TL+2\mu ^2X=0.$$
(30)
From the polar components of the Maxwell’s equations (see MT, Chap. 5 Eqs. (165)-(167) for the Reissner Nördstrom case), it is easy to derive a second order equation for $`B_{23}`$:
$`\left[{\displaystyle \frac{\stackrel{~}{r}^2}{r^2}}e^{2\nu }(r^2B_{23})_{,r}\right]_{,r}+\left[\sigma ^2\stackrel{~}{r}^2e^{2\nu }(\mu ^2+2)\right]B_{23}`$ (31)
$`+{\displaystyle \frac{2Q_e^2}{r^2}}(2B_{23}3L2XN2\mathrm{\Phi })`$ $`=`$ $`0.`$ (32)
The Einstein equations for $`\delta R_{02}`$, $`\delta R_{23}`$, $`\delta G_{22}`$ and $`\delta R_{11}`$ give:
$`\left({\displaystyle \frac{d}{dr}}+(\mathrm{log}\stackrel{~}{r}e^\nu )_{,r}\right)(B_{23}LX)(\mathrm{log}\stackrel{~}{r})_{,r}L+\varphi _{,r}\mathrm{\Phi }`$ $`=`$ $`0,`$ (33)
$`(NL)_{,r}(\mathrm{log}\stackrel{~}{r}e^\nu )_{,r}L(\mathrm{log}\stackrel{~}{r}e^\nu )_{,r}N{\displaystyle \frac{2}{r}}B_{23}+2\varphi _{,r}\mathrm{\Phi }`$ $`=`$ $`0,`$ (34)
$`X_{,r,r}+2(\mathrm{log}\stackrel{~}{r}e^\nu )_{,r}X_{,r}+{\displaystyle \frac{\mu ^2e^{2\nu }}{2\stackrel{~}{r}^2}}(N+L)+\sigma ^2e^{4\nu }X`$ $`=`$ $`0,`$ (35)
$`2(\mathrm{log}\stackrel{~}{r})_{,r}N_{,r}+2(\mathrm{log}\stackrel{~}{r}e^\nu )_{,r}(B_{23}LX)_{,r}{\displaystyle \frac{\mu ^2e^{2\nu }}{\stackrel{~}{r}^2}}TB_{02}`$ (36)
$`{\displaystyle \frac{(\mu ^2+2)e^{2\nu }}{\stackrel{~}{r}^2}}N2(\mathrm{log}\stackrel{~}{r})_{,r}(\mathrm{log}\stackrel{~}{r}e^{2\nu })_{,r}L+2(\varphi _{,r})^2L`$ (37)
$`+2\sigma ^2e^{4\nu }(B_{23}LX){\displaystyle \frac{2Q_e^2e^{2\nu }}{\stackrel{~}{r}^2r^2}}\mathrm{\Phi }2\varphi _{,r}\mathrm{\Phi }_{,r}`$ $`=`$ $`0.`$ (38)
Finally, the perturbed equation for the dilaton is obtained by perturbing eq. (3):
$`{\displaystyle \frac{1}{\stackrel{~}{r}^2}}(\stackrel{~}{r}^2e^{2\nu }\mathrm{\Phi }_{,r})_{,r}+(\sigma ^2e^{2\nu }{\displaystyle \frac{(\mu ^2+2)}{\stackrel{~}{r}^2}}+{\displaystyle \frac{2Q_e^2}{\stackrel{~}{r}^2r^2}})\mathrm{\Phi }`$ (39)
$`+e^{2\nu }(N3L2X+2B_{23})_{,r}\varphi _{,r}2{\displaystyle \frac{(\stackrel{~}{r}^2e^{2\nu }\varphi _{,r})_{,r}}{\stackrel{~}{r}^2}}L+e^{2\nu }B_{02}`$ $`=`$ $`0.`$ (40)
The seventh order linear system composed by equations (31)-(40) has been shown to be reducible to three Schrödinger-like equations by using the following procedure. It is well-known that the order of a system of linear differential equations can be reduced, whenever a particular solution of that system is known. A general algorithm for deriving the particular solution which makes the reduction possible as been obtained by Xanthopulos . Holzhey and Wilczek found that the Xanthopulos solution one gets from the reduction of the polar equations of Schwarzschild and Reissner-Nordström black holes are *pure gauge* solutions i.e. gauge equivalent to the null perturbation. The general form of the metric (7) has, indeed, a gauge degree of freedom, i.e. there exists a one parameter class of coordinate transformations which leaves the form of the metric unchanged. In the case of the GHS black hole the particular *pure gauge* solution which reduces the system (31)-(40) is
$`N^{(0)}`$ $`=`$ $`\sigma ^2e^\nu \stackrel{~}{r}+e^{4\nu }\nu _{,r}(\stackrel{~}{r}e^\nu )_{,r},`$ (41)
$`L^{(0)}`$ $`=`$ $`(e^{4\nu }(\stackrel{~}{r}e^\nu )_{,r})_{,r}e^{4\nu }\nu _{,r}(\stackrel{~}{r}e^\nu )_{,r},`$ (42)
$`T^{(0)}`$ $`=`$ $`e^{4\nu }\stackrel{~}{r}^1\stackrel{~}{r}_{,r}(\stackrel{~}{r}e^\nu )_{,r},`$ (43)
$`X^{(0)}`$ $`=`$ $`{\displaystyle \frac{\mu ^2e^\nu }{2\stackrel{~}{r}}},`$ (44)
$`B_{23}^{(0)}`$ $`=`$ $`T^{(0)}+L^{(0)}2\mu ^2X^{(0)},`$ (45)
$`\mathrm{\Phi }^{(0)}`$ $`=`$ $`\varphi _{,r}e^{4\nu }(\stackrel{~}{r}e^\nu )_{,r}.`$ (46)
Following Holzhey and Wilczek, we now introduce a new variable $`S`$ replacing $`L`$
$$S=B_{23}LX,$$
(47)
and make the following substitutions:
$$N=N^{(0)}s+n,B_{23}=B_{23}^{(0)}s+\frac{b}{r\stackrel{~}{r}}X=X^{(0)}s+\frac{x}{\stackrel{~}{r}},S=S^{(0)}s,\mathrm{\Phi }=\mathrm{\Phi }^{(0)}s+\frac{p}{\stackrel{~}{r}}.$$
(48)
The system is now of order six with respect to the new variables $`n`$, $`x`$, $`s`$, $`b`$ and $`p`$ since $`s`$ does appear through its derivatives only. It is easy to show from Eq. (33) that the first derivative of $`s`$ can be written as a linear combination of $`b`$, $`x`$ and $`p`$ as follows
$$s_{,r}=\frac{1}{S^{(0)}}\left(\frac{\stackrel{~}{r}_{,r}}{\stackrel{~}{r}^2r}b\frac{\stackrel{~}{r}_{,r}}{\stackrel{~}{r}^2}x\frac{\varphi _{,r}}{\stackrel{~}{r}}p\right).$$
(49)
Similarly, using (34) and (38), we can eliminate $`n`$ in favour of $`x`$, $`b`$, $`p`$ and their first derivatives. We are now left with only three perturbation variables, $`x`$, $`b`$ and $`p`$, governed by three second order equations, namely (35), (31) and (40). In (48), $`x`$, $`b`$ and $`p`$ have been defined in such a way that their first derivative *with respect to $`r_{}`$* disappears from the equations they satisfy. Therefore the system has now the desired form
$$\left(\frac{d^2}{dr_{}^2}+\sigma ^2\right)𝐯=\mathrm{𝐀𝐯},$$
(50)
where
$$𝐯=\left(\begin{array}{c}x\\ b\\ p\end{array}\right),$$
(51)
and the components of the symmetric matrix A are complicated functions of $`r`$.
It is convenient to cast A in the following form:
$$𝐀(r)=\frac{1}{D(r)}\left[G(r)𝐓(r)+P(r)𝐈\right],$$
(52)
where the polinomials $`D(r)`$, $`G(r)`$ and $`P(r)`$ and the matrix $`𝐓(r)`$ are given in the Appendix, and $`𝐈`$ is the identity matrix. It is remarkable that the eigenvectors of $`𝐓(r)`$ and $`𝐀(r)`$ are independent of $`r`$. This means that the system can be decoupled with a basis transformation in the $`x`$-$`b`$-$`p`$ space. The three eigenvalues are the potentials of three independent Schrödinger equations:
$$\left(\frac{d^2}{dr_{}^{2}}+\sigma ^2\right)Z_i^p=V_i^pZ_i^p(i=1,\mathrm{\hspace{0.17em}2},\mathrm{\hspace{0.17em}3}).$$
(53)
It should be noted that for $`Q_e=0`$ the equations are already decoupled since the off-diagonal terms of A vanish (see Appendix) while, as already pointed out, the GHS solution becomes the Schwarzschild solution. In this case, as expected, we find the well-known potentials that rule pure gravitational, electromagnetic, and massless scalar perturbations on a Schwarzschild background, respectively.
Proving the independence of $`r`$ of the eigenvectors of $`𝐓(r)`$ directly, even by means of a symbolic procedure, turns out to be awkward. A straighforward manner is instead to expand the polinomial valued matrix as:
$$𝐓=\underset{i=0}{\overset{6}{}}r^i𝐓^{(i)},$$
(54)
where $`𝐓^{(i)}`$ are matrices independent of $`r`$. It is now very easy to show, for instance with the aid of MATHEMATICA, that the $`𝐓^{(i)}`$ commute among themselves, so that they can be diagonalized by the same linear transformation.
The eigenvalue problem for $`𝐀(r)`$ is now considerably reduced by means of (52) and (54). In fact, the three potentials can easily be obtained by fixing *first* the values of the parameters $`M`$, $`Q_e`$ and $`\mu `$, then calculating the eigenvalues of each $`𝐓^{(i)}`$, which are now real numbers, and finally summing them up as coefficients of their respective powers of $`r`$. The eigenvectors of $`𝐀(r)`$ are then obtained straighforwardly by means of (52). The three potentials obtained with the above procedure appear in a rather simple form as functions of $`r`$ only. This fact turns out to be extremely useful, since to calculate the quasinormal frequencies by means of WKB method one needs up to their sixth derivatives. The potentials $`V_i^p`$ associated to the wavefunctions $`Z_i^p`$ (see Eq. 50) are plotted in Fig. 2 for different values of the charge. As in the axial case, when the charge approaches the extremal value, all potentials assume the form of a step. When $`Q_e0,`$ all off-diagonal components of the matrix $`𝐀`$ vanish, and the system of equations (50) decouple into three wave equations for the functions $`(x,b,p)`$. From eqs. (48) it is easy to check that, since $`B_{23}^{(0)}`$ and $`\mathrm{\Phi }^{(0)}`$ vanish, the functions $`b`$ and $`\mathrm{\Phi }`$ become purely electromagnetic and scalar, respectively. Conversely, the variable $`x`$ is not purely gravitational, since it contains $`B_{23}.`$ If, however, we explicitly write the function $`x`$ in terms of $`B_{23},X`$ and $`L`$ in the first Eq. (50), the electromagnetic terms disappear by virtue of the equation satisfied by $`b`$, and we are left with the Zerilli equation.
## V THE QUASI-NORMAL MODES
We have computed the complex frequencies of the quasi-normal modes associated to the axial and polar wave equation by using a WKB approximation devised by Schutz and Will and extended to higher orders by Iyer and Will . This approach has been applied to the Schwarzchild and Reissner Nordström cases, and for the fundamental quadrupole mode $`l=2`$ agrees with other approaches within 1% both for the real and the immaginary parts of the first 3-4 modes. The agreement improves with increasing angular harmonic and decreasing mode number.
The results of our calculations are shown in table I and II, where we tabulate the real and imaginary part of the frequencies of the first five quasi normal modes associated to the axial and polar potentials, for different values of the harmonic index $`\mathrm{}`$, and of the charge $`Q_e`$.
The general behaviour of the quasi-normal frequencies for increasing values of the charge is well described in Figure 3. There we plot the real and the imaginary part of the first eigenfrequency of the $`\mathrm{}=2`$ mode, associated to the potentials $`V_2^a`$ and $`V_1^p,`$ as a function of the charge $`Q_e`$. It should be reminded that in the limit $`Q_e=0,`$ these potentials reduce to the Regge-Wheeler and to the Zerilli potentials for the axial and polar gravitational perturbations of a Schwarzschild black hole, respectively. It is well known that the axial and polar perturbations of a Schwarzschild black hole is a *isospectral*, and this is the reason why, in the limit $`Q_e=0,`$ the axial ($`V_2^a`$) and polar ($`V_1^p`$) frequencies shown in Fig. 3 converge to a unique value.
The upper value of $`Q_e/M`$ we consider in our calculations is $`Q_e/M=1.4,`$ close to the limiting value $`Q_e/M=\sqrt{2},`$ where the potential barrier becomes a step which reflects all incident waves.
For a Reissner-Nordström black hole the perturbed axial and polar equations can be decoupled in terms of four functions $`Z_1^{}`$ and $`Z_2^{}`$, respectively . In the limit $`Q_e=0,`$ the functions $`Z_2^{}`$ reduce to the Regge-Wheeler and to the Zerilli functions, whereas $`Z_1^{}`$ reduce to pure axial and polar electromagnetic functions. Moreover, the potentials governing $`Z_1^{}`$ and $`Z_1^+`$ are *isospectral*, as well as those for $`Z_2^{}`$ and $`Z_2^+`$ . In Figure 4 the same data of Figure 3 are compared with the lowest quasi-normal mode frequency of the functions $`Z_2^{}`$ of a Reissner-Nordström black hole. Since the extremal value of the charge for a Reissner-Nordström black hole is $`Q_e=1,`$ the data stop at that value.
We see that the real part of the frequency increases as a function of the charge both for a dilaton, and for a Reissner-Nordström black hole; similarly, the imaginary part increases and then decreases to a finite value as the charge approaches the limiting value. Since however since the limiting values are different, the imaginary part for Reissner-Nordström begins to decrease while the corresponding GHS values still increase.
Figures 3 and 4 clearly show that the potentials $`V_2^a`$ and $`V_1^p`$ of a dilaton black hole *are not isospectral* as they are instead when $`Q_e=0.`$ In addition, it should be reminded that the potential $`V_2^a`$ rules the equation for $`Z_2^a`$ which is a combination of electromagnetic and gravitational perturbations only (cfr. Eq. 17), whereas $`V_1^p`$ appears in the polar equation for $`Z_1^p,`$ which is a combination of electromagnetic, gravitational and scalar perturbations.
Finally, in Figure 5 we plot the real and imaginary part of the frequency of the lowest $`l=2`$ quasi-normal mode for the remaining potentials, $`V_1^a,`$ $`V_2^p`$ and $`V_3^p,`$ as functions of the charge $`Q_e.`$ When $`Q_e=0,`$ $`V_1^a`$ and $`V_2^p`$ reduce to those governing the pure electomagnetic, axial and polar perturbations of a Schwarzschild black hole, respectively, and indeed, they are isospectral in that limit. In the same limit, $`V_3^p`$ reduces to the potential of the pure scalar perturbations of a Schwarzschild black hole.
## VI CONCLUDING REMARKS
The spectrum of the quasi-normal modes of a charged, dilaton black hole is different from that of a Schwarzschild or a Reissner-Nordström black hole. For a Schwarzschild black hole the perturbations are completely described by the Regge-Wheeler and the Zerilli equations for the axial and the polar perturbations. Although the analytic form of the two potential barriers is different, they are related by a very simple equation (MT, ch. 5, & 43) which allows them to have the same reflexion and trasmission coefficients. Therefore, since the quasi-normal mode frequencies are the singularities of the scattering amplitude, it follows that the two potentials are *isospectral*. Thus, a perturbed Scharzschild black hole emits axial and polar gravitational waves at exactly the same frequencies.
For a Reissner-Nordström black hole the perturbed equations can be reduced to four wave equations, two for the axial and two for the polar perturbations, respectively. The four wave-functions $`Z_1^\pm `$ and $`Z_2^\pm ,`$ where $`+`$ stands for polar and $``$ for axial, are a linear combination of the gravitational and the electromagnetic functions belonging to the corresponding parity. It turns out that the two potentials $`V_1^+`$ and $`V_1^{}`$ are again related in such a way that they have the same reflection and trasmission coefficients, and the same is true for $`V_2^\pm .`$ Thus the coupling GW-EM is such that it preseves the isospectrality of the axial and polar perturbations. However there is an important difference with respect to the Schwarzschild case: no quasi-normal mode exists that is purely electromagnetic or gravitational, which means that the excitation of a mode will be accompanied by the simultaneous emission of both gravitational and electromagnetic waves.
For a charged black hole in a theory described by action (1) the situation is different. Let us consider the axial perturbations first. As shown in Section 3, the perturbed equations can be reduced to two wave equations, but the perturbed dilaton does not couple to the electromagnetic and gravitational fields. This is due to the fact that the dilaton is a scalar, and its axial perturbation vanishes. Consequently, the excitation of an axial mode will be accompanied only by the emission of gravitational and electromagnetic waves. However, the dilaton appears in the unperturbed metric functions that determine the shape of the potentials of the axial wave equations. Thus it affects the scattering properties of the axial potentials. It is interesting to note that the real part of the quasi-normal mode frequencies of the axial potential $`V_2^a`$ are very similar to those of the Reissner-Nordström black hole (see Fig. 4), even though the dilaton solution does not reduce to the Reisnner-Nordström solution in any limiting case.
On the other hand, the two wave equations which describe the polar perturbations of a Reissner-Nordström black hole, are replaced by three wave equations in the case of a dilaton black hole, and they couple the gravitational, electromagnetic and scalar perturbations. This occurrence breaks the symmetry between axial and polar perturbations, and makes the scattering properties of the two parities different (see Fig. 5).
We conclude that for a dilaton black hole the excitation of an axial mode induces the simultaneous emission of gravitational and electromagnetic waves, whereas the excitation of a polar mode is accompanied by the further emission of scalar radiation. In addition, gravitational and electromagnetic polar waves are emitted with frequencies and damping times different from the axial ones.
## Acknowledgments
F.P. whishes to thank Roberto De Pietri, Antonio Scotti and Michele Vallisneri for useful discussions and help in computing technicalities.
## VII APPENDIX
In this Appendix we provide all the elements of the matrix A that appears in equation (50) and in which are contained all the relevant features of the polar perturbations. It proves convenient to cast A in the following form
$$𝐀(r)=\frac{1}{D(r)}\left[G(r)𝐓(r)+P(r)𝐈\right],$$
(55)
where I is the identity matrix, and $`D(r)`$, $`G(r)`$ and $`𝐓(r)`$ are:
$`D(r)`$ $`=`$ $`4\mu Q_er^4\left(Q_{e}^{}{}_{}{}^{2}+Mr\right)^2`$
$`\times \left(3MQ_{e}^{}{}_{}{}^{4}7M^2Q_{e}^{}{}_{}{}^{2}rQ_{e}^{}{}_{}{}^{4}r+6M^3r^2M\mu ^2Q_{e}^{}{}_{}{}^{2}r^2+M^2\mu ^2r^3\right)`$
$`\times \left(4MQ_{e}^{}{}_{}{}^{4}14M^2Q_{e}^{}{}_{}{}^{2}rQ_{e}^{}{}_{}{}^{4}r+12M^3r^22M\mu ^2Q_{e}^{}{}_{}{}^{2}r^2+2M^2\mu ^2r^3\right)^2,`$
$`G(r)`$ $`=`$ $`\mu Q_e\left(2Mr\right)r^2\times \left(Q_{e}^{}{}_{}{}^{2}+Mr\right)`$
$`\left(3MQ_{e}^{}{}_{}{}^{4}7M^2Q_{e}^{}{}_{}{}^{2}rQ_{e}^{}{}_{}{}^{4}r+6M^3r^2M\mu ^2Q_{e}^{}{}_{}{}^{2}r^2+M^2\mu ^2r^3\right)/M,`$
$`P(r)`$ $`=`$ $`(1936M^6Q_{e}^{}{}_{}{}^{6}+324M^4Q_{e}^{}{}_{}{}^{8}160M^4\mu ^2Q_{e}^{}{}_{}{}^{8}+8M^2Q_{e}^{}{}_{}{}^{10}+Q_{e}^{}{}_{}{}^{12}`$
$`2432M^7Q_{e}^{}{}_{}{}^{4}r448M^5Q_{e}^{}{}_{}{}^{6}r+624M^5\mu ^2Q_{e}^{}{}_{}{}^{6}r+88M^3\mu ^2Q_{e}^{}{}_{}{}^{8}r+`$
$`2112M^8Q_{e}^{}{}_{}{}^{2}r^21136M^6Q_{e}^{}{}_{}{}^{4}r^21520M^6\mu ^2Q_{e}^{}{}_{}{}^{4}r^264M^4Q_{e}^{}{}_{}{}^{6}r^2`$
$`232M^4\mu ^2Q_{e}^{}{}_{}{}^{6}r^2+24M^4\mu ^4Q_{e}^{}{}_{}{}^{6}r^28M^2\mu ^2Q_{e}^{}{}_{}{}^{8}r^2`$
$`1152M^9r^3+2688M^7Q_{e}^{}{}_{}{}^{2}r^3+1984M^7\mu ^2Q_{e}^{}{}_{}{}^{2}r^3`$
$`272M^5\mu ^2Q_{e}^{}{}_{}{}^{4}r^3264M^5\mu ^4Q_{e}^{}{}_{}{}^{4}r^3+16M^3\mu ^2Q_{e}^{}{}_{}{}^{6}r^3`$
$`4M^3\mu ^4Q_{e}^{}{}_{}{}^{6}r^31152M^8r^4960M^8\mu ^2r^4+832M^6\mu ^2Q_{e}^{}{}_{}{}^{2}r^4+`$
$`464M^6\mu ^4Q_{e}^{}{}_{}{}^{2}r^416M^4\mu ^2Q_{e}^{}{}_{}{}^{4}r^428M^4\mu ^4Q_{e}^{}{}_{}{}^{4}r^4`$
$`16M^4\mu ^6Q_{e}^{}{}_{}{}^{4}r^4384M^7\mu ^2r^5224M^7\mu ^4r^5+64M^5\mu ^4Q_{e}^{}{}_{}{}^{2}r^5+`$
$`32M^5\mu ^6Q_{e}^{}{}_{}{}^{2}r^532M^6\mu ^4r^616M^6\mu ^6r^6)/8M^2\mu Q_e.`$
$`𝐓_{11}`$ $`=`$ $`2448M^6Q_{e}^{}{}_{}{}^{6}+324M^4Q_{e}^{}{}_{}{}^{8}160M^4\mu ^2Q_{e}^{}{}_{}{}^{8}+8M^2Q_{e}^{}{}_{}{}^{10}+Q_{e}^{}{}_{}{}^{12}+8M^3Q_{e}^{}{}_{}{}^{4}`$
$`(496M^464M^2Q_{e}^{}{}_{}{}^{2}+110M^2\mu ^2Q_{e}^{}{}_{}{}^{2}4Q_{e}^{}{}_{}{}^{4}+11\mu ^2Q_{e}^{}{}_{}{}^{4})r+8M^2Q_{e}^{}{}_{}{}^{2}`$
$`(408M^6+42M^4Q_{e}^{}{}_{}{}^{2}222M^4\mu ^2Q_{e}^{}{}_{}{}^{2}45M^2\mu ^2Q_{e}^{}{}_{}{}^{4}+3M^2\mu ^4Q_{e}^{}{}_{}{}^{4}`$
$`\mu ^2Q_{e}^{}{}_{}{}^{6})r^2+4M^3(288M^6+400M^4\mu ^2Q_{e}^{}{}_{}{}^{2}+108M^2\mu ^2Q_{e}^{}{}_{}{}^{4}`$
$`34M^2\mu ^4Q_{e}^{}{}_{}{}^{4}+4\mu ^2Q_{e}^{}{}_{}{}^{6}\mu ^4Q_{e}^{}{}_{}{}^{6})r^3+4M^4\mu ^2(144M^4`$
$`32M^2Q_{e}^{}{}_{}{}^{2}+52M^2\mu ^2Q_{e}^{}{}_{}{}^{2}4Q_{e}^{}{}_{}{}^{4}7\mu ^2Q_{e}^{}{}_{}{}^{4}4\mu ^4Q_{e}^{}{}_{}{}^{4})r^4+`$
$`32M^5\mu ^4(3M^2+2Q_{e}^{}{}_{}{}^{2}+\mu ^2Q_{e}^{}{}_{}{}^{2})r^516M^6\mu ^4(2+\mu ^2)r^6`$
$`𝐓_{21}`$ $`=`$ $`128M^5\mu Q_{e}^{}{}_{}{}^{7}+8M^2\mu Q_{e}^{}{}_{}{}^{5}(76M^4+8M^2\mu ^2Q_{e}^{}{}_{}{}^{2}`$
$`Q_{e}^{}{}_{}{}^{4})r+16M^3\mu Q_e^3(56M^4+22M^2Q_{e}^{}{}_{}{}^{2}8M^2\mu ^2Q_{e}^{}{}_{}{}^{2}+Q_{e}^{}{}_{}{}^{4}`$
$`2\mu ^2Q_{e}^{}{}_{}{}^{4})r^2+32M^4\mu Q_e(12M^424M^2Q_{e}^{}{}_{}{}^{2}+`$
$`2M^2\mu ^2Q_{e}^{}{}_{}{}^{2}+6\mu ^2Q_{e}^{}{}_{}{}^{4}+\mu ^4Q_{e}^{}{}_{}{}^{4})r^3+32M^5\mu Q_e`$
$`(12M^29\mu ^2Q_{e}^{}{}_{}{}^{2}2\mu ^4Q_{e}^{}{}_{}{}^{2})r^4+32M^6\mu ^3(4+\mu ^2)Q_er^5`$
$`𝐓_{31}`$ $`=`$ $`16M^3\mu \sqrt{2+\mu ^2}Q_{e}^{}{}_{}{}^{2}r(Q_{e}^{}{}_{}{}^{2}Mr)(14M^2Q_{e}^{}{}_{}{}^{2}Q_{e}^{}{}_{}{}^{4}`$
$`24M^3r+4M\mu ^2Q_{e}^{}{}_{}{}^{2}r+6M^2r^26M^2\mu ^2r^2\mu ^2Q_{e}^{}{}_{}{}^{2}r^2+2M\mu ^2r^3)`$
$`𝐓_{22}`$ $`=`$ $`2064M^6Q_{e}^{}{}_{}{}^{6}+260M^4Q_{e}^{}{}_{}{}^{8}160M^4\mu ^2Q_{e}^{}{}_{}{}^{8}+8M^2Q_{e}^{}{}_{}{}^{10}+Q_{e}^{}{}_{}{}^{12}+`$
$`4MQ_{e}^{}{}_{}{}^{4}(536M^6164M^4Q_{e}^{}{}_{}{}^{2}+116M^4\mu ^2Q_{e}^{}{}_{}{}^{2}+`$
$`6M^2Q_{e}^{}{}_{}{}^{4}+18M^2\mu ^2Q_{e}^{}{}_{}{}^{4}+Q_{e}^{}{}_{}{}^{6})r+8M^2Q_e^2`$
$`(72M^6+6M^4Q_{e}^{}{}_{}{}^{2}98M^4\mu ^2Q_{e}^{}{}_{}{}^{2}32M^2Q_{e}^{}{}_{}{}^{4}`$
$`43M^2\mu ^2Q_{e}^{}{}_{}{}^{4}5M^2\mu ^4Q_{e}^{}{}_{}{}^{4}Q_{e}^{}{}_{}{}^{6}+\mu ^2Q_{e}^{}{}_{}{}^{6})r^2+`$
$`4M^3Q_{e}^{}{}_{}{}^{2}(432M^4+256M^4\mu ^2+48M^2Q_{e}^{}{}_{}{}^{2}+`$
$`12M^2\mu ^2Q_{e}^{}{}_{}{}^{2}18M^2\mu ^4Q_{e}^{}{}_{}{}^{2}12\mu ^2Q_{e}^{}{}_{}{}^{4}`$
$`\mu ^4Q_{e}^{}{}_{}{}^{4})r^3+4M^4(288M^4144M^4\mu ^2+160M^2\mu ^2Q_{e}^{}{}_{}{}^{2}+`$
$`76M^2\mu ^4Q_{e}^{}{}_{}{}^{2}+8\mu ^2Q_{e}^{}{}_{}{}^{4}11\mu ^4Q_{e}^{}{}_{}{}^{4}`$
$`4\mu ^6Q_{e}^{}{}_{}{}^{4})r^4+16M^5\mu ^2(24M^212M^2\mu ^2+5\mu ^2Q_{e}^{}{}_{}{}^{2}+`$
$`2\mu ^4Q_{e}^{}{}_{}{}^{2})r^516M^6\mu ^4(2+\mu ^2)r^6`$
$`𝐓_{32}`$ $`=`$ $`8M^2\sqrt{2+\mu ^2}Q_e(16M^3Q_{e}^{}{}_{}{}^{6}36M^4Q_{e}^{}{}_{}{}^{4}r+`$
$`8M^2Q_{e}^{}{}_{}{}^{6}r8M^2\mu ^2Q_{e}^{}{}_{}{}^{6}r+Q_{e}^{}{}_{}{}^{8}r+192M^5Q_{e}^{}{}_{}{}^{2}r^2`$
$`52M^3Q_{e}^{}{}_{}{}^{4}r^216M^3\mu ^2Q_{e}^{}{}_{}{}^{4}r^22MQ_{e}^{}{}_{}{}^{6}r^2+`$
$`4M\mu ^2Q_{e}^{}{}_{}{}^{6}r^2144M^6r^3+48M^4Q_{e}^{}{}_{}{}^{2}r^3+72M^4\mu ^2Q_{e}^{}{}_{}{}^{2}r^3`$
$`16M^2\mu ^2Q_{e}^{}{}_{}{}^{4}r^34M^2\mu ^4Q_{e}^{}{}_{}{}^{4}r^348M^5\mu ^2r^4+`$
$`12M^3\mu ^2Q_{e}^{}{}_{}{}^{2}r^4+8M^3\mu ^4Q_{e}^{}{}_{}{}^{2}r^44M^4\mu ^4r^5)`$
$`𝐓_{33}`$ $`=`$ $`1936M^6Q_{e}^{}{}_{}{}^{6}+324M^4Q_{e}^{}{}_{}{}^{8}160M^4\mu ^2Q_{e}^{}{}_{}{}^{8}+8M^2Q_{e}^{}{}_{}{}^{10}+Q_{e}^{}{}_{}{}^{12}+`$
$`8M^3Q_{e}^{}{}_{}{}^{4}(304M^456M^2Q_{e}^{}{}_{}{}^{2}+78M^2\mu ^2Q_{e}^{}{}_{}{}^{2}+`$
$`11\mu ^2Q_{e}^{}{}_{}{}^{4})r+8M^2Q_e^2(264M^6142M^4Q_{e}^{}{}_{}{}^{2}`$
$`190M^4\mu ^2Q_{e}^{}{}_{}{}^{2}8M^2Q_{e}^{}{}_{}{}^{4}29M^2\mu ^2Q_{e}^{}{}_{}{}^{4}+`$
$`3M^2\mu ^4Q_{e}^{}{}_{}{}^{4}\mu ^2Q_{e}^{}{}_{}{}^{6})r^2+4M^3(288M^6+672M^4Q_{e}^{}{}_{}{}^{2}+`$
$`496M^4\mu ^2Q_{e}^{}{}_{}{}^{2}68M^2\mu ^2Q_{e}^{}{}_{}{}^{4}66M^2\mu ^4Q_{e}^{}{}_{}{}^{4}+`$
$`4\mu ^2Q_{e}^{}{}_{}{}^{6}\mu ^4Q_{e}^{}{}_{}{}^{6})r^3+4M^4(288M^4240M^4\mu ^2+`$
$`208M^2\mu ^2Q_{e}^{}{}_{}{}^{2}+116M^2\mu ^4Q_{e}^{}{}_{}{}^{2}4\mu ^2Q_{e}^{}{}_{}{}^{4}`$
$`7\mu ^4Q_{e}^{}{}_{}{}^{4}4\mu ^6Q_{e}^{}{}_{}{}^{4})r^4+32M^5\mu ^2(12M^2`$
$`7M^2\mu ^2+2\mu ^2Q_{e}^{}{}_{}{}^{2}+\mu ^4Q_{e}^{}{}_{}{}^{2})r^516M^6\mu ^4(2+\mu ^2)r^6`$
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# Meson decay in an independent quark model
## 1 Introduction
Though quantum chromodynamics is considered to be the underlying theory of strong interaction between quarks and gluons at the structural level of hadrons, many low-energy phenomena such as spectroscopy, static electromagnetic properties, can not be explained by first-principles application of QCD. Therefore one needs to resort to phenomenological models. For example, a potential model with an equally mixed scalar-vector harmonic potential<sup>1</sup> of independent quarks in a relativistic Dirac framework has been used to study several low-energy phenomena in the baryonic sector such as octet baryon masses<sup>2</sup>, magnetic moments<sup>3</sup>, weak electric form factors<sup>4</sup>, nucleon electromagnetic form factors and charge radii<sup>5</sup>. In addition to the harmonic potential, with power-law potential such as $`r^\nu `$($`\nu >0`$) heavy-heavy and light-heavy quarkonium states have been studied and the results are in good agreement with experimental results<sup>6</sup>. This model has also been successful in explaining pion mass, its decay constant<sup>7</sup> and the radiative decay<sup>8</sup> of ordinary light and heavy mesons. Because of this wide range of applicability of the model to both baryons and mesons, it has proved to be a rather simple and successful alternative to the cloudy bag model<sup>9</sup>. The purpose of this work is to extend its applicability to the study of the leptonic decay of vector mesons such as $`\rho `$, $`\omega `$, $`\varphi `$ and weak leptonic decay of light and heavy pseudoscalar mesons for an equally mixed scalar-vector power-law potential such as $`Ar^\nu `$. The leptonic decay width of heavier vector mesons in the charm and bottom quark sector have been studied in the nonrelativistic approach through the Van Royen-Weisskopf formula with radiative corrections<sup>10</sup>. But the same approach is not suitable for ordinary vector mesons in the light flavor sector, where the constituent quark dynamics is more relativistic. On the other hand, for weak leptonic decays, while many nonrelativistic quark model calculations<sup>11</sup> suggest that $`f_K>f_D>f_B`$ ($`f_M`$ is the weak leptonic decay constant), some of the models based on QCD sum rules<sup>12</sup> and lattice calculations<sup>13</sup> predict more or less a constant $`f_M`$ between $`K`$ and $`B`$ mesons. Capstick and Godfrey<sup>14</sup> calculate the hadronic matrix elements for the relativitized quark model expression for $`f_M`$. They find $`f_{B_c}>f_{D_s}>f_K>f_D>f_{B_s}>f_B>f_\pi `$, but the calculated value of the ratio $`f_K/f_\pi 1.75`$ is much higher than the experimental value of $`1.22`$. The leptonic decay widths and decay constants $`f_V`$ of the light vector mesons have been calculated by using an equally mixed scalar vector harmonic potential<sup>15</sup>. Also, weak leptonic decay constants, $`f_M`$, of pseudoscalar mesons have been calculated using the same potential and are found to satisfy $`f_{B_c}>f_{D_s}>f_D>f_K>f_{B_s}>f_\pi >f_B`$<sup>16</sup>. Potential used in Ref. 15 and 16 is harmonic. Since the potential between quarks is weaker than the harmonic potential, the potential used in this study is $`Ar^{0.2}+V_0`$ which gives very good results for heavy-heavy and light-heavy quarkonium states and radiative decay of ordinary light and heavy mesons<sup>6,8</sup>. Both in Ref. 15,16 and in this study, the potential does not include Coulomb term because Coulomb-like vector interaction are believed to have less prominent role for the light mesons.
## 2 Potential model
The quark-confining interaction in a hadron, which is believed to be generated by the nonperturbative multigluon mechanism, is not possible to calculate theoretically from first-principles within QCD. Therefore, from a phenomenological point of view, the present model assumes that the quark and antiquark in a hadron core are independently confined by an average flavor-independent potential $`V(r)`$ of the form
$$V(r)=\frac{1}{2}(1+\beta )(Ar^\nu +V_0),A>0\text{and}\nu >0.$$
(1)
For this potential, Dirac equation can be written as
$$\left[\stackrel{}{\alpha }\stackrel{}{p}+\beta m+V(r)\right]\mathrm{\Psi }(r)=E\mathrm{\Psi }(r)$$
(2)
where $`\stackrel{}{\alpha }`$ and $`\beta `$ are Dirac matrixes. Eq. (2) has two solutions with positive and negative energy given respectively in the forms
$$\psi _\mathrm{\Lambda }(r)=\left[\begin{array}{c}ig_\mathrm{\Lambda }(r)/r\\ \stackrel{}{\sigma }\widehat{r}f_\mathrm{\Lambda }(r)/r\end{array}\right]U_\mathrm{\Lambda }(\widehat{r})$$
(3)
$$\varphi _\mathrm{\Lambda }(r)=\left[\begin{array}{c}\stackrel{}{\sigma }\widehat{r}f_\mathrm{\Lambda }(r)/r\\ g_\mathrm{\Lambda }(r)/r\end{array}\right]\stackrel{~}{U}(\widehat{r})$$
(4)
where $`\mathrm{\Lambda }=(nljm)`$ represents the set of Dirac quantum numbers specifying the eigenmodes. The spin angular parts $`U_\mathrm{\Lambda }(\widehat{r})`$ and $`\stackrel{~}{U}_\mathrm{\Lambda }(\widehat{r})`$ are described as
$`U_{ljm}(\widehat{r})`$ $`=`$ $`{\displaystyle \underset{m_l,m_s}{}}lm_l{\displaystyle \frac{1}{2}}m_s|jmY_l^{m_l}(\widehat{r})\chi _{1/2}^{m_s},`$ (5)
$`\stackrel{~}{U}_{ljm}(\widehat{r})`$ $`=`$ $`(1)^{j+ml}U_{ljm}(\widehat{r})`$ (6)
Substituting Eq. (3) or Eq. (4) into Eq. (2), one obtains (for $`n=0`$, $`l=0`$)
$$\left[\frac{1}{2}\frac{d^2}{dr^2}+\lambda _qAr^\upsilon \right]g_\mathrm{\Lambda }(r)=\lambda _q(EmV_0)g_\mathrm{\Lambda }(r),$$
(7)
$$f_\mathrm{\Lambda }(r)=\frac{1}{\lambda _q}(\frac{d}{dr}\frac{1}{r})g_\mathrm{\Lambda }(r),$$
(8)
where $`\lambda _q=E+m`$, $`E`$ is energy of confined quark and $`m`$ is quark mass. Using the substitution $`\rho =(\lambda _qA)^{\frac{1}{\nu +2}}r`$ for convenience, Eq. (7) reduces to the form
$$\left[\frac{1}{2}\frac{d^2}{d\rho ^2}+\rho ^\nu \right]g(\rho )=ϵg(\rho ),$$
(9)
where
$$ϵ=\left[\frac{\lambda _q^\nu }{A^2}\right]^{\frac{1}{\nu +2}}(EmV_0),$$
(10)
and $`g(\rho )`$ is chosen as
$$g(\rho )\rho \mathrm{exp}\left(\left(x\rho \right)^d\right),$$
(11)
where $`x`$ and $`d`$ are variation parameters and they are obtained by minimizing $`ϵ`$. Hence they are solutions of
$$\frac{ϵ}{x}=0,$$
(12)
$$\frac{ϵ}{d}=0.$$
(13)
Using Eq. (9), Eq. (11) and Eq. (12), $`ϵ`$ is found to be
$$ϵ=\left(\frac{\nu +2}{\nu }\right)ax^2$$
(14)
where $`x=\left(\frac{b\nu }{2a}\right)^{\frac{1}{\nu +2}}`$, $`a=\left(\frac{d+1}{8}\right)2^{\frac{2}{d}}\frac{\mathrm{\Gamma }\left(\frac{1}{d}\right)}{\mathrm{\Gamma }\left(\frac{3}{d}\right)}`$, $`b=2^{\frac{\nu }{d}}\frac{\mathrm{\Gamma }\left(\frac{\nu +3}{d}\right)}{\mathrm{\Gamma }\left(\frac{3}{d}\right)}`$. When Eq. (13) is calculated numerically, $`d`$ is obtained as $`1.502`$ and using this value $`ϵ`$ is found as $`1.3268`$ which is very close to the results of the WKB method and $`1/N`$-expansion. Thus, the lowest eigenmodes corresponding to the positive and negative energies have the respective explicit forms
$$\psi _{\mathrm{\Lambda }(1s_{1/2})}(r)=\frac{1}{\sqrt{4\pi }}\left[\begin{array}{c}ig(r)/r\\ \stackrel{}{\sigma }\widehat{r}f(r)/r\end{array}\right]\chi _m,$$
(15)
$$\varphi _{\mathrm{\Lambda }(1s_{1/2})}(r)=\frac{1}{\sqrt{4\pi }}\left[\begin{array}{c}\stackrel{}{\sigma }\widehat{r}f(r)/r\\ ig(r)/r\end{array}\right]\stackrel{~}{\chi }_m,$$
(16)
where the two component spinors $`\chi _m`$ and $`\stackrel{~}{\chi }_m`$ denote $`\chi _{}=\left(\begin{array}{c}1\\ 0\end{array}\right)`$, $`\chi _{}=\left(\begin{array}{c}0\\ 1\end{array}\right)`$ and $`\stackrel{~}{\chi }_{}=\left(\begin{array}{c}0\\ i\end{array}\right)`$, $`\stackrel{~}{\chi }_{}=\left(\begin{array}{c}i\\ 0\end{array}\right)`$. Using Eq. (8) and Eq. (5), the radial parts in the upper and lower component solutions corresponding to a quark flavor $`q`$ are
$`g_q(r)`$ $`=`$ $`N_q\left({\displaystyle \frac{r}{r_{0q}}}\right)\mathrm{exp}\left(\left({\displaystyle \frac{x}{r_{0q}}}\right)^dr^d\right),`$ (17)
$`f_q(r)`$ $`=`$ $`{\displaystyle \frac{N_qd}{\lambda _qr_{0q}}}\left({\displaystyle \frac{x}{r_{0q}}}\right)^dr^d\mathrm{exp}\left(\left({\displaystyle \frac{x}{r_{0q}}}\right)^dr^d\right),`$ (18)
where, $`N_q`$ is normalization constant obtained from the equation
$`N_q^2{\displaystyle _0^{\mathrm{}}}\left(f_q^2(r)+g_q^2(r)\right)𝑑r=1.`$ (19)
## 3 Quark-antiquark momentum distribution
Knowing the quark-antiquark eigenmodes in the ground-state of meson, it is possible to obtain their corresponding momentum distribution amplitude. If $`G_q(p,\lambda ,\lambda ^{})`$ is the amplitude for finding a bound quark of flavor $`q`$ in its eigenmode $`\mathrm{\Phi }_{q\lambda }^{(+)}(r)`$ in a state of definite momentum $`p`$ and spin projection $`\lambda ^{}`$, then it is given by<sup>15,17</sup>
$`\mathrm{\Phi }_{q\lambda }^{(+)}(r)={\displaystyle \underset{\lambda ^{^{}}}{}}{\displaystyle d^3pG_q(p,\lambda ,\lambda ^{})\sqrt{\frac{m}{E_p}}U_q(p,\lambda ^{})\mathrm{exp}\left(i\stackrel{}{p}\stackrel{}{r}\right)}.`$ (20)
Where $`U_q(p,\lambda ^{})`$ is the usual free Dirac spinors. Eq. (20) can be easily inverted to yield
$$G_q(p,\lambda ,\lambda ^{})=\sqrt{\frac{m}{E_p}}U_q^{}(p,\lambda ^{^{}})d^3r\mathrm{\Phi }_{q\lambda }^{(+)}(r)\mathrm{exp}\left(i\stackrel{}{p}\stackrel{}{r}\right).$$
(21)
Thus, it is found as
$$G_q(p,\lambda ,\lambda ^{})=G_q(p)\delta _{\lambda \lambda ^{}},$$
(22)
where
$$G_q(p)=\frac{i\pi N_q}{\sqrt{2}dr_{0q}\beta ^{\frac{3}{d}}\lambda _q}(E_p+E_q)\sqrt{\frac{m+E_p}{E_p}}H(d,z),$$
(23)
where, $`E_p=\sqrt{p^2+m^2}`$, $`E_q`$ is the solution of Eq. (10), $`\beta =\left(\frac{x}{r_{0q}}\right)^d`$, $`z=\frac{pr_{0q}}{2x}`$ and
$`H(d,z)={\displaystyle \frac{d^2}{2^{\frac{2d+3}{2}}}}z^{\frac{3}{2}}{\displaystyle _0^{\mathrm{}}}y^{\frac{2d+1}{2}}\mathrm{exp}\left(\left({\displaystyle \frac{y}{2}}\right)^d\right)J_{3/2}(yz)𝑑y.`$ (24)
## 4 Leptonic decay widths
Now following Margolis and Mendel<sup>17</sup> one can represent the ground state of a neutral vector meson such as ($`\rho `$, $`\omega `$, $`\varphi `$) with a particular spin projection $`S_V`$ and zero momentum as
$$|V(0),S_V=\frac{\sqrt{3}}{\sqrt{N(0)}}\underset{q,\lambda _1,\lambda _2}{}d^3pG_q(p)C_{\lambda _1\lambda _2}^{S_V}\zeta _q^Vb_q^{}(p,\lambda _1)\stackrel{~}{b}_q^{}(p,\lambda _2)|0.$$
(25)
Here, $`b_q^{}(p,\lambda _1)`$ and $`\stackrel{~}{b}_q^{}(p,\lambda _2)`$ operating on the vacuum state are quark and antiquark creation operators, respectively. The summation with the flavor coefficient $`\zeta _q^V`$ and the spin configuration coefficient $`C_{\lambda _1\lambda _2}^{S_V}`$ represents the appropriate SU(6) spin-flavor structure of the particular vector meson $`V`$ with its spin projection $`S_V`$ and zero momentum. The factor $`\sqrt{3}`$ is due to the effective color singlet configuration of the meson. $`N(0)`$ represents the overall normalization, which is given by
$`N(0)={\displaystyle \frac{1}{(2\pi )^3}}{\displaystyle _0^{\mathrm{}}}d^3p\left|G_q(p)\right|^2.`$ (26)
Assuming that the main contribution to the leptonic decay process of neutral vector mesons such as $`\rho `$, $`\omega `$, $`\varphi `$ comes from single virtual photon creation from the annihilation of the bound quark-antiquark pair inside of the meson, $`S`$-matrix element in configuration space can be written as<sup>15</sup>
$`S_{fi}`$ $`=`$ $`e^{}(k_1,\delta _1)e^+(k_2,\delta _2)|ie^2{\displaystyle }d^4x_1d^4x_2(\overline{\psi }_e(x_2)\gamma ^\mu \psi _e(x_2)D_{\mu \nu }(x_2x_1)`$ (27)
$`\times `$ $`{\displaystyle \underset{q}{}}e_q\overline{\psi }_q(x_1)\gamma ^\nu \psi _q(x_1))|V,S_V`$
Where, $`D_{\mu \nu }(x_2x_1)`$ is the photon propagator, $`\psi _e(x)`$ and $`\psi _q(x)`$ are the free lepton and quark fields, respectively. After some standard calculations (details of the calculation can be found in Ref. 15), one obtains
$$\mathrm{\Gamma }\left(Ve^+e^{}\right)=\frac{4\pi }{3}\alpha _{em}^2M_Vf_V^2.$$
(28)
Where, $`f_V`$ is known as the leptonic decay constant and in this model it can be written as
$$f_V^2=\frac{2e_q_V^2I_V^2}{3\pi ^2M_V^3J_V}.$$
(29)
Here,
$`e_q_{\rho ,\omega ,\varphi }`$ $`=`$ $`({\displaystyle \frac{1}{\sqrt{2}}},{\displaystyle \frac{1}{3\sqrt{2}}},{\displaystyle \frac{1}{3}}),`$
$`I_V`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑pp^2\left(2+{\displaystyle \frac{m}{E_p}}\right)G_q(p),`$
$`J_V`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑pp^2\left|G_q(p)\right|^2.`$ (30)
$`I_V`$ and $`J_V`$values can be calculated numerically by computer.
## 5 Weak leptonic decay widths
Here, the weak leptonic decay of charged pseudoscalar mesons such as $`\pi ^\pm `$, $`K^\pm `$, $`D^\pm `$, $`D_s^\pm `$, $`B^\pm `$ and $`B_c^\pm `$ are considered. Assuming that the main contribution to the weak leptonic decay processes comes from the single virtual boson creation from the annihilation of the bound quark-antiquark pair inside the pseudoscalar meson $`M`$, the $`S`$-matrix element in configuration space is written as<sup>16</sup>
$`S_{fi}`$ $`=`$ $`l(k_1,\delta _1)\overline{\nu }_l(k_2,\delta _2)|\left({\displaystyle \frac{iG_F}{\sqrt{2}}}\right){\displaystyle d^4x\overline{\psi }_l(x)\gamma ^\mu \left(1\gamma ^5\right)\psi _l(x)}`$ (31)
$`\times `$ $`{\displaystyle \underset{qm,qn}{}}\nu _{q_mq_n}\psi _{q_m}(x)\gamma _\mu \left(1\gamma ^5\right)\psi _{q_n}(x)|M(0).`$
Where, $`G_F`$ is the Fermi coupling constant, $`\nu _{q_mq_n}`$ are the CKM matrix elements and $`|M(0)`$ is given by<sup>16</sup>
$`|M(0)`$ $`=`$ $`\sqrt{{\displaystyle \frac{3}{N(0)}}}{\displaystyle C_{q_1q_2}^M(\lambda _{1,}\lambda _2)}`$ (32)
$`\times `$ $`{\displaystyle d^3p\left[G_{q_1}(p)G_{q_2}^{}(p)\right]^{\frac{1}{2}}b_{q_1}^{}(p,\lambda _1)\stackrel{~}{b}_{q_2}^{}(p,\lambda _2)|0},`$
where, $`C_{q_1q_2}^M(\lambda _{1,}\lambda _2)`$ stands for the appropriate SU(6) spin-flavor coefficients for the pseudoscalar meson $`M`$. $`N(0)`$ represents the overall normalization, which is given by
$$N(0)=\frac{1}{(2\pi )^3}_0^{\mathrm{}}d^3p\left(G_{q_1}(p)G_{q_2}^{}(p)\right).$$
(33)
After some standard calculations, which can be found in Ref. 16, one obtains
$$\mathrm{\Gamma }\left(Ml\overline{\nu }_l\right)=\frac{G_F^2}{8\pi }\left|\nu _{q_1q_2}\right|^2M_pm_l^2\left(1\frac{m_l^2}{M_p^2}\right)^2f_M^2,$$
(34)
where $`f_M`$ is the weak decay constant, having the form
$$f_M^2=\frac{3I_M^2}{2\pi ^2M_pJ_M}.$$
(35)
Here $`M_p`$ is the mass of pseudoscalar meson, $`m_l`$ is the mass of lepton and $`I_M`$ and $`J_M`$ are found as
$`I_M`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑pp^2A(p)\left[G_{q_1}(p)G_{q_2}^{}(p)\right]^{\frac{1}{2}},`$
$`J_M`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑pp^2\left[G_{q_1}(p)G_{q_2}^{}(p)\right],`$ (36)
and
$$A(p)=\frac{\left(E_{p_1}+m_{q_1}\right)\left(E_{p_2}+m_{q_2}\right)p^2}{\left[E_{p_1}E_{p_2}\left(E_{p_1}+m_{q_1}\right)\left(E_{p_2}+m_{q_2}\right)\right]^{\frac{1}{2}}},$$
(37)
where $`E_{p_i}=\sqrt{p^2+m_{q_i}^2}`$.
## 6 Results
The calculations involve the potential parameters of the model $`(A,V_0,\nu )`$ and the quark masses $`\left(m_u=m_d,m_s,m_c,m_b\right)`$ . The potential parameters are chosen to be $`A=0.68\text{GeV},V_0=0.3961\text{GeV},\nu =0.2`$.
The light-quark masses $`m_u=m_d`$ and $`m_s`$ are obtained from $`\omega `$ and $`\varphi `$ mesons as $`m_u=m_d=0.078\text{GeV}`$, $`m_s=0.3\text{GeV}`$, and from $`D^\pm `$ and $`B^\pm `$ $`m_c=1.3\text{GeV}`$, $`m_b=4.81\text{GeV}`$.
When these parameters are used the masses of $`D,D_s,B,B_s,B_c,w`$ and $`\varphi `$ can be calculated to be almost the same with their experimental values. Also, with the same parameters, radiative decay widths of light and heavy mesons have been calculated and the obtained results are close to their experimental values<sup>8</sup>. Using this potential parameters and quark masses given above, the calculated results are shown in Table I, Table II and Table III.
## 7 Conclusion
In this paper, the independent particle model approach has been used to investigate leptonic decay of light vector mesons and weak leptonic decay of light and heavy pseudoscalar mesons. The Dirac equation in a power-law potential has been solved with variation technique using special type wave function. Leptonic decay widths and decay constant $`f_V`$ and weak leptonic decay widths and decay constant $`f_M`$ have been calculated and are compared with the results of other theoretical investigations and experiments. The results say to us that the potential used in this study can be considered as interaction potential for quarks. Because of the model structure it is assumed that quarks do not interact with each other. However, there are spin-spin and other hyperfine interactions between quarks. Adding these interactions, good results, especially in the calculation of meson mass spectra, could be obtained. Such interaction terms can be found in Ref. 24,25.
## Acknowledgments
We thank T.M. Aliev and H. Akcay for useful discussions.
## References
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P.Leal Ferrera, Lett. Nuovo cimento 20 (1977) 157,
P.Leal Ferrera and N.Zagury, ibid 20 (1977) 511
2. N.Barik and B.K.Dash, Phys. Rev. D33 (1986) 1925
3. N.Barik and B.K.Dash, Phys. Rev. D34 (1986) 2803
4. N.Barik, B.K.Dash and M.Das, Phys. Rev. D32 (1985) 1725
5. N.Barik and B.K.Dash, Phys. Rev. D34 (1986) 2052
6. H.Akcay and H.Ciftci, J. Phys. G22 (1996) 455
7. N.Barik, B.K.Dash and P.C.Dash, Pramana J. Phys. 29 (1987) 543
8. N.Barik, P.C.Dash and A.R.Panda, Phys. Rev. D46 (1992) 3856,
N.Barik, P.C.Dash, Phys.Rev. D49 (1994) 299, H.C̣iftci (unpublished)
9. A.W.Thomas, Adv. Nucl. Phys. 13 (1983) 1
10. R.Van Royen and V.F.Weisskopf, Nuovo cimento A50 (1967) 617
11. S.N.Sinha, Phys. Lett. B178 (1986) 110,
G.Godfrey, Phys. Rev. D33 (1986) 1391
12. L.J.Reinders, Phys. Rev. D38 (1988) 947,
S.Narison, Phys. Lett. B198 (1987) 104,
C.A.Dominguez and N.Power, Phys. Lett. B197 (1987) 423
13. M.B.Govale, et al., Phys. Lett. B206 (1988) 113,
R.M.Wokshyn et al. Phys. Rev. D39 (1989) 978
14. S.Capstick and S.Godfrey, Phys. Rev. D41 (1990) 2856
15. N.Barik, P.C.Dash and A.R.Panda, Phys. Rev. D47 (1993) 1001
16. N.Barik and P.C.Dash, Phys. Rev. D47 (1993) 2788
17. B.Margolis and R.R.Mandal, Phys. Rev. D28 (1983) 468,
C.Hayne and N.Isgnur, Phys.Rev. D25 (1982) 1944
18. H.Kraseman, Phys.Lett. B96 (1980) 397
19. E.Glowich, Phys.Lett. B91 (1980) 271
20. M.Claudsun, Harward Report No. 91 (1981) (unpublished)
21. C.Bernard, et al., Phys. Rev. D38 (1980) 3540
22. H.W.Hamber, Phys. Rev. D39 (1989) 896
23. C.Caso et al., The European Physical Journal C3 (1998) 1
24. Ho-Meoyng Choi and Chueng-Ryong Ji, Phys. Lett. B460 (1999) 461
25. Ho-Meoyng Choi and Chueng-Ryong Ji, Phys. Rev. D59 (1999) 074015
Table I. Leptonic decay widhts and the decay constant $`f_V`$ in keV in comparison with the results of other researchers and the experiment.
| Meson | $`\mathrm{\Gamma }\left(Ve^+e^{}\right)`$ | Experiment<sup>23</sup> | Ref. 15 | Ref. 17 | |
| --- | --- | --- | --- | --- | --- |
| $`\rho `$ | 6.37 | 6.77$`\pm `$0.32 | 6.26(8.1) | 7.8 | |
| $`\omega `$ | 0.684 | 0.6$`\pm `$0.02 | 0.67(0.87) | 0.84 | |
| $`\varphi `$ | 1.46 | 1.37$`\pm `$0.05 | 1.58(1.84) | 1.69 | |
| Meson | $`f_V`$ | Experiment<sup>23</sup> | Ref. 15 | Ref. 17 | |
| $`\rho `$ | 0.193 | 0.2$`\pm `$0.04 | 0.19(0.22) | 0.21 | |
| $`\omega `$ | 0.0624 | 0.06$`\pm `$0.01 | 0.06(0.07) | 0.07 | |
| $`\varphi `$ | 0.0813 | 0.08$`\pm `$0.01 | 0.08(0.09) | 0.07 | |
Table II. Decay constants of pseudoscalar mesons in MeV in comparison with the results of other model and the experiment. Experimental values are taken from Ref. 16.
| Model | $`f_\pi `$ | $`f_K`$ | $`f_D`$ | $`f_{D_s}`$ | $`f_B`$ | $`f_{B_s}`$ | $`f_{B_c}`$ | |
| --- | --- | --- | --- | --- | --- | --- | --- | --- |
| Expt.<sup>16,24</sup> | 131.73$`\pm `$0.15 | 160.6$`\pm `$1.3 | $`<`$219 | 137-304 | - | - | - | |
| This work | 131.4 | 150.2 | 181.1 | 235.9 | 208.3 | 257.3 | 377.2 | |
| Ref. 16 | 138 | 157 | 161 | 205 | 122 | 154 | 221 | |
| Ref. 14 | 100 | 153 | 149 | 160 | 96 | 111 | 141 | |
| Ref. 20 | - | - | 172 | 196 | 149 | 170 | 255 | |
| Ref. 18 | 139 | 176 | 150 | 210 | 125 | 175 | 425 | |
| Ref. 19 | 178 | 182 | 148 | 166 | 98 | - | - | |
| Ref. 21 | - | - | 174$`\pm `$53 | 234$`\pm `$72 | 105$`\pm `$34 | 155$`\pm `$75 | - | |
| Ref. 22 | 141$`\pm `$21 | 155$`\pm `$21 | 282$`\pm `$28 | - | 183$`\pm `$28 | - | - | |
Table III. Partial decay widths $`\mathrm{\Gamma }\left(Ml\overline{\nu _l}\right)`$ in MeV and the branching ratio $`B\left(Ml\overline{\nu _l}\right)`$ of pseudoscalar mesons in comparison with the experiment.(̇$`B\left(Ml\overline{\nu _l}\right)=\tau _M\mathrm{\Gamma }\left(Ml\overline{\nu _l}\right)`$). $`\tau _M`$ is the mean lifetime of meson $`M`$
| Prosses | $`\mathrm{\Gamma }\left(Ml\overline{\nu _l}\right)`$ | $`B\left(Ml\overline{\nu _l}\right)`$ | Expt.<sup>23</sup> $`B\left(Ml\overline{\nu _l}\right)`$ |
| --- | --- | --- | --- |
| $`\pi ^\pm \mu ^\pm \overline{\nu }_\mu `$ | 2.523 10<sup>-14</sup> | 0.994 | 0.999877$`\pm `$0.0000004 |
| $`\pi ^\pm e^\pm \overline{\nu }_e`$ | 3.238 10<sup>-18</sup> | 1.276 10<sup>-4</sup> | (1.23$`\pm `$0.004) 10<sup>-4</sup> |
| $`K^\pm \mu ^\pm \overline{\nu }_\mu `$ | 3.000 10<sup>-14</sup> | 0.563 | 0.6351$`\pm `$0.0018 |
| $`K^\pm e^\pm \overline{\nu }_e`$ | 7.724 10<sup>-19</sup> | 1.45 10<sup>-5</sup> | (1.55$`\pm `$0.07) 10<sup>-5</sup> |
| $`D^\pm \mu ^\pm \overline{\nu }_\mu `$ | 1.795 10<sup>-13</sup> | 2.87 10<sup>-4</sup> | $`<`$7.2 10<sup>-4</sup> |
| $`D_s^\pm \mu ^\pm \overline{\nu }_\mu `$ | 6.240 10<sup>-12</sup> | 4.41 10<sup>-3</sup> | 4$`{}_{2.0}{}^{+2.2}\times `$10<sup>-3</sup> |
| $`B^\pm \mu ^\pm \overline{\nu }_\mu `$ | 1.693 10<sup>-16</sup> | 4.15 10<sup>-7</sup> | $`<`$2.1 10<sup>-5</sup> |
| $`B_c^\pm \mu ^\pm \overline{\nu }_\mu `$ | 8.960 10<sup>-14</sup> | - | - |
| $`D^\pm \tau ^\pm \overline{\nu }_\tau `$ | 4.720 10<sup>-13</sup> | 7.540 10<sup>-4</sup> | - |
| $`D_s^\pm \tau ^\pm \overline{\nu }_\tau `$ | 6.090 10<sup>-11</sup> | 4.3 10<sup>-2</sup> | (7$`\pm `$4) 10<sup>-2</sup> |
| $`B^\pm \tau ^\pm \overline{\nu }_\tau `$ | 3.770 10<sup>-14</sup> | 9.25 10<sup>-5</sup> | $`<`$0.57 10<sup>-3</sup> |
| $`B_c^\pm \tau ^\pm \overline{\nu }_\tau `$ | 2.140 10<sup>-11</sup> | - | - |
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# Commuting self-adjoint extensions of symmetric operators defined from the partial derivatives
## 1. Introduction
Recently several papers have appeared on commuting non-self-adjoint operators and their spectral theory; see, e.g., \[LKMV95\]. The present paper concerns the case when the given commuting operators are unbounded and symmetric, but non-self-adjoint. A concrete class of operators is studied, and we address the questions of when commuting extension operators exist and, when they do exist, what their structural properties are.
The problem of understanding commuting symmetric, but non-self-adjoint, unbounded operators also has an origin in mathematical physics \[AvIv95, Bel89, HaKo91, Pav79\]. The terminology from physics is “hermitian”, or *“formally* self-adjoint”, for symmetry, i.e., for the identity $`Sf\text{ }h=f\text{ }Sh`$ for all vectors $`f,h`$ in the domain of the operator $`S`$. The simplest case of this is the problem of assigning quantum mechanical boundary conditions for free particles confined in a box. More specifically, the problem here corresponds to the quantum-mechanical trajectories of a particle confined to a region of tube type, e.g., a unit cube. It is “free” except for the boundary conditions, and variations of the boundary conditions (as considered here) correspond to different physics. For single operators, von Neumann solved (or made precise) the problem by use of the Cayley transform, and considering instead the extension problem for partial isometries. But this approach does not work well in the case of several operators. Powers (in \[Pow71, Pow74\]) introduced an algebraic approach for understanding several operators, but the present problem is very concrete and does not lend itself easily to the algebraic techniques introduced by Powers.
Closely connected to the problem of finding commuting self-adjoint extensions of $`\frac{1}{i}\frac{}{x_j}`$, $`j=1,\mathrm{},d`$, on $`C_c^{\mathrm{}}\left(\mathrm{\Omega }\right)`$ is the corresponding spectral question: If commuting self-adjoint extensions do exist, then it is known that the common eigenfunctions of the extension operators must be of the form $`e_\lambda :=e^{i\lambda x}`$ for special values of $`\lambda ^d`$. Hence the spectral problem is that of finding when a given pair $`(\mathrm{\Omega },\mathrm{\Lambda })`$ satisfies the condition that $`\left\{e_\lambda |_\mathrm{\Omega }:\lambda \mathrm{\Lambda }\right\}`$ is an orthogonal basis in the Hilbert space $`^2\left(\mathrm{\Omega }\right)`$. We note that this so-called *spectral pair* condition is very restrictive, and so it explains the rigid geometric configurations $`(\mathrm{\Omega },\mathrm{\Lambda })`$ which admit solutions. But it also serves to motivate recent very interesting developments on overcomplete systems; see, e.g., \[Kem99a, Kem99b\].
The setting of *spectral pairs* in $`d`$ real dimensions involves two subsets $`\mathrm{\Omega }`$ and $`\mathrm{\Lambda }`$ in $`^d`$ such that $`\mathrm{\Omega }`$ has finite and positive $`d`$-dimensional Lebesgue measure, and $`\mathrm{\Lambda }`$ is an index set for an orthogonal $`^2\left(\mathrm{\Omega }\right)`$-basis $`e_\lambda `$ of exponentials, i.e.,
(1.1)
$$e_\lambda \left(x\right)=e^{i2\pi \lambda x},x\mathrm{\Omega },\lambda \mathrm{\Lambda }$$
where $`\lambda x=_{j=1}^d\lambda _jx_j`$. We use vector notation $`x=(x_1,\mathrm{},x_d)`$, $`\lambda =(\lambda _1,\mathrm{},\lambda _d)`$, $`x_j,\lambda _j`$, $`j=1,\mathrm{},d`$. The basis property refers to the Hilbert space $`^2\left(\mathrm{\Omega }\right)`$ with inner product
(1.2)
$$f\text{ }g_\mathrm{\Omega }:=_\mathrm{\Omega }\overline{f\left(x\right)}g\left(x\right)𝑑x$$
where $`dx=dx_1\mathrm{}dx_d`$, and $`f,g^2\left(\mathrm{\Omega }\right)`$. The corresponding norm is
(1.3)
$$f_\mathrm{\Omega }^2:=f\text{ }f_\mathrm{\Omega }=_\mathrm{\Omega }\left|f\left(x\right)\right|^2𝑑x,$$
as usual. It follows that the spectral pair property for a pair $`(\mathrm{\Omega },\mathrm{\Lambda })`$ is equivalent to
$$\mathrm{\Lambda }\mathrm{\Lambda }=\{\lambda \lambda ^{}:\lambda ,\lambda ^{}\mathrm{\Lambda }\}$$
being contained in the *zero-set* of the complex function
(1.4)
$$z_\mathrm{\Omega }e^{i2\pi zx}dx=:F_\mathrm{\Omega }\left(z\right)$$
where $`z=(z_1,\mathrm{},z_d)^d`$, and $`zx:=_{j=1}^dz_jx_j`$, and the corresponding $`e_\lambda `$-set $`\{e_\lambda :\lambda \mathrm{\Lambda }\}`$ being *total* in $`^2\left(\mathrm{\Omega }\right)`$. Recall, totality means that the span of the $`e_\lambda `$’s is dense in $`^2\left(\mathrm{\Omega }\right)`$ relative to the $`_\mathrm{\Omega }`$-norm, or, equivalently, that $`f=0`$ is the only $`^2\left(\mathrm{\Omega }\right)`$-solution to:
$$f\text{ }e_\lambda _\mathrm{\Omega }=0\text{,\hspace{1em}for all }\lambda \mathrm{\Lambda }.$$
## 2. Spectral pairs
The theory of *spectral pairs* was developed in previous joint papers by the coauthors \[JoPe92, JoPe94, JoPe96\]. A set $`\mathrm{\Omega }`$ with finite nonzero Lebesgue measure is called a *spectral set* if $`(\mathrm{\Omega },\mathrm{\Lambda })`$ is a spectral pair for some set $`\mathrm{\Lambda }`$. We recall that Fuglede showed \[Fug74\] that the disk and the triangle in two dimensions are *not spectral sets.* By the disk and the triangle we mean the usual versions, respectively, $`\{(x_1,x_2)^2:x_1^2+x_2^2<1\}`$ and $`\{(x_1,x_2)^2:0<x_1,\mathrm{\hspace{0.33em}0}<x_2,x_1+x_2<1\}.`$ Note that, for the present discussion, it is inessential whether or not the sets $`\mathrm{\Omega }`$ are taken to be open, but it is essential for the following theorem which we will need. It is due to Fuglede and the coauthors; see \[Fug74, Jor82, Ped87, JoPe92\].
If $`\mathrm{\Omega }^d`$ is open, then we consider the partial derivatives $`\frac{}{x_j}`$, $`j=1,\mathrm{},d`$, defined on $`C_c^{\mathrm{}}\left(\mathrm{\Omega }\right)`$ as unbounded skew-symmetric operators in $`^2\left(\mathrm{\Omega }\right)`$. The corresponding versions $`\frac{1}{2\pi \sqrt{1}}\frac{}{x_j}`$ are symmetric of course. We say that $`\mathrm{\Omega }`$ has the *extension property* if there are commuting self-adjoint extension operators $`H_j`$, i.e.,
(2.1)
$$\frac{1}{2\pi i}\frac{}{x_j}H_j,j=1,\mathrm{},d.$$
We say that the containment $`AB`$ holds for two operators $`A`$ and $`B`$ if the graph of $`A`$ is contained in that of $`B`$. (For details, see \[ReSi\] and \[DS2\].) *Commutativity* for the extension operators $`H_j`$ is in the strong sense of spectral resolutions. Since the $`H_j`$’s are assumed self-adjoint, each one has a projection-valued *spectral resolution* $`E_j`$, i.e., an $`^2\left(\mathrm{\Omega }\right)`$-projection-valued Borel measure on $``$, such that $`E_j\left(\right)=I_{^2\left(\mathrm{\Omega }\right)}`$, and
(2.2)
$$H_j=_{\mathrm{}}^{\mathrm{}}\lambda E_j\left(d\lambda \right),$$
for $`j=1,\mathrm{},d`$. The strong commutativity is taken to mean
(2.3)
$$E_j\left(\mathrm{\Delta }\right)E_j^{}\left(\mathrm{\Delta }^{}\right)=E_j^{}\left(\mathrm{\Delta }^{}\right)E_j\left(\mathrm{\Delta }\right)$$
for all $`j,j^{}=1,\mathrm{},d`$, and all Borel subsets $`\mathrm{\Delta },\mathrm{\Delta }^{}`$. Extensions commuting in a weaker sense were considered in \[Fri87\].
Our analysis is based on von Neumann’s deficiency-space characterization of the self-adjoint extensions of a given symmetric operator \[vNeu29\]. Let $`\mathrm{\Omega }`$ be an open set with finite Lebesgue measure. For each $`j`$, the deficiency spaces corresponding to $`\frac{1}{i}\frac{}{x_j}`$ are infinite-dimensional. It follows that each $`\frac{1}{i}\frac{}{x_j}`$ has “many” self-adjoint extensions. The main problem (not addressed by von Neumann’s theory) is the selection of a commuting set $`H_1,H_2,\mathrm{},H_d`$ of extensions. In fact, for some $`\mathrm{\Omega }`$ (e.g., when $`d=2`$, the disk and the triangle) it is impossible to select a commuting set $`H_1,H_2,\mathrm{},H_d`$ of extensions.
We have (see \[Fug74, Jor82, Ped87, JoPe92\])
###### Theorem 2.1.
(Fuglede, Jorgensen, Pedersen) Let $`\mathrm{\Omega }^d`$ be open and connected with finite and positive Lebesgue measure. Then $`\mathrm{\Omega }`$ has the extension property if and only if it is a spectral set. Moreover, with $`\mathrm{\Omega }`$ given, there is a one-to-one correspondence between the two sets of subsets:
(2.4)
$$\{\mathrm{\Lambda }^d:(\mathrm{\Omega },\mathrm{\Lambda })\text{ is a spectral pair}\}$$
and
(2.5)
$$\begin{array}{c}\{\mathrm{\Lambda }^d:\mathrm{\Lambda }\text{ is the joint spectrum of some commutative}\hfill \\ \hfill \text{family }(H_1,\mathrm{},H_d)\text{ of self-adjoint etensions}\}.\end{array}$$
This correspondence is determined as follows:
1. If the extensions $`(H_1,\mathrm{},H_d)`$ are given, then $`\lambda \mathrm{\Lambda }`$ if and only if
(2.6)
$$e_\lambda \underset{j}{}domain\left(H_j\right).$$
2. If, conversely, $`(\mathrm{\Omega },\mathrm{\Lambda })`$ is a spectral pair at the outset, then the ansatz (2.6) and
(2.7)
$$H_je_\lambda =\lambda _je_\lambda ,\lambda \mathrm{\Lambda }$$
determine uniquely a set of commuting extensions.
If $`\mathrm{\Omega }`$ is only assumed open, then the spectral-set property implies the extension property, but not conversely.
###### Corollary 2.2.
Suppose $`\mathrm{\Omega }`$ is open and connected. It follows then that a discrete set $`\mathrm{\Lambda }`$ is the joint spectrum of some commuting self-adjoint extension operators $`H_j`$, $`j=1,\mathrm{},d`$, if and only if $`(\mathrm{\Omega },\mathrm{\Lambda })`$ is a spectral pair.
###### Remark 2.3.
The simplest case of a disconnected $`\mathrm{\Omega }`$ which has the extension property, but which is not a spectral set, is in $`d=1`$, and we may take $`\mathrm{\Omega }=0,12,4`$, i.e., the union of two intervals with a doubling and separation. The example was noted first in \[Fug74\] and is based on the simple observation that the polynomial $`1+z^2+z^3`$ has no roots $`z`$ on the circle $`|z|=1`$.
Some of the interest in spectral pairs derives from their connection to *tilings.* A subset $`\mathrm{\Omega }^d`$ with nonzero measure is said to *tile* $`^d`$ if there is a set $`L^d`$ such that the translates $`\{\mathrm{\Omega }+l:lL\}`$ cover $`^d`$ up to measure zero, and if the intersections
(2.8)
$$\left(\mathrm{\Omega }+l\right)\left(\mathrm{\Omega }+l^{}\right)\text{ for }ll^{}\text{ in }L$$
have measure zero. We will call $`(\mathrm{\Omega },L)`$ a *tiling pair* and we will say that $`L`$ is a *set of translations*. The *Spectral-Set conjecture* (see \[Fug74, Jor82, Ped87, LaWa96c, LaWa97a\]) states:
###### Conjecture 2.4.
Let $`\mathrm{\Omega }^d`$ have positive and finite Lebesgue measure. Then $`\mathrm{\Omega }`$ is a spectral set if and only if there is a set $`L`$ of translations which make $`\mathrm{\Omega }`$ tile $`^d`$.
###### Lemma 2.5.
If $`\mathrm{\Omega }=I^d`$, then the zero-set for the function $`F_\mathrm{\Omega }`$ in (1.4) is
(2.9)
$$𝐙_{I^d}=\{z^d\mathrm{}\left\{0\right\}:j\{1,\mathrm{},d\}\mathrm{s}.\mathrm{t}.z_j\mathrm{}\left\{0\right\}\}.$$
###### Proof.
The function $`F_{I^d}()`$ factors as follows.
(2.10)
$$F_{I^d}\left(z\right)=\underset{j=1}{\overset{d}{}}e^{i\pi z_j}\frac{\mathrm{sin}\pi z_j}{\pi z_j}$$
for $`z=(z_1,\mathrm{},z_d)^d`$, with the interpretation that the function $`z\frac{\mathrm{sin}\pi z}{\pi z}`$ is $`1`$ when $`z=0`$ in $``$.∎
###### Remark 2.6.
What is special about $`𝐙_\mathrm{\Omega }`$ for $`\mathrm{\Omega }=I^d`$, as opposed to the general form of $`\mathrm{\Omega }`$, is that $`𝐙_{I^d}\left\{0\right\}`$ is the *Cayley graph* of the group $`\mathrm{\Gamma }=^d`$ with generators
$$S=\{(\pm 1,0,\mathrm{},0),\mathrm{},(0,\mathrm{},\pm 1,0,\mathrm{},0),\mathrm{},(0,\mathrm{},0,\pm 1)\}.$$
We recall from \[BKS, Chapter 10\] the definition of the Cayley graph $`G(\mathrm{\Gamma },S)`$ of a discrete group $`\mathrm{\Gamma }`$ with generators $`S`$, $`eS`$. When $`\mathrm{\Gamma },S`$ are given, $`G(\mathrm{\Gamma },S)`$ is the graph with vertex set $`\mathrm{\Gamma }`$ in which two vertices $`\gamma _1,\gamma _2`$ are the two ends of an edge iff $`\gamma _1^1\gamma _2S`$. This gives a non-oriented graph, without any loop or multiple edge.
## 3. Two dimensions
We begin with the following simple observation in one dimension for $`\mathrm{\Omega }=I=[0,1`$. (For details, see \[JoPe92, ReSi\].)
###### Proposition 3.1.
The only subsets $`\mathrm{\Lambda }`$ such that $`(I,\mathrm{\Lambda })`$ is a spectral pair are the translates
(3.1)
$$\mathrm{\Lambda }_\alpha :=\alpha +=\{\alpha +n:n\}$$
where $`\alpha `$ is some fixed real number.
In two dimensions, the corresponding result is more subtle, but the possibilities may still be enumerated as follows:
###### Theorem 3.2.
(\[JoPe99\]) The only subsets $`\mathrm{\Lambda }^2`$ such that $`(I^2,\mathrm{\Lambda })`$ is a spectral pair must belong to either one or the other of the two classes, indexed by a number $`\alpha `$, and a sequence $`\{\beta _m[0,1:m\}`$, where
(3.2) $`\mathrm{\Lambda }`$ $`=\{\left(\begin{array}{c}\alpha +m\\ \beta _m+n\end{array}\right):m,n\}`$
or
(3.3) $`\mathrm{\Lambda }`$ $`=\{\left(\begin{array}{c}\beta _n+m\\ \alpha +n\end{array}\right):m,n\}.`$
Each of the two types occurs as the spectrum of a pair for the cube $`I^2`$, and each of the sets $`\mathrm{\Lambda }`$ as specified is a set of translation vectors which produces a tiling of $`^2`$ by the cube $`I^2`$.
###### Proof.
See \[JoPe99\] for details. The following are some remarks of relevance to the general extension problem for operators.
The assertion in the theorem about $`\mathrm{\Lambda }`$-translations tiling the plane with $`I^2`$ is also clear from (3.2)–(3.3), and it is illustrated graphically in Figures 1 and 2.
It is clear that the pattern (3.2) for $`d=2`$ continues to higher dimensions as follows:
(3.4)
$$\left(\begin{array}{c}\alpha +k_1\\ \beta \left(k_1\right)+k_2\\ \gamma (k_1,k_2)+k_3\\ \mathrm{}\\ \zeta (k_1,k_2,\mathrm{},k_{d1})+k_d\end{array}\right)$$
with $`k_1,k_2,\mathrm{},k_d`$, and
$`\beta `$ $`:[0,1,`$
$`\gamma `$ $`:\times [0,1,`$
$`\mathrm{}`$
$`\zeta `$ $`:^{d1}[0,1.`$
Of course, then there are the obvious modifications of those cases resulting from permutation of the $`d`$ coordinates; but the assertion is that, when $`d10`$, these configurations do *not* suffice for cataloguing all the possible spectra $`\mathrm{\Lambda }`$ which turn $`(I^d,\mathrm{\Lambda })`$ into an $`^d`$-spectral pair.
We now turn to the non-trivial spectral-theoretic content of the conclusion of the theorem. We claim that the two cases (3.2)–(3.3) suffice when $`d=2`$. Note that the sequence $`\beta :[0,1`$ is completely arbitrary.
We will show in Theorem 5.1 below that, up to a single translation in the plane, the possibilities for the coordinates of points in a spectrum $`\mathrm{\Lambda }`$ for $`I^2`$ are given by two sequences $`\xi _m`$, $`\eta _n`$ satisfying the following two cocycle relations:
(3.5) $`\left(e^{i\xi _{m+k}}e^{i\xi _m}\right)\left(1e^{i\eta _n}\right)`$ $`=0`$
and
(3.6) $`\left(e^{i\eta _{n+l}}e^{i\eta _n}\right)\left(1e^{i\xi _m}\right)`$ $`=0`$
as identities in $`m,n`$, and $`k,l\mathrm{}\left\{0\right\}`$. Note that the respective sequences are determined from this only up to $`2\pi `$ at each coordinate place.
Simple algebra shows that the two identities (3.5)–(3.6) imply the following single identity
(3.7)
$$\left(1e^{i\xi _{m+k}}\right)\left(1e^{i\eta _n}\right)=\left(1e^{i\xi _m}\right)\left(1e^{i\eta _{n+l}}\right)$$
again for all $`m,n`$ and $`k,l\mathrm{}\left\{0\right\}`$. But it follows from (3.7) that at least one of the two sequences, $`1e^{i\xi _m}`$ or $`1e^{i\eta _n}`$, must then vanish identically. This yields the connection to the two cases for $`\mathrm{\Lambda }`$ stated in (3.2)–(3.3) of the theorem.
Hence the result giving two classes for $`\mathrm{\Lambda }`$ in Theorem 3.2 may be derived from our more general result in Section 5.
The proof sketch of Theorem 3.2 is completed for now, but details will be resumed in Section 5 below.∎
## 4. Operator extensions
We saw in Theorem 2.1 that in some cases the existence problem for spectral pairs, i.e., the question of when some given open subset $`\mathrm{\Omega }`$ in $`^d`$ has an orthogonal basis $`\{e_\lambda :\lambda \mathrm{\Lambda }\}`$ in $`^2\left(\mathrm{\Omega }\right)`$ for some set $`\mathrm{\Lambda }`$ in $`^d`$, may be reformulated as a problem about existence of commuting self-adjoint extensions of the operators $`\{\frac{1}{i}\frac{}{x_j}:j=1,\mathrm{},d\}`$ with common (dense) domain $`C_c^{\mathrm{}}\left(\mathrm{\Omega }\right)`$ in $`^2\left(\mathrm{\Omega }\right)`$. Suppose for the moment that $`\mathrm{\Omega }=0,1\times \mathrm{\Omega }_2`$ where $`\mathrm{\Omega }_2`$ is some subset in $`^{d1}`$ of finite positive $`\left(d1\right)`$-dimensional Lebesgue measure. We then have the following classification of the self-adjoint extensions $`H`$ of $`\frac{1}{i}\frac{}{x_1}`$.
###### Theorem 4.1.
The symmetric operator $`\frac{1}{i}\frac{}{x_1}`$ in $`^2\left(0,1\times \mathrm{\Omega }_2\right)`$ with dense domain $`𝒟`$ consisting of $`\phi ^2\left(0,1\times \mathrm{\Omega }_2\right)`$ such that $`\phi (,y)C_c^{\mathrm{}}\left(0,1\right)`$ for all $`y\mathrm{\Omega }_2`$, has self-adjoint extensions indexed by unitary operators $`V`$ in $`^2\left(\mathrm{\Omega }_2\right)`$ in such a way that the (unique) extension $`H_V`$ is determined by its core domain being of the form
(4.1)
$$𝒟_V=\{\phi (x_1,y)+e^{x_1}h\left(y\right)+e^{1x_1}\left(Vh\right)\left(y\right):\phi 𝒟,h^2\left(\mathrm{\Omega }_2\right)\}$$
and
(4.2)
$$iH_V\left(\phi +e^{x_1}h+e^{1x_1}Vh\right)=\frac{\phi }{x_1}+e^{x_1}he^{1x_1}Vh,$$
for $`\phi 𝒟`$ and $`h^2\left(\mathrm{\Omega }_2\right)`$. We shall interpret the implicit boundary condition dictating some extension $`H_V`$ as
(4.3)
$$f(1,)=U_V\left(f(0,)\right),$$
$`f𝒟_V`$ where the partial isometry $`U_V`$ is given by
(4.4)
$$W_V=\left(eI+V\right)\left(I+eV\right)^1,U_V=\mathrm{exp}W_V.$$
Conversely, $`V`$ may be calculated from $`U_V`$ by
(4.5)
$$V=\left(IeW_V\right)^1\left(W_VeI\right),$$
and in each case, the fractional linear transform, and its inverse, are well defined.
###### Proof.
The proof is based on von Neumann’s deficiency-space analysis of self-adjoint extensions, and we refer to \[vNeu29\], \[ReSi\], and \[Jor79\] for background material on the theory of operator extensions. If $`S`$ is a symmetric operator with dense domain $`𝒟`$ in a Hilbert space $``$, then it has self-adjoint extensions if and only if the two spaces
(4.6)
$$\left((iI\pm S)𝒟\right)^{}=:𝒟_\pm $$
have the same dimensions. In that case, the corresponding extensions are given by *partial isometries* between the respective defect spaces $`𝒟_+`$ and $`𝒟_{}`$ (see \[vNeu29\], \[ReSi\], or \[DS2\]). For convenience, we have chosen a slightly different “normalization” in our treatment of the Cayley transform (4.4) and its inverse (4.5). We did not normalize the functions $`e^{x_1}`$ and $`e^{1x_1}`$ in the defect spaces. They have $`^2\left(I\right)`$-norm equal to $`\left(\frac{e^21}{2}\right)^{\frac{1}{2}}`$. The fact that $`U_V`$ in (4.4) then defines a partial isometry as claimed amounts to the identities:
If $`\psi (x_1,y)=e^{x_1}h\left(y\right)+e^{1x_1}\left(Vh\right)\left(y\right)`$ as in (4.1), then
$$\psi (1,y)=eh\left(y\right)+\left(Vh\right)\left(y\right)=\left(eI+V\right)h,$$
and
$$\psi (0,y)=h\left(y\right)+eVh\left(y\right)=\left(I+eV\right)h.$$
This means that the vectors in the domain (4.1) are given by the boundary conditions (4.3) which in turn determine the unitary one-parameter group
$$𝒰_V\left(t\right):=\mathrm{exp}\left(itH_V\right),t.$$
This group is defined from (4.3) by using translation modulo $``$ in the $`x_1`$-variable. Then the operator $`U_V`$ in (4.4) is used in defining the representation $`t𝒰_V\left(t\right)`$ via induction from $``$.
If $`V:𝒟_+𝒟_{}`$ is a partial isometry, then the domain of the corresponding extension $`H`$ ($`H=H_V`$) is
$$\{\phi +h_++Vh_+:\phi 𝒟,h_+𝒟_+\}$$
and
(4.7)
$$iH_V\left(\phi +h_++Vh_+\right)=iS\phi +h_+Vh_+.$$
It follows that the lemma amounts to an identification of the *defect spaces* $`𝒟_\pm `$ when the symmetric operator is as specified. When the variables in $`\mathrm{\Omega }=0,1\times \mathrm{\Omega }_2`$ are separated as $`(x_1,y)`$, $`0<x_1<1`$, $`y=(x_2,\mathrm{},x_d)\mathrm{\Omega }_2`$, then vectors $`h_\pm 𝒟_\pm `$ are precisely the solutions to
(4.8)
$$S^{}h_\pm =\pm ih_\pm .$$
This amounts to solving
$$\frac{}{x_1}h_\pm (x_1,y)=\pm h_\pm (x_1,y)$$
in the sense of distributions, but with the restrictions $`h_\pm ^2\left(0,1\times \mathrm{\Omega }_2\right)`$. The result of the lemma then follows from von Neumann’s characterization. If the minimal operator is not closed at the outset, then the resulting self-adjoint extension comes from passing to the operator closure in the formulas (4.2) and (4.7).∎
###### Corollary 4.2.
Let $`V`$ be a unitary operator in $`^2\left(\mathrm{\Omega }_2\right)`$ and let $`H_V`$ be the self-adjoint extension operator described in Theorem 4.1 in (4.2)–(4.3). Then $`H_V`$ generates a unitary one-parameter group $`\{U_V\left(t\right):t\}`$ in $`^2\left(0,1\times \mathrm{\Omega }_2\right)`$ which may be realized (up to unitary equivalence) in the Hilbert space $`_V`$ of measurable functions $`f:^2\left(\mathrm{\Omega }_2\right)`$, satisfying
(4.9)
$$f\left(x_1+1\right)=U_V\left(f\left(x_1\right)\right),$$
for all $`x_1`$, where $`U_V`$ is the operator from (4.4) in Theorem 4.1, and the norm on $`_V`$ is defined by
$$f__V^2=_0^1f\left(x_1\right)_{^2\left(\mathrm{\Omega }_2\right)}^2𝑑x_1.$$
In this space the group $`U_V\left(t\right):_V_V`$ is given by
(4.10)
$$\left(U_V\left(t\right)f\right)\left(x_1\right)=f\left(x_1t\right),\text{ for }x_1,t.$$
The unitary isomorphism of $`_V`$ onto $`^2\left(0,1\times \mathrm{\Omega }_2\right)=^2(0,1,^2\left(\mathrm{\Omega }_2\right))`$ is simply the restriction to $`0,1`$ in the $`x_1`$-variable. Finally, if $`U_V\left(t\right)`$ is computed in $`^2\left(0,1\times \mathrm{\Omega }_2\right)`$, the formula is
(4.11)
$$\left(U_V\left(t\right)f\right)f(x_1,)=\{\begin{array}{cc}f(x_1t,)\hfill & \text{ if }0t<x_1<1,\hfill \\ U_V\left(f(x_1t,)\right)\hfill & \text{ if }0<x_1t1.\hfill \end{array}$$
###### Proof.
The realization on the space $`_V`$ is the interpretation of $`U_V`$ as a unitary representation of the group $``$ which is induced from the subgroup $``$ via formula (4.10). The advantage of this viewpoint is that the spectral resolution of the unitary operator $`U_V`$ leads directly to an associated direct integral decomposition for the unitary one-parameter group $`\{U_V\left(t\right):t\}`$ which is generated by the extension operator $`H_V`$.∎
When the corollary is applied to $`^2\left(I\times I\right)`$ from Section 3 we note that the respective unitary one-parameter groups, $`U_x\left(s\right)`$ and $`U_y\left(t\right)`$, on $`^2\left(I^2\right)`$ which are generated by self-adjoint extension operators of $`\frac{1}{i}\frac{}{x}`$ and $`\frac{1}{i}\frac{}{y}`$ with domain $`C_c^{\mathrm{}}\left(I^2\right)`$, are induced representations in the sense of (4.9)–(4.10). For the extensions of $`\frac{1}{i}\frac{}{x}`$, the boundary-unitary from (4.9) is acting on $`^2\left(\left\{0<y<1\right\}\right)`$. But we shall view it as a unitary operator in $`^2\left(I\times I\right)=^2\left(I_x\right)^2\left(I_y\right)`$ via $`UIU_2`$ with $`U_2`$ acting in the $`y`$-variable. A similar observation applies to the unitary one-parameter group $`\{U_y\left(t\right):t\}`$ acting on $`^2\left(I^2\right)`$ and generated by one of the self-adjoint extensions of $`\frac{1}{i}\frac{}{y}`$. Hence the boundary conditions for $`\{U_x\left(s\right):s\}`$ are given by a unitary $`UIU_2`$ with $`U_2`$ acting in the second variable, while those of $`\{U_y\left(t\right):t\}`$ are determined by a second unitary operator $`V`$ in $`^2\left(I^2\right)`$, now of the form $`VV_1I`$ with $`V_1`$ acting in the first variable of $`^2\left(I\times I\right)`$.
With this terminology we have the following preliminary result for the square $`I^2`$ in the plane.
###### Theorem 4.3.
Let $`U_x\left(s\right)`$ be the unitary one-parameter group on $`^2\left(I\times I\right)`$, and let $`U_2`$ be the corresponding unitary boundary operator acting in the second variable $`y`$. Then $`U_2`$ commutes with the phase-periodic translation in the $`y`$-variable for a phase angle $`\beta `$ if and only if there is a real-valued sequence $`\{\phi _n:n\}`$ such that
(4.12)
$$U_x\left(s\right)e_{m+\phi _n}e_{n+\beta }=e^{i2\pi \left(m+\phi _n\right)s}e_{m+\phi _n}e_{n+\beta }$$
for all $`s`$ and $`m,n`$, where for $`(\xi ,\eta )^2`$, $`e_\xi e_\eta (x,y)=e_\xi \left(x\right)e_\eta \left(y\right)=e^{i2\pi \left(\xi x+\eta y\right)}`$, restricted to $`(x,y)I^2`$.
###### Proof.
Recall that some fixed unitary one-parameter group $`\{U_x\left(s\right):s\}`$ on $`^2\left(I\times I\right)`$ is determined uniquely by the corresponding boundary operator $`IU_2`$. But it follows from Proposition 3.1 that $`U_2`$ satisfies the commutativity property of the theorem if and only if it is diagonalized by the basis functions $`\{e_{n+\beta }:n\}`$ in $`^2\left(I_y\right)`$ for some $`\beta [0,1`$, i.e., if, for some sequence $`\phi _n`$,
(4.13)
$$\left(U_2e_{n+\beta }\right)\left(y\right)=e^{i2\pi \phi _n}e_{n+\beta }\left(y\right).$$
But, according to Corollary 4.2, this means that $`U_x\left(s\right)`$ as an induced representation decomposes accordingly, which is to say that the basis vectors $`e_{m+\phi _n}e_{n+\beta }`$ simultaneously diagonalize each operator $`U_x\left(s\right)`$ as stated in formula (4.12).∎
###### Remark 4.4.
For more details on the operator-theoretic approach to spectrum and to tiles, we refer to \[Jor87b, Jor89b, Ped87, Ped96\]
## 5. Cocycles in two dimensions
In this section, we continue with the self-adjoint extensions of the two commuting minimal operators $`\frac{1}{i}\frac{}{x}`$ and $`\frac{1}{i}\frac{}{y}`$ with common dense domain $`C_c^{\mathrm{}}\left(I^2\right)`$ in $`^2\left(I^2\right)`$.
###### Theorem 5.1.
Consider two commuting unitary one-parameter groups $`U_x\left(s\right)`$ and $`U_y\left(t\right)`$ on $`^2\left(I\times I\right)`$ with respective boundary unitaries $`U_2`$ and $`V_1`$. Then:
1. Either $`U_2`$ is of the form $`aI_{^2\left(I_y\right)}`$ for a scalar $`a`$, or else $`V_1`$ commutes with periodic translation in the $`x`$-variable.
2. Either $`V_1`$ is of the form $`bI_{^2\left(I_x\right)}`$ for some scalar $`b`$, or else $`U_2`$ commutes with periodic translation in the $`y`$-variable.
3. In case $`U_2=e^{i2\pi \alpha }I_{^2\left(I_y\right)}`$, then
(5.1)
$$U_x\left(s\right)\left(e_{\alpha +m}g\right)=e^{i2\pi \left(\alpha +m\right)s}e_{\alpha +m}g$$
for all $`m`$ and $`g^2\left(I_y\right)`$.
4. In case $`V_1=e^{i2\pi \beta }I_{^2\left(I_x\right)}`$, then
(5.2)
$$U_y\left(t\right)\left(fe_{\beta +n}\right)=e^{i2\pi \left(\beta +n\right)t}fe_{\beta +n}$$
for all $`f^2\left(I_x\right)`$ and $`n`$.
###### Remark 5.2.
It follows that the conclusion in Theorem 4.3 is satisfied when the two one-parameter groups commute, i.e., when
(5.3)
$$U_x\left(s\right)U_y\left(t\right)=U_y\left(t\right)U_x\left(s\right)$$
is assumed, $`s,t`$. Specifically, it will then always be the case that $`U_2`$ commutes with some phase-periodic translation in the $`y`$-variable, while $`V_1`$ commutes with some (possibly different) phase-periodic translation in the $`x`$-variable. (Also note that (5.3) is a reformulation of (2.3) in the case $`d=2`$. Furthermore (5.3) signifies the presence of a unitary representation of $`^2.`$)
###### Proof of Theorem 5.1.
When the two one-parameter groups $`U_x\left(s\right)`$ and $`U_y\left(t\right)`$ are written in the form (4.11) from Corollary 4.2, then the alternatives in (4.11) may be expanded as follows. Let $`\tau _s`$ denote periodic translation in $`^2\left(0,1\right)`$, and let $`P_s`$ denote the projection of $`^2\left(0,1\right)`$ onto $`^2\left(0,s\right)`$, with $`P_s^{}=IP_s`$ denoting then the projection onto the complement $`^2\left(s,1\right)`$, for $`s[0,1]`$. We have $`P_0=0`$ and $`P_1=I_{^2\left(0,1\right)}`$. Then from (4.11) we get
(5.4) $`U_x\left(s\right)`$ $`=\tau _sP_s^{}I+\tau _sP_sU_2`$
and
(5.5) $`U_y\left(t\right)`$ $`=I\tau _tP_t^{}+V_1\tau _tP_t.`$
The assumed commutativity (5.3) then takes the form:
$$\begin{array}{c}\tau _sP_s^{}V_1\tau _tP_t+\tau _sP_sU_2\tau _tP_t^{}+\tau _sP_sV_1U_2\tau _tP_t\hfill \\ \hfill =\tau _sP_s\tau _tP_t^{}U_2+V_1\tau _sP_s^{}\tau _tP_t+V_1\tau _sP_s\tau _tP_tU_2.\end{array}$$
If $`V_1`$ is not a scalar times $`I_{^2\left(I_x\right)}`$ then two terms on either side are independent when evaluated on $`fg`$. Hence both $`U_2\tau _tP_t^{}=\tau _tP_t^{}U_2`$ and $`U_2\tau _tP_t=\tau _tP_tU_2`$ hold. Addition of these two identities yields $`U_2\tau _t=\tau _tU_2`$ which is the commutativity of $`U_2`$ with periodic translation.
If on the other hand $`V_1`$ is a scalar, then it follows from the argument in Section 4 that (iv) must hold.
The two possibilities for the other boundary operator $`U_2`$ lead to cases (i) and (iii) by symmetry.∎
###### Corollary 5.3.
Consider unitary one-parameter groups $`U_x\left(s\right)`$ and $`U_y\left(t\right)`$ as in Theorem 5.1 and suppose the corresponding boundary operators $`U_2`$ and $`V_1`$ diagonalize as follows (identities in $`n,m`$):
(5.6) $`U_2e_{n+\beta }`$ $`=e^{i2\pi \alpha _n}e_{n+\beta }`$
and
(5.7) $`V_1e_{m+\alpha }`$ $`=e^{i2\pi \beta _m}e_{m+\alpha }`$
for some $`\alpha ,\beta `$. The sequences $`\alpha _n,\beta _m`$ will be chosen taking values in $`[0,1`$. Then the commutativity (5.3) for the two groups holds if and only if the two sequences satisfy a certain cocycle property: Let $`a_n:=e^{i2\pi \alpha _n}`$ and $`b_m:=e^{i2\pi \beta _m}`$. Then the two identities
(5.8) $`\left(b_mb_{m+k}\right)\left(1a_n\right)`$ $`=0,m,n,k\mathrm{}\left\{0\right\}`$
and
(5.9) $`\left(a_na_{n+l}\right)\left(1b_m\right)`$ $`=0,m,n,l\mathrm{}\left\{0\right\}`$
are equivalent to the commutativity (5.3). If commutativity holds, we must have $`\left(1a_n\right)\left(1b_m\right)0`$, $`n,m`$. Hence we get a spectral pair with spectrum $`\mathrm{\Lambda }`$ having one of the two forms
(i) $`\{\left(\begin{array}{c}\alpha +m\\ n+\beta _m\end{array}\right):m,n\}\text{ if }\alpha _n`$ $`0,`$
or
(ii) $`\{\left(\begin{array}{c}m+\alpha _n\\ \beta +n\end{array}\right):m,n\}\text{ if }\beta _m`$ $`0.`$
The derivation of the two cocycle identities (5.8)–(5.9) from commutativity (5.3) at the end of the proof is based on the following corollary of independent interest:
###### Corollary 5.4.
Let $`U=IU_2`$ and $`V=V_1I`$ be the respective boundary operators of the one-parameter unitary groups $`U_x\left(s\right)`$ and $`U_y\left(t\right)`$ acting on $`^2\left(I\times I\right)`$. Then, if (5.6)–(5.7) hold for some $`\alpha ,\beta `$ and some sequences as specified, it follows that the respective one-parameter groups may be expanded in the common basis $`E(m,n)=E_{(\alpha ,\beta )}(m,n):=e_{m+\alpha }^{\left(1\right)}e_{n+\beta }^{\left(2\right)}`$, $`(m,n)^2`$, as follows: There are complex sequences $`\left\{s_k\right\}_k`$ and $`\left\{t_l\right\}_l`$ so that, if we define
(5.10)
$$s_0^{}:=1s_0,t_0^{}:=1t_0$$
and
(5.11)
$$s_k^{}:=s_k\text{ }\text{(}\text{for }k0\text{), }t_l^{}:=t_l\text{ }\text{(}\text{for }l0\text{),}$$
then
(5.12) $`U_x\left(s\right)E(m,n)`$ $`={\displaystyle \underset{k}{}}e^{i2\pi \left(m+\alpha +k\right)s}\left(s_k^{}+s_ka_n\right)E(m+k,n)`$
and
(5.13) $`U_y\left(t\right)E(m,n)`$ $`={\displaystyle \underset{l}{}}e^{i2\pi \left(n+\beta +l\right)t}\left(t_l^{}+t_lb_m\right)E(m,n+l).`$
The two one-parameter groups $`U_x\left(s\right)`$ and $`U_y\left(t\right)`$ commute if and only if the cocycle identities (5.8)–(5.9) hold.
###### Proof.
Recall from (5.4)–(5.5) that the two one-parameter groups are expressed in terms of multiplication operators on $`^2\left(0,1\right)`$ with the respective indicator functions $`\chi _{0,s}`$ and $`\chi _{0,t}`$. The sequences (5.10)–(5.11) are the Fourier coefficients of these indicator functions, acting by multiplication in $`^2\left(I\right)`$, and the relations (5.10)–(5.11) simply reflect the following two obvious identities,
$`\chi _{0,s}+\chi _{[s,1}`$ $`=1`$
and
$`\chi _{0,t}+\chi _{[t,1}`$ $`=1,`$
as functions on the unit interval. When the resulting formulas (5.12)–(5.13) are substituted into
(5.14)
$$U_x\left(s\right)U_y\left(t\right)E(m,n)=U_y\left(t\right)U_x\left(s\right)E(m,n)$$
the equivalence to (5.8)–(5.9) results.∎
## 6. Quasicrystals
For the spectral pairs $`(I^d,\mathrm{\Lambda })`$ in dimensions $`d=2,3`$, we noted that each of the candidates for spectrum $`\mathrm{\Lambda }`$ tiles $`^d`$ with $`\mathrm{\Lambda }`$-translates of $`I^d`$. (See Theorems 3.2 and 4.3.) But reviewing formulas (3.2)–(3.3) and (3.4), and (7.4) in the next section, for the possible sets $`\mathrm{\Lambda }`$ which serve as $`I^d`$-spectrum, we find functions $`\alpha ,\beta ,\mathrm{}`$ on $``$ or $`^k`$ which describe the particular set $`\mathrm{\Lambda }`$. Since all the candidates for $`\mathrm{\Lambda }`$ make tilings, there is a direct *geometric* interpretation for these functions; but we note in the present section that there is also a *spectral-theoretic* significance which derives from diffraction considerations of quasicrystals; see \[Sen95\], \[Hof95\], and \[BoTa87\].
In this setting, diffractions show up as discrete components of the spectral distribution
$$D_\mathrm{\Lambda }\left(x\right)=\underset{\lambda \mathrm{\Lambda }}{}e_\lambda \left(x\right)=\underset{\lambda \mathrm{\Lambda }}{}e^{i2\pi \lambda x}.$$
We say that a spectrum $`\mathrm{\Lambda }`$ ($`^d`$) has a *diffraction pattern* if there is a pair $`(M,c)`$ where $`M`$ is a subset of $`^d`$ and $`c`$ is a function (measuring intensity) defined on $`M`$ such that
$$D_\mathrm{\Lambda }\left(x\right)=\underset{mM}{}c\left(\mu \right)\delta \left(x\mu \right),$$
i.e., the spectral distribution is a weighted sum of point-masses, supported on some (discrete) subset $`M`$ in $`^d`$. Note that the interpretation in both of the summations involving $`D_\mathrm{\Lambda }()`$ is to be understood as Schwartz distributions; that is if the respective sums are evaluated on a testing function $`\phi C_c^{\mathrm{}}\left(^d\right)`$, then the first sum yields $`_{\lambda \mathrm{\Lambda }}\stackrel{~}{\phi }\left(\lambda \right)`$ where $`\stackrel{~}{\phi }\left(\lambda \right)=_^de_\lambda \left(x\right)\phi \left(x\right)𝑑x`$, while the second sum is $`_{\mu M}c\left(\mu \right)\phi \left(\mu \right)`$. We also note that, by the Poisson summation formula, the condition is satisfied if $`\mathrm{\Lambda }=^d=M`$, and the density (intensity) function $`c`$ is $`c1`$ on $`M`$.
We shall also need the following definition: A function $`\xi `$ on $``$ is said to be *quasi-periodic* if there are positive numbers $`\omega _1,\mathrm{},\omega _r`$, which are independent over $``$, and functions $`\xi _1,\mathrm{},\xi _r`$ such that $`\xi _j`$ has $`\omega _j`$ as period, and $`\xi =_{j=1}^r\xi _j`$. The condition on $`\xi _j`$ amounts to the generalized Fourier expansion
$$\xi _j\left(x\right)=\underset{n}{}c_j\left(n\right)e^{i2\pi \frac{nx}{\omega _j}}.$$
In the following result we show that, if the functions which define a spectrum $`\mathrm{\Lambda }`$ for some $`I^d`$ are quasi-periodic, then it follows that $`\mathrm{\Lambda }`$ has a diffraction pattern. We will not state the result in the widest generality as it will be clear that the idea in the simplest case carries over to the variations in higher dimensions. Even for $`d=2`$, Theorem 3.2 shows that there are two classes of $`\mathrm{\Lambda }`$ corresponding to (3.2) and (3.3) respectively. In the following we will treat only (3.2), but the result applies to (3.3) *mutatis mutandis.*
###### Theorem 6.1.
Let
$$\mathrm{\Lambda }=\{\left(\begin{array}{c}m\\ \beta \left(m\right)+n\end{array}\right):m,n\}$$
for some function $`\beta :`$ and suppose $`\beta `$ extends to a function on $``$ which is quasi-periodic with periods $`\omega _1,\mathrm{},\omega _r`$, independent over $``$. Then it follows that $`(I^2,\mathrm{\Lambda })`$ is a spectral pair with diffraction pattern; specifically, there is a density function $`c:^r\times `$ such that
$$D_\mathrm{\Lambda }(x,y)=\underset{k^r}{}\underset{n}{}c(k,n)\underset{m}{}\delta \left(x\underset{i=1}{\overset{r}{}}\frac{k_i}{\omega _i}m\right)\delta \left(yn\right)$$
with the density $`c(k_1,\mathrm{},k_r,n)`$ derived from the Bohr almost periodic Fourier expansion applied to $`\beta `$.
###### Proof.
Consider the formula $`D_\mathrm{\Lambda }(x,y)=_m_ne^{i2\pi \left(mx+\left(\beta \left(m\right)+n\right)y\right)}`$ and expand the inside function, $`me^{i2\pi \beta \left(m\right)y}`$ according to the quasi-periodicity assumption on $`\beta `$: specifically,
$`e^{i2\pi \beta \left(m\right)y}`$ $`={\displaystyle \underset{j=1}{\overset{r}{}}}e^{i2\pi \xi _j\left(m\right)y}`$
$`={\displaystyle \underset{j=1}{\overset{r}{}}}{\displaystyle \underset{k_j}{}}c^{\left(j\right)}\left(k_j\right)e^{i2\pi \frac{mk_j}{\omega _j}}`$
$`={\displaystyle \underset{k_1}{}}\mathrm{}{\displaystyle \underset{k_r}{}}c^{\left(1\right)}\left(k_1\right)\mathrm{}c^{\left(r\right)}\left(k_r\right)e^{i2\pi m_{j=1}^r\frac{k_j}{\omega _j}}.`$
Setting $`c\left(k\right):=_{j=1}^rc^{\left(j\right)}\left(k_j\right)`$ and using
$$\underset{m}{}e^{i2\pi m\left(x+_{j=1}^r\frac{k_j}{\omega _j}\right)}=\underset{m}{}\delta \left(x\underset{j=1}{\overset{r}{}}\frac{k_j}{\omega _j}m\right)$$
together with Poisson summation (also in the second variable) we arrive at the desired formula.∎
## 7. Higher dimensions
The following definitions help summarize the results for $`d=2`$: We say that the one-parameter unitary groups on $`^2\left(I\times I\right)`$ generated by self-adjoint extensions of the respective partial derivatives $`\frac{1}{i}\frac{}{x}`$ and $`\frac{1}{i}\frac{}{y}`$ on $`C_c^{\mathrm{}}\left(I\times I\right)`$ are *quasi-commuting* if the conditions (5.6)–(5.7) hold. Recall this means that the respective boundary operators commute with some phase-periodic translation in the opposite variable. We then showed in Theorem 5.1 that the commutativity property (5.3), for the unitary groups $`U_x\left(s\right)`$ and $`U_y\left(t\right)`$, implies *quasi-commutativity.* Finally we showed in Corollary 5.3 that, among the quasi-commuting extensions, those that in fact commute (in the sense of (5.3)) are characterized by the two cocycle identities (5.8)–(5.9).
It is clear that *quasi-commutativity* can be defined analogously for $`d>2`$. It follows from Theorem 2.1 that commutativity of $`d`$ self-adjoint extensions of the respective partial derivatives $`\{\frac{1}{i}\frac{}{x_j}:j=1,\mathrm{},d\}`$, on $`C_c^{\mathrm{}}\left(I^d\right)^2\left(I^d\right)`$, is equivalent to the *spectral-pair* condition for $`(I^d,\mathrm{\Lambda })`$. Moreover, if commuting self-adjoint extensions exist (i.e., $`\frac{1}{i}\frac{}{x_j}H_j`$, $`H_j^{}=H_j`$, $`j=1,\mathrm{},d`$), then we may take $`\mathrm{\Lambda }`$ to be the joint spectrum of the family $`\left\{H_j\right\}_{j=1}^d`$. Conversely, commuting operators $`H_j`$ may easily be associated with some spectrum $`\mathrm{\Lambda }`$ in a spectral pair $`(I^d,\mathrm{\Lambda })`$. Hence, for $`d=2`$, our results in Section 5 provide a complete classification of the commuting (and also the quasi-commuting) self-adjoint extensions of $`\left\{\frac{1}{i}\frac{}{x_j}\right\}_{j=1}^d`$.
In higher dimensions, we still have boundary operators corresponding to each self-adjoint extension of the partials $`\frac{1}{i}\frac{}{x_j}`$ (on $`C_c^{\mathrm{}}\left(I^d\right)^2\left(I^d\right)`$, $`j=1,\mathrm{},d`$), by Corollary 4.2. If for each $`j`$, $`U_j\left(t\right)`$ denotes the unitary one-parameter group on $`^2\left(I^d\right)`$ generated by some self-adjoint extension $`H_j`$, then Corollary 4.2 states that $`U_j\left(t\right)`$ is induced by some unitary operator $`V_j`$ acting in the remaining variables $`(x_1,\mathrm{},x_{j1},x_{j+1},\mathrm{},x_d)`$ (i.e., with omission of the variable on the $`j`$’th place): specifically, $`U_j\left(t\right)=ind_{}^{}\left(V_j\right)`$ as a representation of $`(,+)`$; or equivalently the domain of $`H_j`$ is, for each $`j`$, given by the boundary condition
$$f(x_1\mathrm{},x_{j1},1,x_{j+1},\mathrm{},x_d)=V_j\left(f(x_1\mathrm{},x_{j1},0,x_{j+1},\mathrm{},x_d)\right).$$
(Note that the more precise interpretation of this set of boundary conditions is given in formula (4.9) of Corollary 4.2. This is the interpretation of the unitary one-parameter groups in the respective coordinate directions as *induced unitary representations* (see \[Mac53, Mac62\]), with induction $``$ for each direction.) We say that a family of self-adjoint extension operators $`H_j`$, with corresponding boundary unitaries $`V_j`$, is *quasi-commuting* if there are phase angles $`\alpha _j[0,1`$, $`j=1,\mathrm{},d`$, such that each $`V_j`$ is diagonalized by
(7.1)
$$e_{\alpha _1+n_1}^{\left(1\right)}\mathrm{}e_{\alpha _{j1}+n_{j1}}^{\left(j1\right)}e_{\alpha _{j+1}+n_{j+1}}^{\left(j+1\right)}\mathrm{}e_{\alpha _d+n_d}^{\left(d\right)}$$
as $`(n_1,\mathrm{},n_{d1},n_{d+1},\mathrm{},n_d)`$ vary over $`^{d1}`$; i.e., the lattice resulting from $`^d`$ with the $`j`$’th coordinate place omitted. It follows that the quasi-commutative case is characterized by the phase angles $`\alpha _1,\mathrm{},\alpha _d`$, and by functions $`v_j:^{d1}𝕋`$ such that, for $`n=(n_1,\mathrm{},\widehat{ȷ},\mathrm{},n_d)`$, $`v_j\left(n\right)=v_j(n_1,\mathrm{},\widehat{ȷ},\mathrm{},n_d)`$ is the eigenvalue of $`V_j`$ coresponding to the eigenvector in (7.1). (The notation $`(n_1,\mathrm{},\widehat{ȷ},\mathrm{},n_d)`$ means that the $`j`$’th place is omitted.)
###### Theorem 7.1.
Let $`\left\{H_j\right\}_{j=1}^d`$ be a family of self-adjoint extensions of the respective partials $`\frac{1}{i}\frac{}{x_j}`$ ($`j=1,\mathrm{},d`$, on $`C_c^{\mathrm{}}\left(I^d\right)^2\left(I^d\right)`$), which is assumed quasi-commutative with eigenvalue functions $`v_j(n_1,\mathrm{},\widehat{ȷ},\mathrm{},n_d)`$ from $`^{d1}`$ to $`𝕋`$. Then the extensions are commutative if and only if the following pair of cocycle conditions is satisfied for all $`j,k`$ such that $`1j<kd`$, all $`(n_1,\mathrm{},\widehat{ȷ},\mathrm{},n_d)`$, and all $`l,m\mathrm{}\left\{0\right\}`$:
(7.2)
$$\begin{array}{c}\left(v_j(n_1,\mathrm{},\widehat{ȷ},\mathrm{},n_k+l,\mathrm{},n_d)v_j(n_1,\mathrm{},\widehat{ȷ},\mathrm{},n_d)\right)\hfill \\ \hfill \times (1v_k(n_1,\mathrm{},\widehat{k},\mathrm{},n_d))=0\end{array}$$
and
(7.3)
$$\begin{array}{c}\left(v_k(n_1,\mathrm{},n_j+m,\mathrm{},\widehat{k},\mathrm{},n_d)v_k(n_1,\mathrm{},\widehat{k},\mathrm{},n_d)\right)\hfill \\ \hfill \times (1v_j(n_1,\mathrm{},\widehat{ȷ},\mathrm{},n_d))=0.\end{array}$$
###### Proof.
Since the commutativity for the one-parameter groups of unitary operators $`U_j\left(t_j\right)`$ may be stated for pairs, i.e., $`U_j\left(t_j\right)U_k\left(t_k\right)=U_k\left(t_k\right)U_j\left(t_j\right)`$, $`j<k`$, $`t_j`$, $`t_k`$, the argument for the general case $`d>2`$ is the same as for $`d=3`$. To see this, just use the formulas for the respective one-parameter groups which are analogues to (5.12)–(5.13) in the proof of Corollary 5.4. For $`d=3`$, we may introduce the leg-notation: $`v_1v_{23}`$, $`v_2v_{13}`$, $`v_3v_{12}`$. When evaluated at a general point in $`^3`$ of the form $`(k,l,m)`$, the respective eigenvalues are:
(i) $`v_{23}(l,m)`$ $`\text{ for }V_{23},`$
(ii) $`v_{13}(k,m)`$ $`\text{ for }V_{13},`$
(iii) $`v_{12}(k,l)`$ $`\text{ for }V_{12}.`$
Specifically,
(i) $`V_{23}e_{\beta +l}^{\left(2\right)}e_{\gamma +m}^{\left(3\right)}`$ $`=v_{23}(l,m)e_{\beta +l}^{\left(2\right)}e_{\gamma +m}^{\left(3\right)},`$
(ii) $`V_{13}e_{\alpha +k}^{\left(1\right)}e_{\gamma +m}^{\left(3\right)}`$ $`=v_{13}(k,m)e_{\alpha +k}^{\left(1\right)}e_{\gamma +m}^{\left(3\right)},`$
(iii) $`V_{12}e_{\alpha +k}^{\left(1\right)}e_{\beta +l}^{\left(2\right)}`$ $`=v_{12}(k,l)e_{\alpha +k}^{\left(1\right)}e_{\beta +l}^{\left(2\right)},`$
where $`\alpha ,\beta ,\gamma `$ are the fixed phase angles from the quasi-commutativity. Then the three pairs of cocycle identities from the theorem are as follows: (i a)–(i b), (ii a)–(ii b), and (iii a)–(iii b) below. The argument for the equivalence of commutativity and the cocycle identities is essentially the same as the one used in the proof of Corollary 5.4 above. The cocycle identities for $`d=3`$ are:
(i a) $`\left(v_{13}(k,m)v_{13}(k+n_1,m)\right)\left(1v_{23}(l,m)\right)`$ $`=0`$
(i b) $`\left(v_{23}(l,m)v_{23}(l+n_2,m)\right)\left(1v_{13}(k,m)\right)`$ $`=0,`$
(ii a) $`\left(v_{12}(k,l)v_{12}(k+n_1,l)\right)\left(1v_{23}(l,m)\right)`$ $`=0`$
(ii b) $`\left(v_{23}(l,m)v_{23}(l,m+n_3)\right)\left(1v_{12}(k,l)\right)`$ $`=0,`$
and
(iii a) $`\left(v_{13}(k,m)v_{13}(k,m+n_3)\right)\left(1v_{12}(k,l)\right)`$ $`=0`$
(iii b) $`\left(v_{12}(k,l)v_{12}(k,l+n_2)\right)\left(1v_{13}(k,m)\right)`$ $`=0.\text{}`$
###### Example 7.2.
*Not* all the spectral pairs in three dimensions are quasi-commutative (although this is true in $`d=2`$). Take for example the case (3.4) of Section 3 with
(7.4)
$$\mathrm{\Lambda }=\{\left(\begin{array}{c}k\\ \beta \left(k\right)+l\\ \gamma (k,l)+m\end{array}\right):k,l,m\}$$
with $`\beta :[0,1`$ and $`\gamma :^2[0,1`$ arbitrarily given functions. Then the three operators $`V_{23}`$, $`V_{13}`$ and $`V_{12}`$ are as follows:
(i) $`V_{23}=I\text{ (the identity operator in the two marked tensor slots),}`$
(ii) $`V_{13}\left(e_k^{\left(1\right)}e_{\gamma (k,l)+m}^{\left(3\right)}\right)=e^{i2\pi \beta \left(k\right)}e_k^{\left(1\right)}e_{\gamma (k,l)+m}^{\left(3\right)},`$
and
(iii) $`V_{12}\left(e_k^{\left(1\right)}e_{\beta \left(k\right)+l}^{\left(2\right)}\right)=e^{i2\pi \gamma (k,l)}e_k^{\left(1\right)}e_{\beta \left(k\right)+l}^{\left(2\right)}.`$
It follows that the three commuting unitary one-parameter groups associated with $`\mathrm{\Lambda }`$, via Theorem 2.1, are not quasi-commuting if the two functions $`\beta `$ and $`\gamma `$ in formula (7.4) are both non-constant.
###### Acknowledgements .
The authors gratefully acknowledge excellent typesetting and graphics production by Brian Treadway.
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# POST-NEWTONIAN FRAMES OF REFERENCE
## 1 Introduction
The problem of defining Post-Newtonian reference frames and developing the corresponding theory of adapted coordinate transformations has been discussed in two main papers: , and in the framework of covariant metric theories of gravity and related fields, and by in one of the so-called PPN frameworks. None of these papers is based on a general theory and they all suffer of a long-lasting confusion between frames of reference and systems of coordinates. As a consequence they rely more on improvisations suggested on the way by details of the approximate calculations than on a well-defined logical line. The first paper considers frames of reference with rotation but it does it in an oversimplified and insufficiently accurate way. The two other papers ignore the rotating frames of reference and are for this reason of restricted interest as far as the theory of such frames is concerned..
In and a general theory of frames reference and a particular post-Newtonian application was considered. This tentative theory has been subsequently modified and a few other applications have been reviewed in . This paper deals with an explicit description of this renewed theory in the framework of the post-Newtonian approximation.
The second section contains a short account of the main ingredients built in the concept of a frame of reference, namely the definition of a meta-rigid motion and a chorodesic synchronization. It includes also the definition of a particular type of frames of reference: the quo-inertial ones. The third section describes the general framework of the post-Newtonian approximation, slightly enlarged to include all types of inertial fields at the Newtonian approximation. The fourth and fifth sections deal with the main object of this paper: to use the general definition of Sect. 2 to develop a restricted but explicit theory of frames of reference as as far as the post-Newtonian approximation allows it. The Appendix contains a summary of a few basic definitions.
## 2 Frames of Reference
A frame of reference is a pair of geometrical objects intrinsically defined in the space-time being considered:
* A time-like congruence $``$ of a particular type, that we shall call the Meta-rigid motion of the frame of reference,
* A space-like foliation of a particular type $`𝒮`$, that we shall call the Chorodesic synchronization of the frame of reference.,
### 2.1 Meta-rigid motions
.
A meta-rigid motion $``$ has to possess at least three general properties which are meant to make this concept to inherit as much as possible of the formal properties of rigid motions in classical physics:
i) Their local characterization must be intrinsic and have a meaning independently of the space-time being considered.
ii) The knowledge of any open sub-bundle of the congruence must be sufficient to characterize the whole congruence.
iii) Its definition does not have to distinguish any particular world-line of the congruence.
To implement these three properties the most natural approach is to require the vector field $`u^\alpha `$ of the meta-rigid motions to satisfy a set of differential conditions. Obvious candidates which satisfy this requirement are the differential equations defining the Killing congruences <sup>1</sup><sup>1</sup>1Definitions in the Appendix.:
$$\mathrm{\Sigma }_{\alpha \beta }=0,_\alpha \mathrm{\Lambda }_\beta _\beta \mathrm{\Lambda }_\alpha =0.$$
(1)
Other congruences to be acceptable as meta-rigid motions are the Born congruences, which are a generalization of the Killing congruences and are defined by the single group of conditions:
$$\mathrm{\Sigma }_{\alpha \beta }=0.$$
(2)
We shall continue to call these congruences rigid congruences, as usual. Meta-rigid ones will be an appropriate generalization of them.
The Born conditions, and a fortiori the Killing conditions, are very restrictive in any space-time including Minkowski’s one (, ), and the question may be raised of generalizing the concept of meta-rigid motion in such a way that the family of congruences $``$ generalize Born congruences.
When facing this problem confronted with a physical application which requires the sublimated modelisation of the behavior of a rigid body, many authors resort to use Fermi congruences. By this we mean the time-like congruences which accept adapted Fermi coordinates based on a distinguished world-line seed. Because of this distinction Fermi congruences are not in general acceptable candidates to the concept of a meta-rigid motion.
Harmonic congruences can be defined as those time-like congruences, with unit tangent vector $`u^\alpha `$, for which there exist three space-like functions $`f^a(x^\alpha )`$ satisfying the following equations:
$$\mathrm{}f^a=0,u^\alpha _\alpha f^a=0,a,b,c=1,2,3$$
(3)
where $`\mathrm{}`$ is the d’Alembertian operator corresponding to the space-time metric. Harmonic congruences, with the corresponding harmonic coordinates, have been used extensively in the literature for several reasons, including mainly the fact that they simplify many calculations. Harmonic congruences could in principle be considered good representatives of meta-rigid motions. In fact they are not. They are acceptable generalizations of Killing congruences but they are not generalizations of the Born congruences. More precisely it has been proved that irrotational, and not Killing, Born congruences are never harmonic congruences.
To solve the problem just mentioned concerning the harmonic congruences it has been proposed to identify the meta-rigid motions with appropriately selected Quo-harmonic congruences (, , ). These being those congruences for which there exist three independent space-like solutions $`f^a(x^\alpha )`$ satisfying the following equations:
$$\widehat{\mathrm{}}f^a(\mathrm{}+\mathrm{\Lambda }^\alpha _\alpha )f^a=0,u^\alpha _\alpha f^a=0,a,b,c=1,2,3$$
(4)
or equivalently, using adapted coordinates and the notations of the Appendix:
$$\widehat{\mathrm{}}f^a=\widehat{g}^{ij}(\widehat{}_i\widehat{}_j\widehat{\mathrm{\Gamma }}_{ij}^k\widehat{}_k)f^a=0,f^a=f^a(x^i)$$
(5)
These are proper generalizations of Born congruences in the sense that any Born congruence is also a quo-harmonic one. In fact when using adapted coordinates the operator $`\widehat{\mathrm{}}`$ in the Eq. above becomes the usual Laplacian operator corresponding to a 3-dimensional Riemannian metric.
Besides the above mentioned general local conditions to be imposed to the concept of Meta-rigid motions global specific conditions will be necessary depending on the particular global structure of the space-time being considered. Also particular but important will be the case, where the congruence contains a time-like geodesic and it make sense, physically, to implement the principle of geodesic equivalence.
Notice that quo-harmonic congruences can be described and studied using any system of coordinates. But quo-harmonic coordinates, i.e. any system of three independent functions $`f^a`$ of Eq. 4, are in general the most convenient ones to use.
### 2.2 Chorodesic synchronizations
.
The choice of a foliation to define the synchronization time of a frame of reference is to some extent less essential than the choice of a quo-harmonic congruence to define its meta-rigid motion. Besides, time is a much better understood concept in General relativity than the concept of space. The role of a synchronization is to be a first step towards a convenient definition of a universal scale of time, and the properties that are to be required from a synchronization will depend on the use for which it is intended. Atomic and sidereal time are two scales of common use which correspond to different synchronizations of the frame of reference co-moving with the Earth. Not to mention many other scales of time used in astronomy. On the contrary, for instance, it would not make sense to use different types of congruences to describe the motion of the Earth, at the approximation which considers it as rigid.
We consider here briefly how the International Atomic Time (IAT) scale is defined. It is based on the definition of the second in the international system of units SI as a duration derived from the frequency of a particular atomic transition. This definition is universal, in the sense that any physicist is able to use a well defined second, but at the same time is local because the identification of atomic time with proper-time duration implies that this definition of the second can be only used to define a scale of time along the world-line of the clock being used as a standard reference.
To define a scale of time on Earth to be used on navigation systems like the GPS for instance requires a different approach. In this case <sup>2</sup><sup>2</sup>2See for instance the second is defined as above with some precisions added. The resulting definition of IAT is:
IAT is a coordinate time scale defined in a geocentric reference frame with an SI second realized on the rotating geoid. <sup>3</sup><sup>3</sup>3The geoid is a level surface of constant geo-potential (gravity and centrifugal potential) at “mean sea level”, extended below the continents.
It follows from this that the operational definition of the second is relative to some particular locations, i.e. world-lines of a particular congruence, which in this case <sup>4</sup><sup>4</sup>4If the influence of the Sun and the Moon is neglected is the rotating Killing congruence corresponding to the frame of reference co-moving with the Earth. A second at some other location is then defined as to nullify the relativistic red-shift between any pair of clocks of reference co-moving with the Earth. Therefore the second at any other location, not in the geoid, is then the interval of time separating the arrival of two light signals sent, one second apart, from a standard clock located on the geoid. The red-shift formulas of General relativity allow then to compare the rhythm of any two clocks that can be joined with light signals.
Let $``$ be any time-like congruence. A Chorodesic $`C`$ of $``$ is by definition, , a space-like line such that its tangent vector $`p^\alpha `$ satisfy the following equations:
$$\frac{p^\alpha }{d\lambda }=\frac{1}{2}u^\alpha \mathrm{\Sigma }_{\mu \nu }p^\mu p^\nu ,p^\alpha =\frac{dx^\alpha }{d\lambda }$$
(6)
where $`\lambda `$ is the proper length along $`C`$. Obviously if $`\mathrm{\Sigma }_{\mu \nu }=0`$, i.e. if $``$ is a Born congruence then any chorodesic of $``$ is also a geodesic of the space-time.
Chorodesics of a congruence are important mainly because of the following result:
If a space-like chorodesic $`C`$ is orthogonal to a world-line of a congruence $``$ then it is orthogonal to all the world-lines of the congruence $``$ that it crosses.
This follows from differentiating the scalar product $`p^\rho u_\rho `$ along $`C`$. We thus get:
$$\frac{d}{d\lambda }(p^\rho u_\rho )=\frac{1}{2}\mathrm{\Sigma }_{\mu \nu }p^\mu p^\nu +\frac{1}{2}p^\rho p^\sigma (_\rho u_\sigma +_\sigma u_\rho )$$
(7)
and using the definition of $`\mathrm{\Sigma }_{\mu \nu }`$ this is equivalent to:
$$\frac{d}{d\lambda }(p^\rho u_\rho )=2(p^\rho u_\rho )(p^\rho \mathrm{\Lambda }_\rho )$$
(8)
from where the statement above follows.
Particular foliations associated with a congruence $``$ are the one parameter family of hyper-surfaces generated by the chorodesics orthogonal to any particular world-line of the congruence. We call these foliations Chorodesic synchronizations. In general a chorodesic synchronization depends on a world-line seed, except if the congruence is integrable in which case all its chorodesic synchronizations coincide with the family of hyper-surfaces orthogonal to the congruence.
Let $`𝒮`$ be the chorodesic synchronization orthogonal to some world-line $`W`$ of a frame of reference $``$. Then the associated Atomic (or Proper) Time Coordination (ATC) scale is by definition a time coordinate which value at any event $`E`$ is the proper time interval between some arbitrary event on $`W`$ and the intersection of $`W`$ with the leave of $`𝒮`$ which contains $`E`$.
Notice that, even when the congruence of reference $``$ is hypersurface orthogonal, i.e. there exists a single chorodesic synchronization, the ATC may depend on the particular world-line $`W`$ which is used to identify the coordinate time with proper-time along $`W`$.
Switching from a world-line seed, $`W`$, of an ATC to another $`W^{}`$ gives rise to a sub-group of adapted coordinate transformations <sup>5</sup><sup>5</sup>5An analytical example of such sub-groups can be seen in . A complete definition of the ATC requires therefore to choose a particular world-line $`𝒮`$, or eventually a bunch of equivalent ones. On Earth the ATC, i.e. the IAT, is associated to the bunch of world-lines $`𝒮`$ of locations on the geoid.
The interest of the consideration of chorodesic synchronizations stems from the fact that they generalize the hyper-surfaces orthogonal to integrable congruences, as well as being a particular case of synchronizations of equal cyclic adapted time for both static and stationary space-times.
## 3 Inertial and gravitational fields: The Post-Newtonian framework
Our general framework will consist of the following items:
* A domain $`D`$ of $`R_4`$ to be called the primary domain, or global domain if global conditions are required.
* A four dimensional space-time metric with signature $`+2`$, $`g_{\alpha \beta }`$, which can be expanded in power of $`1/c`$ as follows:
$`g_{00}`$ $`=`$ $`1+{\displaystyle \frac{1}{c^2}}f_{00}+{\displaystyle \frac{1}{c^4}}h_{00}`$ (9)
$`g_{0j}`$ $`=`$ $`{\displaystyle \frac{1}{c}}f_{0j}+{\displaystyle \frac{1}{c^3}}h_{0j}`$ (10)
$`g_{ij}`$ $`=`$ $`\delta _{ij}+{\displaystyle \frac{1}{c^2}}h_{ij}`$ (11)
The corresponding developments of $`\xi ,\phi _i`$, and the Fermat quo-tensor $`\widehat{g}_{ij}`$ are:
$`\xi `$ $`=`$ $`1{\displaystyle \frac{1}{2c^2}}f_{00}{\displaystyle \frac{1}{2c^4}}(h_{00}+{\displaystyle \frac{1}{4}}f_{00}^2)`$ (12)
$`\phi _i`$ $`=`$ $`{\displaystyle \frac{1}{c}}f_{0i}+{\displaystyle \frac{1}{c^3}}(h_{0i}+f_{0i}f_{00})`$ (13)
$`\widehat{g}_{ij}`$ $`=`$ $`\delta _{ij}+{\displaystyle \frac{1}{c^2}}\alpha _{ij}\alpha _{ij}=h_{ij}+f_{0i}f_{0j}`$ (14)
Every further term will be systematically neglected.
Our general requirements on both $`D`$ and $`g_{\alpha \beta }`$ are the following:
* The time-like congruence $``$ defined by the parametric equations $`x^i=const.`$ is, at the required approximation, a quo-harmonic congruence and the $`x^i`$ are the space quo-harmonic coordinates adapted to $``$. That is to say we assume that we have:
$$_i\alpha _j^i\frac{1}{2}_j\alpha =0,\alpha =\delta ^{ij}\alpha _{ij}$$
(15)
This follows from making explicit that the functions $`x^i`$ themselves are solutions of Eqs. 5.
* The foliation $`𝒮`$ defined by the hyper-surfaces $`t=const.`$ is the chorodesic synchronization of $``$ with seed on a world-line $`W`$ with parametric equations $`x^i=x_0^i=const.`$. To see what this means at the required approximation let us follow the steps below:
+ Since by definition (See Sect. 2) on $`W`$ $`t`$ is its proper time we have:
$$g_{00}1$$
(16)
where here and in the remaining of this paragraph $``$ means that the corresponding relation is satisfied on $`W`$.
+ Since by construction the chorodesics at any event $`P`$ of $`W`$ remain on a hypersuface $`x^0=const`$, their tangent vector have a zero time-component, $`p^0=0`$. On the other hand the unit tangent vector at $`P`$ to $`W`$, with components $`u^0=1,u^i=0`$, is orthogonal to every chorodesic through $`P`$ therefore we shall have:
$$g_{0i}0$$
(17)
+ From $`p^0=0`$ all along the path of every chorodesic passing through $`P`$ and from 6 it follows also that:
$$\frac{d}{d\lambda }p^00\text{or}\mathrm{\Gamma }_{ij}^0p^ip^j0$$
(18)
Using the explicit expressions of $`\mathrm{\Gamma }_{ij}^0`$ and $`\mathrm{\Sigma }_{ij}`$ as functionals of $`g_{\alpha \beta }`$, and taking into account that the relations above have to hold whatever the values of $`p^i`$, we readily find that:
$$_ig_{0j}+_jg_{0i}0$$
(19)
This process could be continued indefinitely to prove that whatever the number of derivatives one has:
$$_{(ijk\mathrm{}}g_{s)0}0$$
(20)
the round brackets meaning complete symmetrization. But to implement this construction at the post-Newtonian level only the three first steps are necessary.
* $`f_{0i}`$ in 9 will be required to be a linear function of $`x^i`$ so that the rotation rate of $``$ at the order $`1/c`$ be always a function of time only. This will allows to include in our framework the general theory of Newtonian frames of reference.
$$_if_{0j}_jf_{0i}=\omega _{ij}(t)$$
(21)
This general framework will be called The post-Newtonian formalism. It can be further restricted to deal with more specific situations corresponding to two broad categories:
* The first one is that for which the Riemann tensor of the space-time metric is zero at the appropriate approximation, this meaning here that:
$$R_{0i0j}=O(c^6),R_{0ijk}=O(c^5),R_{ijkl}=O(c^4)$$
(22)
in which case 9 is an approximation to an inertial field.
* the second broad category includes all cases for which 22 does not hold, and among them all metric theories of gravitation, alone or coupled to other fields. No use whatsoever of specific field equations will be made in this paper that will cover also therefore any appropriate parameterized formalism.
Two particular cases have to be mentioned because the general framework starts from them. They are the following:
* Inertial and gravitational fields at the Newtonian approximation can be described assuming that:
$$h_{\alpha \beta }=0$$
(23)
* An extended Newtonian theory of gravitation that we discussed in can be described assuming that:
$$h_{00}=0,h_{ij}=0$$
(24)
Because in both cases $`h_{ij}=0`$ these two approximate <sup>6</sup><sup>6</sup>6In both theories where described as exact. Newton’s theory as an affine theory and the extended theory as a semi-metric one. frameworks can be said, in a generalized sense, to be invariant under the group of rigid motions:
$$t=t,x^i=s^i(t)+R_j^i(t)y^i$$
(25)
where $`R_j^i(t)`$ is a time-dependent rotation matrix. It is this family of transformations that we generalize in the two next sections.
At the post-Newtonian approximation the Newtonian field is:
$`\mathrm{\Lambda }_i`$ $`=`$ $`\lambda _i+{\displaystyle \frac{1}{c^2}}\mathrm{\Xi }_i`$ (26)
$`\lambda _i`$ $`=`$ $`{\displaystyle \frac{1}{2}}_if_{00}_tf_{0i}`$ (27)
$`\mathrm{\Xi }_i`$ $`=`$ $`({\displaystyle \frac{1}{2}}_ih_{00}+{\displaystyle \frac{3}{2}}_tf_{00}f_{0i}+_th_{0i}f_{00}\lambda _i)`$ (28)
The Coriolis, or rate of rotation field, is:
$`\mathrm{\Omega }_{ij}`$ $`=`$ $`\omega _{ij}+{\displaystyle \frac{1}{c^2}}\mathrm{\Psi }_{ij}`$ (29)
$`\omega _{ij}`$ $`=`$ $`_if_{0j}_jf_{0i}`$ (30)
$`\mathrm{\Psi }_{ij}`$ $`=`$ $`_ih_{0j}_jh_{0i}+{\displaystyle \frac{1}{2}}f_{00}\omega _{ij}`$ (31)
$`(2\lambda _i+3_tf_{0i})f_{0j}+(2\lambda _j+3_tf_{0j})f_{0i}`$
And the so-called rate of deformation field is:
$$\mathrm{\Sigma }_{ij}=\frac{1}{c^2}_t\alpha _{ij}$$
(32)
## 4 Post-Newtonian quo-harmonic coordinate transformations
We assume in this section that the potentials 9 considered in Sect. 2 are defined in the appropriate global domain and are referred to a frame of reference of a particular type that we shall call quo-inertial. By definition this type is defined as equivalent to the class of frames of reference satisfying the following conditions:
* The quo-harmonic congruence defining the meta-rigid motion of the frame of reference contains one or many time-like geodesics.
* The restricted Principle of geodesic equivalence can be implemented on at least one of these geodesics, say $`𝒞`$, to be called the center of motion of the frame of reference. This meaning that a system of global adapted quo-harmonic coordinates can be found such that on this geodesic the Fermat tensor is the unit tensor, and the Zel’ manov-Cattaneo connection is zero.
$$\widehat{g}_{ij}(x_0^s)=\delta _{ij},\widehat{\mathrm{\Gamma }}_{jk}^i(x_0^s)=0$$
(33)
$`x_0^i`$ being the coordinates of $`𝒞`$. As we shall see this is locally always possible. It is an assumption to say that a globally well-defined quo-harmonic congruence has this property.
We consider below the construction of any other quo-inertial frame of reference as a second order approximate development around its center world-line. This construction can be split in two steps:
* The first step defines the meta-rigid motion by giving the parameterized equations of a quo-harmonic congruence referring this congruence to a system of quo-harmonic coordinates.
* The second step constructs the chorodesic synchronization seeded on the center geodesic of the quo-inertial frame of reference.
As already mentioned in Sect. 2, these two steps have very different physical meanings and geometrical interpretations. The first step has to do with a modeling of an idealized rigid body with its motion being referred to the global frame of reference. The second step in making a particular choice of a possible time-distribution protocol. Both are necessary to be able to define physically meaningful tensor transformations between the two frames of reference.
Let us consider a geodesic with parameterized equations:
$$x^i=s^i(t)$$
(34)
To implement the construction of a quo-inertial meta-rigid motion centered on this geodesic we consider the congruence defined by equations:
$`t`$ $`=`$ $`t`$ (35)
$`x^i`$ $`=`$ $`s^i(t)+R_p^i(t)y^p+{\displaystyle \frac{1}{c^2}}R_p^i(t)\delta _2y^p`$ (36)
$`\delta _2y^p(t)`$ $`=`$ $`L_r^p(t)y^r+{\displaystyle \frac{1}{2}}Q_{rk}^p(t)y^ry^k+{\displaystyle \frac{1}{3!}}C_{rkl}^p(t)y^ry^ky^l`$ (37)
where $`R_p^i(t)`$ is a time dependent rotation matrix and $`y^i`$ can equivalently be considered either as parameters or as a system of adapted coordinates of the new congruence. Obviously the zero order choice in these expressions has been chosen to guarantee the correct Newtonian limit. Also to simplify the writing we have chosen the coordinates $`y^i`$ such that the parametric equations of the geodesic $`𝒞`$ of the new congruence be $`y^i=0`$.
Differentiating the expressions above we obtain:
$$dt=dt,dx^i=W^idt+R_j^idy^j+\frac{1}{c^2}(A_k^idy^k+B^idt),$$
(38)
where:
$`W^i`$ $`=`$ $`v^i+\dot{R}_p^iy^p`$ (39)
$`A_j^i`$ $`=`$ $`R_p^i(L_j^p+Q_{jk}^py^k+{\displaystyle \frac{1}{2}}C_{jkl}^py^ky^l)`$ (40)
$`B^i`$ $`=`$ $`R_p^i(\dot{L}_r^py^r+{\displaystyle \frac{1}{2}}\dot{Q}_{rk}^py^ry^k+{\displaystyle \frac{1}{3!}}\dot{C}_{rkl}^py^ry^ky^l)`$ (41)
where $`v^i=\dot{s}^i`$, the over-head dot meaning a derivative with respect to $`t`$.
A straightforward calculation gives the transformed basic metric quantities:
$`\overline{g}_{00}`$ $`=`$ $`1+{\displaystyle \frac{1}{c^2}}\overline{f}_{00}+{\displaystyle \frac{1}{c^4}}\overline{h}_{00}`$ (42)
$`\overline{f}_{00}`$ $`=`$ $`f_{00}+2f_{0i}W^i+W^2`$ (43)
$`\overline{h}_{00}`$ $`=`$ $`h_{00}+2h_{0i}W^i+2(f_{0i}+W_i)B^i+h_{ij}W^iW^j+h_{00}^{}`$ (44)
$`h_{00}^{}`$ $`=`$ $`+(_kf_{00}+\omega _{ki}R_p^kW^i)\delta _2y^p`$ (45)
$`\overline{g}_{0i}`$ $`=`$ $`{\displaystyle \frac{1}{c}}\overline{f}_{0i}+{\displaystyle \frac{1}{c^3}}\overline{h}_{0i}`$ (46)
$`\overline{f}_{0i}`$ $`=`$ $`(f_{0s}+W_s)R_i^s`$ (47)
$`\overline{h}_{0i}`$ $`=`$ $`(h_{0s}+B_s+h_{sk}W^k)R_i^s+(f_{0k}+W_k)A_i^k+h_{0i}^{}`$ (48)
$`h_{0i}^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\omega _{rs}R_k^rR_i^s\delta _2y^k`$ (49)
$`\overline{g}_{ij}`$ $`=`$ $`\delta _{ij}+{\displaystyle \frac{1}{c^2}}\overline{h}_{ij}`$ (50)
$`\overline{h}_{ij}`$ $`=`$ $`h_{rs}R_i^rR_j^s+\delta _{rs}(A_i^rR_j^s+A_j^rR_i^s),`$ (51)
$`\widehat{\overline{g}}_{ij}`$ $`=`$ $`\delta _{ij}+{\displaystyle \frac{1}{c^2}}\overline{\alpha }_{ij}`$ (52)
$`\overline{\alpha }_{ij}`$ $`=`$ $`\overline{h}_{ij}+\overline{f}_{0i}\overline{f}_{0j}`$ (53)
the starred terms in 42 and 46 coming from substituting $`x^i`$ in $`\overline{f}_{00}`$ and $`\overline{f}_{0i}`$ according to 35.
The congruence defined by the parametric Eqs. 35 will be a quo-harmonic one if we choose the free functions of time $`L_j^i(t),Q_{rk}^p(t)`$ and $`C_{rkl}^p(t)`$ in such way that:
$$\frac{}{y^i}\overline{\alpha }_j^i\frac{1}{2}\frac{}{y^j}\overline{\alpha }=0$$
(54)
i.e., in such a way that the parameters $`y^i`$ are quo-harmonic coordinates.
The conditions stated in 33 are further restrictions on these conditions. We shall rewrite them at the appropriate approximation as:
$$\alpha _{ij}0,\frac{}{y^k}\overline{\alpha }_{ij}0$$
(55)
where from now on the symbol $``$ will mean that the coordinates $`x^i`$ have been replaced by $`s^i(t)`$ and $`y^k`$ by $`0`$.
Taking into account 46 \- 52 the first group of conditions require that:
$$L_{(ij)}\frac{1}{2}(h_{rs}+(f_{0s}+v_s)(f_{0r}+v_r))R_i^sR_j^r$$
(56)
where:
$$L_{ij}=\delta _{ik}L_j^k,L_{(ij)}=\frac{1}{2}(L_{ij}+L_{ji})$$
(57)
We shall see at the end of this section that the antisymmetric part of $`L_{ij}`$ can be chosen as to give an unambiguous physical meaning to the rotation matrix $`R_j^i`$.
The second group of conditions 33 leads to:
$$Q_{j,ik}+Q_{i,jk}+X_{k,ij}0$$
(58)
or:
$$Q_{i,jk}=\frac{1}{2}(X_{j,ki}+X_{k,ji}X_{i,jk})$$
(59)
where:
$`X_{k,ij}`$ $`=`$ $`[_lh_{rs}+{\displaystyle \frac{1}{2}}(\omega _{ls}(f_{0r}+v_r)+(f_{0s}+v_s)\omega _{lr}]R_k^lR_i^rR_j^s`$ (61)
$`+{\displaystyle \frac{1}{2}}(\mathrm{\Delta }_{ki}(f_{0r}+v_r)R_j^r+\mathrm{\Delta }_{kj}(f_{0s}+v_s)R_i^r)`$
and:
$$\mathrm{\Delta }_{ij}=2\dot{R}_i^sR_{sj}$$
(62)
Notice that the conditions 55 could have been implemented on any world-line. But it is only because we have assumed that the world-line $`x^i=s^i(t)`$ is a geodesic that these conditions are physically justified and required by the principle of geodesic equivalence.
Consider now the next step in the approximation. Using again 46 \- 52 we have:
$$\frac{^2}{y^ky^l}\overline{\alpha }_{ij}Y_{ij,kl}+C_{i,jkl}+C_{j,ikl}$$
(63)
where:
$`Y_{ij,kl}`$ $`=`$ $`_{mn}^2\alpha _{rs}R_i^rR_j^sR_k^mR_l^n`$ (65)
$`+{\displaystyle \frac{1}{4}}(\overline{\omega }_{ki}\overline{\omega }_{lj}+\overline{\omega }_{kj}\overline{\omega }_{li}){\displaystyle \frac{1}{4}}(\stackrel{~}{\omega }_{ki}\stackrel{~}{\omega }_{lj}+\stackrel{~}{\omega }_{kj}\stackrel{~}{\omega }_{li})`$
where:
$$\overline{\omega }_{ij}=\stackrel{~}{\omega }_{ij}+\mathrm{\Delta }_{ij},\stackrel{~}{\omega }_{ij}=\omega _{rs}R_i^rR_j^s$$
(66)
At the order $`1/c^2`$ we have:
$$\widehat{\overline{R}}_{ijkl}\widehat{\stackrel{~}{R}}_{ijkl}+\frac{3}{4}(\overline{\omega }_{ij}\overline{\omega }_{kl}\stackrel{~}{\omega }_{ij}\stackrel{~}{\omega }_{kl})$$
(67)
where:
$`\widehat{\stackrel{~}{R}}_{ijkl}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(_{mn}^2\alpha _{rs}+_{rs}^2\alpha _{mn}_{mr}^2\alpha _{ns}_{ns}^2\alpha _{mr})R_i^rR_j^sR_k^mR_l^n`$ (68)
$`\widehat{\overline{R}}_{ijkl}`$ $`=`$ $`{\displaystyle \frac{1}{2}}({\displaystyle \frac{^2}{y^iy^k}}\overline{\alpha }_{jl}+{\displaystyle \frac{^2}{y^jy^l}}\overline{\alpha }_{ik}{\displaystyle \frac{^2}{y^iy^l}}\overline{\alpha }_{jk}{\displaystyle \frac{^2}{y^jy^k}}\overline{\alpha }_{il})`$ (69)
Since 67 is independent of $`C_{jkl}^i`$ and $`\mathrm{\Delta }_{ij}`$ is an arbitrary skew-symmetric quo-tensor, it is not in general possible to nullify the second derivatives of the Fermat tensor on the center of a quo-inertial motion.
On the contrary it is always possible to implement the second order quo-harmonic condition:
$$\frac{}{y^k}(\frac{}{y^i}\overline{\alpha }_j^i\frac{1}{2}\frac{}{y^j}\overline{\alpha })0$$
(70)
In fact from 63 and 65 it follows that the preceding condition is equivalent to:
$$C_{ijkl}\delta ^{kl}\frac{1}{2}(\overline{\omega }_{ik}\overline{\omega }_{jl}\stackrel{~}{\omega }_{ik}\stackrel{~}{\omega }_{jl})\delta ^{kl}$$
(71)
The solution of this equation compatible with the symmetry of the problem is:
$$C_{ijk}^s=\frac{1}{4}(\overline{\mathrm{\Theta }}_i^s\overline{\delta }_{jk}+\overline{\mathrm{\Theta }}_j^s\overline{\delta }_{ki}+\overline{\mathrm{\Theta }}_k^s\overline{\delta }_{ij})\frac{1}{4}(\stackrel{~}{\mathrm{\Theta }}_i^s\stackrel{~}{\delta }_{jk}+\stackrel{~}{\mathrm{\Theta }}_j^s\stackrel{~}{\delta }_{ki}+\stackrel{~}{\mathrm{\Theta }}_k^s\stackrel{~}{\delta }_{ij})$$
(72)
where:
$$\overline{\mathrm{\Theta }}_{ij}=\frac{1}{2}ϵ^k\overline{\omega }_{ik}\overline{\omega }_{kj},\stackrel{~}{\mathrm{\Theta }}_{ij}=\frac{1}{2}ϵ^k\stackrel{~}{\omega }_{ik}\stackrel{~}{\omega }_{kj},ϵ^k=1$$
(73)
$$\overline{\delta }_{jk}=2\overline{\mathrm{\Theta }}_{jk}/(\overline{\omega }_{rs}\overline{\omega }^{rs})$$
(74)
$$\stackrel{~}{\delta }_{jk}=2\stackrel{~}{\mathrm{\Theta }}_{jk}/(\stackrel{~}{\omega }_{rs}\stackrel{~}{\omega }^{rs})$$
(75)
To understand the meaning of this solution let us look at the coordinate transformations 35 as the composition of two transformations:
$`x^i`$ $`=`$ $`s^p+R_p^i(y^p+{\displaystyle \frac{1}{c^2}}(L_r^p(t)y^r+{\displaystyle \frac{1}{2}}Q_{rk}^p(t)y^ry^k)+{\displaystyle \frac{1}{3!}}C_{jkl}^py^jy^ky^l))`$ (76)
$`y^i`$ $`=`$ $`R_p^{\prime \prime i}(y^p+{\displaystyle \frac{1}{3!c^2}}C_{jkl}^{\prime \prime p}y^jy^ky^l)`$ (77)
with:
$$R_k^iR_j^{\prime \prime k}=R_j^i$$
(78)
which leads to the composition law:
$$C_{jkl}^i=C_{jkl}^{\prime \prime i}+\stackrel{~}{R}_p^iC_{mnr}^pR_j^mR_k^nR_l^r$$
(79)
If we choose $`R_k^i`$ such that:
$$2\dot{R}_i^k=\omega _{mk}R_i^m$$
(80)
i.e. if choose $`R_k^i`$ such that, according to 62 and 66, the rotation rate of the congruence defined by the first group of equations 76 is zero, then from 72 we have:
$$C_{jkl}^i=\frac{1}{4}(\mathrm{\Theta }_i^s\delta _{jk}^{}+\mathrm{\Theta }_j^s\delta _{ki}^{}+\mathrm{\Theta }_k^s\delta _{ij}^{})$$
(81)
with:
$$\mathrm{\Theta }_{rs}^{}=\frac{1}{2}ϵ^k\omega _{ik}^{}\omega _{kj}^{},\delta _{jk}^{}=2\mathrm{\Theta }_{jk}^{}/(\omega _{rs}^{}\omega ^{rs}),\omega _{rs}^{}=\omega _{ij}R_r^iR_s^j$$
(82)
this being the unique acceptable solution of 71 if $`\overline{\omega }_{ij}^{}=0`$, because one has to take into account the complete symmetry of $`C_{jkl}^i`$ with respect to its covariant indices and the facts that $`\mathrm{\Theta }_{ij}^{}`$ is symmetric and that the dual of $`\omega _{ij}^{}`$ is an eigen-vector with eigen-value zero.
Using again 71 and 72 with $`\stackrel{~}{\omega }_{ij}^{}=0`$, we get:
$$C_{ijk}^{\prime \prime s}=\frac{1}{4}(\overline{\mathrm{\Theta }}_i^s\overline{\delta }_{jk}^{}+\overline{\mathrm{\Theta }}_j^s\overline{\delta }_{ki}^{}+\overline{\mathrm{\Theta }}_k^s\overline{\delta }_{ij}^{})$$
(83)
Combining 78, 79, 82 and 83 we uncover the meaning of $`R_j^i`$ and $`R_j^{\prime \prime i}`$. The first is the rotation leading to the non rotating congruence centered on the geodesic of the transformed congruence; the second is the rotation of the latter congruence as measured from the non rotating reference defined by the first one.
To determine the skew-symmetric part of $`L_{ij}`$ we shall demand on the center geodesic $`C`$ of the new congruence that:
$$\overline{\mathrm{\Psi }}_{ij}0$$
(84)
i.e. that the post-Newtonian correction to $`\omega _{ij}`$ on $`C`$ be zero. This requirement does not have to be considered as an additional demand to the quo-harmonic condition but as a precision made to the meaning of the matrix $`R_j^i`$ in the transformations 35. Equivalently we could say that the real parameter in these transformations is not $`R_j^i`$ but the rotation rate $`\overline{\omega }_{ij}`$, on $`𝒞`$, of the new congruence. In fact considering the latter as a free parameter it does not make sense to consider separately its classical and relativistic correction. It has to be considered as a whole and have the meaning that it has at the classical approximation.
From 29 and from 39 \- 50 it follows that $`\overline{\mathrm{\Psi }}_{ij}`$ can be written as:
$$\overline{\mathrm{\Psi }}_{ij}\dot{L}_{[ij]}+\overline{\omega }_{ik}L_j^k\overline{\omega }_{jk}L_i^k+Z_{ik},L_{[ij]}=\frac{1}{2}(L_{ij}L_{ji})$$
(85)
where both the symmetric part of $`L_{ij}`$, as given by 56, and the skew-symmetric object $`Z_{ij}`$ depend only on quantities that have been already calculated. It follows that the requirement 84 is equivalent to a differential equation for $`L_{[ij]}`$ whose general solution will depend on its value at a given instant. These constants can be considered to be zero without any loss of generality because whatever their value they will all lead to the same congruence.
## 5 Post-Newtonian chorodesic synchronizations
We consider in this section a chorodesic foliation $``$ orthogonal to the center geodesic $`C`$ of the meta-rigid motion we are considering, i.e. the family of hyper-surfaces spanned by the chorodesics orthogonal to this geodesic. The corresponding IAT will be a time coordinate $`t^{}=t^{}(t,y^i)`$ such that $`t^{}=const`$ be the local equations of the chorodesic leaves of the foliation.
Let $`t^{}`$ be defined as the inverse of the following time transformation:
$`t`$ $`=`$ $`t^{}+{\displaystyle \frac{1}{c^2}}\delta _2t^{}+{\displaystyle \frac{1}{c^4}}\delta _4t^{}`$ (86)
$`\delta _2t^{}`$ $`=`$ $`I(t^{})+L_i(t^{})y^i+{\displaystyle \frac{1}{2}}Q_{ij}(t^{})y^iy^j`$ (87)
$`\delta _4t^{}`$ $`=`$ $`H(t^{})+K_i(t^{})y^i+{\displaystyle \frac{1}{2}}P_{ij}(t^{})y^iy^j`$ (88)
Differentiating this expression we get:
$$dt=dt^{}+\frac{1}{c^2}(Ddt^{}+C_idy^i)+\frac{1}{c^4}(Fdt^{}+D_idy^i)$$
(89)
where:
$`D`$ $`=`$ $`\dot{I}^{}+\dot{L}_i^{}y^i+{\displaystyle \frac{1}{2}}\dot{Q}_{ij}^{}y^iy^j`$ (90)
$`C_i`$ $`=`$ $`L_i+Q_{ij}y^j`$ (91)
$`F`$ $`=`$ $`\dot{H}^{}+\dot{K}_i^{}y^i+{\displaystyle \frac{1}{2}}\dot{P}_{ij}^{}y^iy^j`$ (92)
$`E_i`$ $`=`$ $`K_i+P_{ij}y^j`$ (93)
where a dotted and primed quantity means that it has been differentiated with respect to $`t^{}`$. With these notations the new potentials are:
$`g_{00}^{}`$ $`=`$ $`1+{\displaystyle \frac{1}{c^2}}f_{00}^{}+{\displaystyle \frac{1}{c^4}}h_{00}^{}`$ (94)
$`f_{00}^{}`$ $`=`$ $`\overline{f}_{00}2D`$ (95)
$`h_{00}^{}`$ $`=`$ $`\overline{h}_{00}+2\overline{f}_{00}D2FD^2+h_{00}^{}`$ (96)
$`h_{00}^{}`$ $`=`$ $`_t\overline{f}_{00}\delta _2t^{}`$ (97)
$`g_{0i}^{}`$ $`=`$ $`{\displaystyle \frac{1}{c}}f_{0i}^{}+{\displaystyle \frac{1}{c^3}}h_{0i}^{}`$ (98)
$`f_{0i}^{}`$ $`=`$ $`\overline{f}_{0i}C_i`$ (99)
$`h_{0i}^{}`$ $`=`$ $`\overline{h}_{0i}+\overline{f}_{00}C_i+\overline{f}_{0i}DE_i+h_{0i}^{}`$ (100)
$`h_{0i}^{}`$ $`=`$ $`_t\overline{f}_{0i}\delta _2t^{}`$ (101)
the starred terms in $`h_{00}^{}`$ and $`h_{0i}^{}`$ coming from using 86 in the lower order terms $`f_{00}^{}`$ and $`f_{0i}^{}`$. And also:
$`g_{ij}^{}`$ $`=`$ $`\delta _{ij}+{\displaystyle \frac{1}{c^2}}h_{ij}^{}`$ (102)
$`h_{ij}^{}`$ $`=`$ $`\overline{h}_{ij}+\overline{f}_{0i}C_j+\overline{f}_{0j}C_iC_iC_j`$ (103)
and:
$$\widehat{g}_{ij}^{}=\delta _{ij}+\frac{1}{c^2}\alpha _{ij}^{},\alpha _{ij}^{}=\overline{\alpha }_{ij}$$
(104)
this last result following directly from the fact that the Fermat object, $`\widehat{g}_{ij}`$, is a quo-tensor, i.e. a three dimensional tensor under coordinate transformations that leave invariant the congruence.
The conditions to be demanded to 86 were derived in Sect. 3. Taking into account 94, 16 requires that:
$`\dot{I}^{}`$ $``$ $`{\displaystyle \frac{1}{2}}\overline{f}_{00}`$ (105)
$`\dot{H}^{}`$ $``$ $`{\displaystyle \frac{1}{2}}\overline{h}_{00}+{\displaystyle \frac{3}{4}}\overline{f}_{00}^2`$ (106)
The functions $`I`$ and $`H`$ have to be obtained integrating these differential equations and each will depend on an arbitrary constant. The second group of conditions 17 and 98 require that:
$`L_i`$ $``$ $`\overline{f}_{0i}`$ (107)
$`K_i`$ $``$ $`\overline{h}_{0i}+{\displaystyle \frac{1}{2}}\overline{f}_{00}(f_{0i}+L_i)`$ (108)
Finally the third group of conditions 19 and the relations obtained differentiating 98 with respect to $`y^i`$ yield:
$`Q_{ij}`$ $``$ $`{\displaystyle \frac{1}{2}}(\overline{}_i\overline{f}_{0j}+\overline{}_j\overline{f}_{0i})`$ (109)
$`P_{ij}`$ $``$ $`{\displaystyle \frac{1}{2}}(N_{ij}+N_{ji})`$ (110)
with:
$$N_{ij}=\overline{}_i\overline{h}_{0j}+\overline{}_i\overline{f}_{00}L_i+\frac{1}{2}\overline{f}_{00}Q_{ij}+\dot{L}_i^{}(\overline{f}_{0j}L_j)$$
(111)
This concludes our modelisation of the chorodesic atomic time distribution protocol at the post-Newtonian approximation.
## Appendix: Some definitions and notations
Given any space-time with line element:
$$ds^2=g_{\alpha \beta }(x^\rho )dx^\alpha dx^\beta ,\alpha ,\beta ,\mathrm{}=0,1,2,3,x^0=ct,$$
let us consider a time-like congruence $``$, $`u^\alpha `$ being its unit tangent vector field ($`u_\alpha u^\alpha =1`$). The Projector into the plane orthogonal to $`u^\alpha `$ is:
$$\widehat{g}_{\alpha \beta }=g_{\alpha \beta }+u_\alpha u_\beta $$
By definition the Newtonian field is the opposite to the acceleration field:
$$\mathrm{\Lambda }_\alpha =u^\rho _\rho u_\alpha ,\mathrm{\Lambda }_\alpha u^\alpha =0.$$
The Coriolis field, or the rotation rate field, is the skew-symmetric 2-rank tensor orthogonal to $`u^\alpha `$:
$$\mathrm{\Omega }_{\alpha \beta }=\widehat{}_\alpha u_\beta \widehat{}_\beta u_\alpha ,\mathrm{\Omega }_{\alpha \beta }u^\alpha =0.$$
where:
$$\widehat{}_\alpha u_\beta \widehat{g}_\alpha ^\rho \widehat{g}_\beta ^\sigma _\alpha u_\beta .$$
And Born’s deformation rate field is the symmetric 2-rank tensor orthogonal to $`u^\alpha `$:
$$\mathrm{\Sigma }_{\alpha \beta }=\widehat{}_\alpha u_\beta +\widehat{}_\beta u_\alpha ,\mathrm{\Sigma }_{\alpha \beta }u^\alpha =0,$$
Let $`x^\alpha `$ be a system of adapted coordinates, i.e., such that:
$$u^i=0,i,j,\mathrm{}=1,2,3$$
We use the following notations:
$$\xi =\sqrt{g_{00}},\phi _i=\xi ^2g_{0i},$$
and:
$$\widehat{g}_{ij}=g_{ij}+\xi ^2\phi _i\phi _j,\widehat{g}^{ij}=g^{ij}$$
which we shall call the Fermat quo-tensor of the congruence. Here and below quo-tensor refers to an object, well defined on the quotient manifold $`𝒱_3=𝒱_4/`$, whose covariant components are the space components of a tensor of $`𝒱_4`$ orthogonal to $`u^\alpha `$.
The Newtonian field is then the quo-vector:
$$\mathrm{\Lambda }_i=c^2(\widehat{}_i\mathrm{ln}\xi +\frac{1}{c}_t\phi _i),\widehat{}_i_i+\frac{1}{c}\phi _i_t$$
The Coriolis field, or Rotation rate field, is the skew-symmetric quo-tensor:
$$\mathrm{\Omega }_{ij}=c\xi (\widehat{}_i\phi _j\widehat{}_j\phi _i)$$
and the Born’s deformation rate field is the symmetric quo-tensor:
$$\mathrm{\Sigma }_{ij}=\widehat{}_t\widehat{g}_{ij},\widehat{}_t=\xi ^1_t$$
To these familiar geometrical objects it is necessary to add the following ones, , , :
$$\widehat{\mathrm{\Gamma }}_{jk}^i=\frac{1}{2}\widehat{g}^{is}(\stackrel{~}{}_j\widehat{g}_{ks}+\stackrel{~}{}_k\widehat{g}_{js}\stackrel{~}{}_s\widehat{g}_{jk}),$$
which are the Zel’manov-Cattaneo symbols, and:
$$\widehat{R}_{jkl}^i=\stackrel{~}{}_k\widehat{\mathrm{\Gamma }}_{jl}^i\stackrel{~}{}_l\widehat{\mathrm{\Gamma }}_{jk}^i+\widehat{\mathrm{\Gamma }}_{sk}^i\widehat{\mathrm{\Gamma }}_{jl}^s\widehat{\mathrm{\Gamma }}_{sl}^i\widehat{\mathrm{\Gamma }}_{jk}^s$$
which is the Zel’manov-Cattaneo quo-tensor,
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# Exploring the high frequency emission of radio loud X-ray binaries
## 1 Introduction
An interesting group of X-ray binary systems in our Galaxy are known to be powerful and efficient sources of radio waves. The number of radio emitting X-ray binaries (REXRBs) detected so far is about 10 $`\%`$ of the total $``$ 200 systems catalogued (Hjellming & Han 1995). Although far from representing a numerous population, their remarkable properties and scaled down similarity with extragalactic AGNs and quasars makes them to deserve a careful study based on multi-wavelength monitoring programs.
Radio emission in REXRBs is normally highly variable and of non-thermal synchrotron origin. Radio outbursts with different amplitude are frequently detected and interpreted as synchrotron radiation due to the ejection and expansion of ionized plasma clouds (plasmons), usually following a super Eddington accretion event. Recent multi-wavelength monitoring of radio outbursts from the microquasar system GRS 1915+105 (Fender et al. 1997a; Mirabel et al. 1998) have revealed that the flaring synchrotron emission extends well beyond the centimetric domain, reaching up to infrared wavelengths. The energetic implications of this fact are considerable (Mirabel et al. 1998), and it would be important to investigate if similar behaviour is observed in other REXRBs. In an attempt to better assess this issue, we undertook a daily monitoring campaign in the mm domain for some well known objects in the REXRB class. Our main goal here was to study the variability and spectral index properties of selected sources in a wide frequency range. Whenever possible, we have taken advantatge from the availability of daily monitorings in the cm domain, thanks to the Ryle Telescope and to the Green Bank Interferometer. This allowed us to estimate the source spectral indices between radio frequencies separated by two orders of magnitude.
The targets for the observing program were chosen among the brightest REXRBs with luminous massive companions and declination $`\delta >30^{}`$. They include: Cygnus X-3, SS 433, LSI+61303, Cygnus X-1 and GRS 1915+105. A summary of their main physical properties is condensed in Table 1. Previous cm observations for all of them are abundant in the literature, and a few mm detections have been reported as well. Nevertheless, no extended and truly simultaneous cm/mm monitoring has been systematically carried out to our knowledge. The present work is an exploratory step in this direction.
## 2 Observations
The observations of the REXRBs listed in Table 1 were carried out, during the interval 1998 March 14 to March 20, using the following astronomical facilities:
### 2.1 IRAM 30 m-telescope
Observing sessions were conducted at the 30 m telescope of the Institut de Radio Astronomie Millimétrique (IRAM) in Pico Veleta (Spain). The backend installed was the Max Planck Institut für Radioastronomie (MPIfR) 39 channel bolometer array, operating at the 250 GHz frequency (1.25 mm). The technical details of this instrument and those of the 30 m telescope are described by Wild (1995) and Kramer et al. (1998). Each daily session was started with a skydip, pointing and focus sequence. Such sequence was repeated every hour minimum and very often before a new target source. In this way, we were able to ensure an continuous monitoring on the atmospheric opacity, pointing offsets and focusing parameters of the telescope. The method used to determine the pointing corrections was by cross-scanning on a bright calibrator as close as possible to the target source. For the program sources, the ON-OFF technique was preferred in order to obtain higher sensitivity. The ON-OFF procedure should also help to remove possible background extended emission around some of our sources, specially in the case of SS 433 and its associated radio nebula W50. A typical target observation consisted of several tens of symmetric ON-OFFs subscans, during which the antenna beam was nutated on and off the source by 46$`\mathrm{}`$ at a sampling rate of 0.5 s. The duration of each subscan was 10 s. The resulting count rate, for a given source, is taken as the weighted average of all these individual subscans. Source counts were finally converted to flux density using planets as calibrators.
The data reduction was all carried out with the New Imaging Concept (NIC) package, available at the telescope site and originally written at the MPIfR and IRAM. The NIC tasks include spike removal, bad subscan flagging, gain channel and elevation corrections, and correlated channel noise subtraction.
### 2.2 Ryle Telescope
The Ryle Telescope at the Mullard Radio Astronomy Observatory (MRAO) was used to monitor the daily flux variations of four sources in our list: Cygnus X-3, Cygnus X-1, GRS 1915+105 and LSI+61303. The Ryle Telescope operates at the frequency of 15 GHz (2.0 cm). Details of the observing procedure are given in Pooley & Fender (1997). The data reported here were measured with linearly-polarized feed-horns and represent the Stokes parameters $`I+Q`$.
### 2.3 Green Bank Interferometer
The public data from the Green Bank Interferometer<sup>1</sup><sup>1</sup>1 The GBI is a facility of the USA National Science Foundation operated by the NRAO in support of NASA High Energy Astrophysics programs. (GBI) were also retrieved to complete our study. The GBI consists of two 26 m antennas on a 2.4 km baseline observing simultaneously at 2.25 and 8.3 GHz (13 and 3.6 cm). The data taken are made available to the public immediately. All our targets are included among the radio sources routinely monitored by the GBI. Some of them are also observed more than five times daily. Typical errors in GBI data are 4 mJy at 2.25 GHz and 6 mJy at 8.3 GHz for fluxes less than 100 mJy. For fluxes of about 1 Jy (as in SS433), the errors at 2.25 and 8.3 GHz are 15 and 50 mJy respectively. In order to better compare the GBI light curves with those obtained at other instruments, the GBI observations were averaged on a daily basis.
## 3 Results
The summary of our main results, those obtained with the IRAM 30 m-telescope at 250 GHz (1.25 mm), is presented in Table 2. First column lists the star name; the second one indicates the date of the observation and the third column lists the number of subscans performed. Finally, fourth column gives the measured flux density and its error. The error quoted is based on the standard deviation for the estimated flux density. For the cases where the radio star has not been detected, an upper limit of three times the rms noise is given. Among the sources listed in Table 1, the only ones that we could never detect at 250 GHz were Cygnus X-1 and GRS 1915+105. In the cm domain, all targets were found to be detectable with the Ryle and GBI facilities. The cm flux density levels appeared mostly consistent with expectations based on the available literature.
The light curves for Cygnus X-3, SS 433 and LSI+61303 provided several positive mm detections and we additionally plot them in Figs. 1, 3 and 4. These figures also include panels with selected radio spectra from cm to mm wavelengths. They have been computed by interpolating (when possible) the daily Ryle and GBI monitorings at the times of the corresponding IRAM observations. When discussing this spectral information, the spectral index $`\alpha `$ will be defined as $`S_\nu \nu ^\alpha `$, where $`S_\nu `$ is the flux density and $`\nu `$ the frequency.
### 3.1 Cygnus X-3
This system is currently regarded as a Wolf-Rayet star plus a compact object (van Kerkwijk et al. 1992; Fender et al. 1999a). The orbital cycle is assumed to be 4.8 h based on the strong modulation observed with this period, specially in the X-ray domain (Parsignault et al. 1972). This modulation has been also reported at infrared (Becklin et al. 1973) and radio wavelengths (Molnar 1985). Cygnus X-3 is well known for its strong radio flares reaching cm and mm peak flux densities of several Jy (Gregory et al. 1972 and references therein; Nesterov 1992), that may be interpreted in terms of collimated ejection events (e.g. Martí et al. 1992). Waltman et al. (1994) have shown that Cygnus X-3 exhibits periods of normal quiescent emission, varying on time scales of months from 60 to 140 mJy. The average emission level is often around 80 mJy and 90 mJy at 2.25 GHz and 8.3 GHz, respectively. From the GBI data presented in Fig. 1, it is clear that Cygnus X-3 was in a quiescent state during our observations with only minor flaring events taking place.
At mm wavelengths, only five observations have been published in the past. The first two were carried out during a strong flaring period in the radio (Pompherey & Epstein 1972; Baars et al. 1986). Fender et al. (1995) and Tsutsumi et al. (1996) do detect the Cygnus X-3 quiescent radio emission at mm wavelengths using the JCMT. A low level flare peaking at 263 mJy was detected by Altenhoff et al. (1994) with the IRAM 30 m-telescope at 250 GHz. Finally, there is also a sub-mm detection of Cygnus X-3 by Fender et al. (1997b).
The data in this paper, that extend over nearly one week, confirm that Cygnus X-3 is continuously active at 250 GHz even during quiescent emission periods. The peak frequency of the spectrum is most likely to be between 8.3 and 15 GHz most of the time, as for example on 1998 March 16. Although sparsely sampled, our IRAM flux densities are suggestive of a possible mm flare around 1998 March 18.
The occurrence of this minor flaring event is also supported from the hard X-ray (HXR) and soft X-ray (SXR) behaviour of Cygnus X-3 as recorded by the BATSE (20-100 keV) and ASM (1.5-12 keV) instruments on board the CGRO and RXTE satellites, respectively. The HXRs are known to be anticorrelated with quiescent and low level flaring emission in the radio (McCollough et al. 1997). HXRs are also generally anticorrelated with the SXRs. The top panel in Fig. 2 represents the one day average from a number of individual ASM dwells. There are typically 10 dwells of $`90`$ s each per day, that provide an acceptable sampling of the 4.8 h orbital modulation. Therefore, the ASM daily averages may be considered as representative of the unmodulated level of emission. Similarly, the bottom panel in this figure contains the daily averaged HXR data reported by BATSE. We notice here that the HRX flux dropped significantly in coincidence with our bright mm detection, as expected when a low level radio flare occurs.
The 1.25 mm emission on 1998 March 18 was higher than average by a factor of 3. This apparent flaring increase seems to have induced an important evolution in the Cygnus X-3 spectrum. Its main stages are sketched in Fig. 1. We can see here how the high frequency end of the spectrum experiences significative changes. The 8.3-250 GHz spectral index evolves from $`\alpha 0.6`$ during the first three days (March 14-16) to a nearly flat value during the mm ‘maximum’ (March 18), and back to the original negative value after its decay. Unfortunately, the 2 d gap in the Ryle monitoring prevents us from interpolating a reliable 15 GHz measurement at the time of the mm ‘maximum’, that could provide a better view of this event.
At low frequencies, the spectrum remains always optically thick with the 2.25-8.3 GHz spectral index being $`\alpha +0.3`$. The slower evolution and consequent overlapping of flaring events at optically thick frequencies causes the mm flare to be barely evident in the GBI curves.
### 3.2 SS 433
SS 433 is well known to exhibit highly collimated jets flowing out at $`0.26c`$ and precessing every 164 days (e.g. Hjellming & Johnston 1981). The accepted interpretation of the observed sub-arcsecond radio structure assumes that, during the SS 433 radio outbursts, twin plasmons or “bullets”, containing relativistic plasma, are ejected from a central unresolved core into the jet opposite directions (Vermeulen et al. 1987). The source cm emission is consistent with non-thermal synchrotron radiation with optically thin properties. This can be seen clearly in the 2 yr GBI radio flux history (see e.g. Fender et al. 1997c). At mm wavelengths, a steep optically thin spectrum has been also reported (Band & Gordon 1989; Tsutsumi et al. 1996).
At the time of our observations, it is evident from the left panel in Fig. 3 that SS 433 was undergoing one of its outburst events at cm wavelengths. The onset time can be estimated around JD 2450885.5 (1998 March 13), peaking 4 d later in the GBI data. At mm wavelengths, our sampling is unfortunately not very good. Nevertheless, it is likely that the 250 GHz peak took place between JD 2450888-889. The mm flux density on these dates appears to be higher than in the following days, when the outburst decay is already in progress according to GBI.
It is known that SS 433 always exhibits an important optically thin quiescent emission of non-thermal synchrotron origin. According to Seaquist et al. (1982), it can be estimated on average as $`S_\nu ^{\mathrm{quies}}=1.23`$ Jy $`(\nu /`$GHz$`)^{0.6}`$. From the GBI data during the $`20`$ d long period of quiescence, prior to the outburst event, we derive the similar power law $`S_\nu ^{\mathrm{quies}}=1.27\pm 0.03`$ Jy $`(\nu /`$GHz$`)^{0.67\pm 0.02}`$.
In order to better appreciate the spectral evolution of SS 433 during the outburst, we have subtracted this quiescent component from all spectra in the right panel of Fig. 3 and, hereafter, only the flaring emission will be considered. The three representative epochs shown correspond to the outburst decaying part, on JD 2450889, 891 and 892 (March 16, 18 and 19, respectively). For each of these epochs, we plot two estimates of the intrinsic flaring spectrum. They correspond to the use of the Seaquist et al. (1982) and the contemporaneous GBI quiescent spectra, respectively. The difference between both estimates is indicative of the uncertainty involved in the subtraction process. This problem is only important at the IRAM frequency. Probably, the most reliable subtraction is the one involving contemporaneous GBI data.
The results obtained are suggestive that the synchrotron flaring spectrum of SS 433 extends up to the mm regime throughout the full radio outburst. In addition we also find possible evidence of spectral steepening during the outburst decay, at least concerning the GBI frequencies. The spectral index between 2.25 and 8.3 GHz varied from $`\alpha =0.5\pm 0.1`$ (March 16) to $`\alpha =0.9\pm 0.1`$ (March 19). This result does not depend very much on which quiescent spectrum is subtracted. On the other hand, the evolution of the spectral index between 8.3 and 250 GHz also seems to indicate a steepening evolution. Unfortunately, the errors in the subtraction at 250 GHz prevent us from being very certain about it.
### 3.3 LSI+61303
This object is a Be REXRB with periodic and strong outburst events in the radio every 26.5 d. This recurrence is assumed to reflect the system orbital period and was discovered by Taylor & Gregory (1982). LSI+61303 is thus the only selected stars where the onset of radio outbursts can be predicted (see e.g. Paredes et al. 1990).
In the millimetre region, two previous LSI+61303 observations are available in the literature. Altenhoff et al. (1994) observed this star at 250 GHz but they did not detect it. Considering their sensitivity, this can be possibly understood because their observation did not coincide with an outburst episode. On the contrary, Tsutsumi et al. (1996) do detected LSI+61303 at 230 GHz at a flux density level of 10 mJy. This occurred on two consecutive days close to the expected outburst phases, when the source was reported to peak at $`100`$ mJy at cm wavelengths.
In view of these facts, our observing runs were specially scheduled in coincidence with one of LSI+61303 periodic radio outbursts. As it can be seen in the Fig. 4 left panel, the expected flaring event actually took place at cm wavelengths according to the GBI monitoring. Its amplitude and duration were, however, lower than expected. Variations in the outburst shape from cycle to cycle are in fact not unusual, as shown by the extended GBI light curves (see e.g. Ray et al. 1997). In our IRAM measurements, we managed to detect the source at 1.25 mm on two consecutive days that happen to be in close coincidence with the radio outburst peak. The right panel in Fig. 4 shows that, at this time (1998 March 15), the synchrotron spectrum of LSI+61303 extended from cm to mm wavelengths over two frequency decades. A least square fit gives $`S_\nu =(150\pm 3)`$mJy $`(\nu /\mathrm{GHz})^{0.47\pm 0.01}`$ as the best power law representing the data.
During the decay, no mm observation of LSI+61303 resulted in a positive detection and only upper limit estimates are available (see Table 2). All these upper limits are consistent with the extrapolation of the observed 2.25-8.3 GHz spectral indices in the corresponding GBI data ($`\alpha 0.55`$). For example, at the times of lowest GBI emission, the expected 250 GHz flux density should be $`2.4`$ mJy. This value is below our sensitivity.
### 3.4 Cygnus X-1
The cm radio emission of this massive REXRB, and classical black hole candidate, is usually persistent at the $``$ 15 mJy level. Its spectral index is usually flat and some radio outbursts, reaching up to 30-40 mJy, have been occasionally witnessed (Hjellming 1973; Martí et al. 1996). Radio variations at the binary period (5.6 d) and at a longer period (150 d) have been reported (Pooley et al. 1999). The only mm detection so far published is the value of 10$`\pm `$3 mJy given by Altenhoff et al. (1994), suggesting that the generally flat radio spectrum extends into the mm regime. This suspicion has been recently confirmed by Fender et al. (1999b), who found no evidence for a high frequency cut-off up to 220 GHz.
During our monitoring, the centimetric radio emission of Cygnus X-1 behaved in the expected way, i.e., GBI flux densities in the range $`5`$-20 mJy with a flat spectral index being observed. In contrast, very little can be said at 250 GHz. The source remained below our IRAM detection limits throughout all our observations.
### 3.5 GRS 1915+105
This is one of the two superluminal jet sources in the Galaxy as discovered by Mirabel & Rodríguez (1994). The system has been proposed to be a high mass X-ray binary (Mirabel et al. 1997). Its radio emission displays highly active episodes, lasting several months during which multiple and strong ($`1`$ Jy) radio outbursts have been reported. During one of these events, a 234 GHz flux density of 123 mJy was observed with the IRAM 30 m telescope (Rodríguez et al. 1995). As in SS 433, the GRS 1915+105 flaring episodes correspond to the ejection of twin plasmons, at relativistic speeds, sometimes producing apparent superluminal motion on the sky. The most recent sequence of superluminal motion in GRS 1915+105 has been obtained by Fender et al. (1999c) with unprecedented resolution.
We monitored GRS 1915+105 on five different days at 250 GHz. The level of activity indicated by the GBI and Ryle monitoring was relatively weak, i.e., a few tens of mJy. Consistent with this state, no positive 250 GHz detection was achieved in a reliable way.
### 3.6 Other X-ray binaries
On 1998 March 17, most of our program sources were not in good elevation conditions during the scheduled time. Therefore, we devoted part of our run to try to detect three additional X-ray binaries at 250 GHz: the stars LSI+65010, LSI+61235 and X Per.
We observed the first star, LSI+65010, for a total of 37 subscans around 18.5h UT. No positive detection was achieved, with our upper limit estimate being $`<`$17 mJy. Previous attempts to detect this system at cm wavelengths have failed as well. Nelson & Spencer (1988) used the Jodrell Bank Lovell-MkII interferometer at 5 GHz to provide upper limits of $`<`$1.2 mJy and $`<`$3.7 mJy in 1986 and 1987, respectively. Our group has also carried out cm observations of LSI+65010 with the VLA interferometer of NRAO<sup>2</sup><sup>2</sup>2The National Radio Astronomy Observatory is a facility of the USA National Science Foundation operated under cooperative agreement by Associated Universities, Inc. on two consecutive dates, 1993 March 9 and 10. The 20 cm wavelength was used with the array being in its B configuration. A flux density upper limit of $`<`$0.34 mJy was obtained after concatenating all the visibility data.
The second star, LSI+61235, was observed for a total number of 80 subscans around 18.8h UT. Again, no positive detection occurred. The 250 GHz upper limit in this case is $`<`$11 mJy. At the 20 cm wavelength, a $`<0.25`$ mJy upper limit is also available for LSI+61235 from the same VLA runs mentioned above. Finally, X Per appeared to be below our detection limits as well when observed with the IRAM antenna at 18.1h UT. The number of subscans devoted to it was 60 and the resulting upper limit $`<`$15 mJy.
## 4 Discussion
### 4.1 How extended is the synchrotron spectrum in REXRBs?
We have already mentioned that non-thermal radio emission in REXRBs is attributed to the ejection of ionized plasmons containing relativistic electrons. These particles generate synchrotron emission due to the plasmon magnetic field $`B`$ acting on them. In this scenario, the changes in the radio spectrum are actually reflecting those occurring in the electron energy distribution. The energy spectra of the electrons is often well represented by a truncated power law $`N(E)dE=KE^pdE`$ ($`E_{\mathrm{min}}EE_{\mathrm{max}}`$). At any time, the synchrotron emission will extend up to a break frequency $`\nu _{\mathrm{break}}`$ corresponding approximately to the critical frequency of electrons with the highest energy available. In c.g.s. units:
$$\nu _{\mathrm{break}}=6\times 10^{18}BE_{\mathrm{max}}^2.$$
(1)
Based on the data of this paper, the fact that both Cygnus X-3 and SS 433 are persistently detected at 1.25 mm implies for these sources that $`\nu _{\mathrm{break}}>250`$ GHz most of the time. Consequently, relativistic electrons need to be accelerated up to energies of $`E_{\mathrm{max}}>2.0\times 10^4B^{1/2}`$ erg, equivalent to gamma factors of $`\gamma _{\mathrm{max}}>240B^{1/2}`$, to account for the observed high frequency emission. Estimates of the magnetic field in REXRB plasmons come mainly from equipartition arguments, electron age considerations and theoretical fits to light curves during radio outbursts. The resulting values are mostly in the range $`10^3`$-$`10^1`$ G (see e.g. Martí et al. 1992; Martí 1993; Mirabel et al. 1998; Ogley et al. 1998). Adopting $`B10^1`$ G as a representative estimate, the maximum gamma factor of the electrons is typically expected to be higher than $`10^3`$. These order of magnitude considerations are also valid for LSI+61303 and possibly other REXRBs when in outburst.
### 4.2 Interpreting the spectral steepening in Cygnus X-3 and SS 433
In both Cygnus X-3 and SS 433, we have witnessed episodes of spectral steepening during the decaying part of flaring events. As explained below, we interpret this behavior in terms of the energy distribution of the relativistic electrons evolving in time during a flaring event. Of course we cannot rule out that other effects may be work, although the proposed explanation appears to be reproduce this kind of behavior.
Our interpretation is as follows. Since the electrons are undergoing energetic losses of different kind (expansion losses, synchrotron radiation, inverse Compton scattering, etc.), the distribution $`N(E)dE`$ will be shifted towards lower energies. The values of $`E_{\mathrm{max}}`$, and consequently that of $`\nu _{\mathrm{break}}`$ will be decreasing functions of time. A progressive steepening of the radio spectrum is thus expected as the outburst evolves and decays. Theoretical modeling of plasmon evolution, within spherical and jet geometries, has been addressed among others by Paredes et al. (1991), Martí et al. (1992) and Peracaula (1997). When these models are used to compute multi-epoch spectra, we are able to reproduce a clear steepening behavior in strong resemblance to what is observed.
The plots in Fig. 5 are intended to be an illustrative example of what we have described in a qualitative way. They correspond to a modelling attempt of the SS 433 radio outburst previously discussed, i.e., the best sampled event that is well suitable for this purpose. The best fit plasmon physical parameters found are listed in Table 3. The Paredes et al. (1991) formulation has been used in the calculations as improved by Peracaula et al. (1997). Two plasmons ejected into opposite directions at a velocity of $`0.26c`$ are considered. The ejection center coincides with the compact companion of the binary system with an adopted orbital separation of $`1.4\times 10^{15}`$ cm. The assumed luminosity of the optical companion is $`4\times 10^{39}`$ erg s<sup>-1</sup>. It is this radiation field which causes most of the steepening seen in Fig. 5 through strong inverse Compton losses.
## 5 Conclusions
1. We have conducted a series of millimetre and centimetre observations of a sample of REXRBs during a week long interval. Our target list included both Cygnus X-3 and SS 433, that were persistently detected at 250 GHz (1.25 mm) throughout the whole run. The REXRB LSI+61303 was also detected at 250 GHz near the peak of one of its periodic radio outbursts. Several 250 GHz upper limits for other REXRBs are also reported.
2. For the detected sources, our results are in agreement with the synchrotron spectrum in REXRBs extending commonly up to millimetre wavelengths and possibly beyond. This observed fact reinforces the idea that these systems are able to accelerate relativistic electrons to very high energies, at least $`\gamma 10^3`$.
3. The high frequency radio spectrum of Cygnus X-3 and SS 433 was also observed to steepen noticeably during the decay of flaring events. This behavior is interpreted in terms of energetic losses of the synchrotron emitting electrons.
###### Acknowledgements.
JMP and JM acknowledge partial support by DGICYT (PB97-0903). JM is in addition supported by Junta de Andalucía (Spain), and wishes to thank as well the hospitality and support of the Service d’Astrophysique (CEA/Saclay, France) during the early stages of this work.
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# The strange quark mass from flavor breaking in hadronic 𝜏 decays
## I Introduction
The light quark masses, $`m_s`$, $`m_u+m_d`$, are among the least well determined of the fundamental parameters of the Standard Model and, as such, have been the subject of much recent attention, in both the QCD sum rule and lattice communities.
Recent attempts to extract $`m_u+m_d`$ and $`m_s`$ via sum rule analyses of, in the former case, the light quark ($`ud`$) pseudoscalar correlator, and in the latter case, the light-strange ($`us`$) scalar or pseudoscalar correlators, suffer from the problem that the relevant spectral functions are not fully determined experimentally in the region required for the analyses.
Analyses based on vector current correlators involving various pieces of the light quark electromagnetic (EM) current suffer from analogous problems. In the case of Narison’s sum rule based on the difference of the flavor $`33`$ (isovector) and $`88`$ (hypercharge, or isoscalar) correlators , the $`G`$-parity-based identification of the $`33`$ and $`88`$ contributions to the EM hadroproduction cross-section, which would allow the difference of $`33`$ and $`88`$ spectral functions to be determined from experimental data, is valid only in the absence of isospin breaking (IB). The high degree of cancellation (to the level of $`1015\%`$) between the $`33`$ and $`88`$ spectral integrals makes the analysis rather sensitive to the neglect of IB . This sensitivity is compounded by the fact that a sum rule determination of the corrections required to remove the $`38`$ contributions from the experimental data shows that, for reasons which are easily understood , the dominant corrections, associated with the $`\omega `$ contribution to the nominal $`88`$ spectral function, are larger than one would naively expect.<sup>*</sup><sup>*</sup>*The central value $`m_s(1\mathrm{GeV})=176\mathrm{MeV}`$ , obtained neglecting IB corrections, is reduced to $`146\mathrm{MeV}`$ when one applies the IB corrections obtained in the sum rule analysis of Ref. . The necessity of determining the IB corrections theoretically thus prevents one from working with a sum rule whose spectral side is determined solely by experimental data.
A similar problem exists for the sum rule based on the difference of $`33`$ and $`ss`$ vector current correlators , since the portion of the EM hadroproduction cross-section associated with the $`ss`$ part of the EM spectral function is not an experimental observable. In Ref. , it is assumed to be given by the cross-section for the production of the various $`\varphi `$ resonances. This approximation, while no doubt a reasonable one, is exactly valid only if both (1) the Zweig rule is $`100\%`$ satisfied and (2) the $`\varphi `$ resonances are all pure flavor $`\overline{s}s`$ states. The close cancellation (to the $`15\%`$ level) between the $`33`$ and $`ss`$ spectral integrals again makes the analysis sensitive to even small (few $`\%`$) Zweig rule violations (ZRV). To illustrate this sensitivity, let us take the deviation from ideal mixing in the vector meson sector as a measure of the natural scale of ZRV, From Ref. one has that the vector meson mixing angle is either $`36^o`$ or $`39^o`$, depending on whether one uses the linear or quadratic mass formula. and consider a scenario in which ZRV occurs dominantly in the mass matrix and not in the vacuum-to-vector-meson matrix elements of the vector currents. The strange (light) quark part of the EM current then couples only to the strange (light) part of any given resonance. If the flavor content of a given $`\varphi `$ resonance is $`\alpha \overline{s}s+\beta (\overline{u}u+\overline{d}d)/\sqrt{2}`$ (with $`\alpha 1`$ and $`\beta `$ small), the ratio of the square of the full EM $`\varphi `$ decay constant to that of the decay constant describing the coupling only to the $`ss`$ part of the EM current is then $`1\sqrt{2}\beta /\alpha `$. For either the linear or quadratic versions of mixing this ratio is less than $`1`$; including ZRV corrections will thus increase the $`ss`$ spectral function and hence lower the extracted value of $`m_s`$. Taking, to be specific, the case that the radius of the circular part of the FESR contour is $`(1.6\mathrm{GeV})^2`$, we find that, using an identical method of analysis and identical higher dimensional condensate values to those employed in Ref. (and including, for completeness, the small IB isovector contribution to the $`\varphi (1020)`$ EM decay constant determined in Ref. ), the central value of $`m_s(1\mathrm{GeV})`$ obtained ignoring IB and ZRV ($`196\mathrm{MeV}`$) is lowered to $`177\mathrm{MeV}`$ ($`108\mathrm{MeV}`$) for the linear (quadratic) cases, respectively. We stress that the point of this exercise is not to attempt a realistic estimate of ZRV corrections but rather to point out that, given the scale at which such violations are already known to occur, the uncertainties in the extraction of $`m_s`$ associated with the neglect of ZRV are large, and, moreover, cannot be significantly reduced without a major improvement in our theoretical understanding of the precise nature and magnitude of ZRV.In Ref. , the agreement of the $`33`$-$`88`$ and $`33`$-$`ss`$ determinations of $`m_s`$ obtained ignoring IB and ZRV, respectively, was taken as evidence against the size of the IB corrections obtained in Ref. . Note, however, that (1) within errors, the latter result is compatible with either the IB-corrected or uncorrected $`33`$-$`88`$ determination, and (2) two inverse moment sum rule determinations of the $`6^{th}`$ order chiral low-energy constant, $`Q`$, one based on the $`33`$-$`88`$ , and one on the $`\overline{s}u`$-$`33`$ correlator difference , are brought into almost perfect agreement once the IB corrections of Ref. are applied to the former analysis.
In light of the fact that, in each of the analyses above, it is not possible to work with sum rules for which the hadronic spectral function is determined entirely by experimental data, we will, in this paper, instead construct finite energy sum rules (FESR’s) based on the flavor-breaking difference between the sum of the $`ud`$ vector and axial vector correlators and the corresponding sum of $`us`$ correlators, for which, up to $`s=m_\tau ^2`$, the spectral function can be taken from experimental hadronic $`\tau `$ decay data . The rest of the paper is organized as follows. In Section II we provide a brief review, and discuss the practical difficulties to be overcome in arriving at a reliable implementation of this approach. In Section III we describe a construction which leads to FESR’s which successfully overcome these difficulties, and in Section IV we give numerical details and discuss our results.
## II Flavor-Breaking Sum Rules Involving Hadronic $`\tau `$ Decay Data
For a general correlator, $`\mathrm{\Pi }(s)`$, with a cut beginning at $`s=s_{th}`$ and running along the timelike real axis, one obtains from Cauchy’s theorem, defining the spectral function, as usual, by $`\rho \mathrm{Im}\mathrm{\Pi }/\pi `$, the general FESR relation
$$_{s_{th}}^{s_0}𝑑s\rho (s)w(s)=\frac{1}{2\pi i}_{|s|=s_0}𝑑s\mathrm{\Pi }(s)w(s)$$
(1)
where $`w(s)`$ is any function analytic in the region of the contour, $`C`$, consisting of the union of the circle of radius $`s_0`$ in the complex $`s`$-plane and the lines above and below the physical cut, running from $`s_{th}`$ to $`s_0`$.
As is well known, the ratios of $`ud`$ and $`us`$ inclusive hadronic $`\tau `$ decay widths to the $`\tau `$ electronic decay width,
$$R_\tau ^{ij}\frac{\mathrm{\Gamma }[\tau ^{}\nu _\tau \mathrm{hadrons}_{\mathrm{ij}}(\gamma )]}{\mathrm{\Gamma }[\tau ^{}\nu _\tau e^{}\overline{\nu }_e(\gamma )]},$$
(2)
where $`(\gamma )`$ indicates additional photons or lepton pairs, and $`ij=ud,us`$ labels the flavors of the relevant portion of the hadronic weak current, can be expressed as weighted integrals over the relevant spectral functions. Eq (1) then allows these ratios to be recast into a form appropriate for the use of techniques based on the OPE and perturbative QCD . Letting $`J_{ij;V,A}^\mu `$ be the usual vector and axial vector currents with flavor content $`ij`$, and defining the scalar $`J=0,1`$ parts of the corresponding correlators by
$`i{\displaystyle d^4xe^{iqx}}`$ $`0|T\left(J_{ij;V,A}^\mu (x)J_{ij;V,A}^\nu (0)^{}\right)|0`$ (4)
$`\left(g^{\mu \nu }q^2+q^\mu q^\nu \right)\mathrm{\Pi }_{ij;V,A}^{(1)}(q^2)+q^\mu q^\nu \mathrm{\Pi }_{ij;V,A}^{(0)}(q^2),`$
one has
$`R_\tau ^{ij}`$ $`=`$ $`12\pi ^2S_{EW}|V_{ij}|^2{\displaystyle _0^{m_\tau ^2}}{\displaystyle \frac{ds}{m_\tau ^2}}\left(1{\displaystyle \frac{s}{m_\tau ^2}}\right)^2\left[\left(1+2{\displaystyle \frac{s}{m_\tau ^2}}\right)\rho _{ij}^{(1)}(s)+\rho _{ij}^{(0)}(s)\right]`$ (5)
$`=`$ $`6\pi S_{EW}|V_{ij}|^2i{\displaystyle _{|s|=m_\tau ^2}}{\displaystyle \frac{ds}{m_\tau ^2}}\left(1{\displaystyle \frac{s}{m_\tau ^2}}\right)^2\left[\left(1+2{\displaystyle \frac{s}{m_\tau ^2}}\right)\mathrm{\Pi }_{ij}^{(0+1)}(s)2{\displaystyle \frac{s}{m_\tau ^2}}\mathrm{\Pi }_{ij}^{(0)}(s)\right],`$ (6)
where $`\mathrm{\Pi }_{ij}^{(J)}\mathrm{\Pi }_{ij;V}^{(J)}+\mathrm{\Pi }_{ij;A}^{(J)}`$, $`\rho _{ij}^{(J)}(s)`$ are the corresponding spectral functions, $`S_{EW}=1.0194`$ represents the leading electroweak corrections, and $`V_{ij}`$ are the usual CKM matrix elements. Since $`m_\tau ^23\mathrm{GeV}^2`$, the second expression in Eq. (6) is amenable to evaluation using the OPE. Dividing both the hadronic and OPE expressions by $`|V_{ij}|^2`$, and taking the difference of the $`ij=ud`$ and $`us`$ cases, one arrives at a flavor-breaking FESR
$`{\displaystyle _0^1}𝑑y\left(w_{L+T}(y)\mathrm{\Delta }\rho ^{(0+1)}(s)+w_L(y)\mathrm{\Delta }\rho ^{(0)}(s)\right)`$ (7)
$`={\displaystyle \frac{1}{2\pi i}}{\displaystyle _{|y|=1}}𝑑y\left(w_{L+T}(y)\mathrm{\Delta }\mathrm{\Pi }^{(0+1)}(s)+w_L(y)\mathrm{\Delta }\mathrm{\Pi }^{(0)}(s)\right)`$ (8)
where $`ys/m_\tau ^2`$, $`\mathrm{\Delta }\mathrm{\Pi }^{(J)}\mathrm{\Pi }_{ud}^{(J)}\mathrm{\Pi }_{us}^{(J)}`$, $`\mathrm{\Delta }\rho ^{(J)}\rho _{ud}^{(J)}\rho _{us}^{(J)}`$, and $`w_{L+T}`$, $`w_L`$ refer to the longitudinal-plus-transverse ($`(J=0)+(J=1)`$, or “$`L+T`$”) and “longitudinal” ($`(J=0)`$) kinematic weights $`w_{L+T}(y)\left(1y\right)^2\left(1+2y\right)`$ and $`w_L(y)=2y\left(1y\right)^2`$, respectively. The mass-independent ($`D=0`$) piece of the correlator difference $`\mathrm{\Delta }\mathrm{\Pi }^{(J)}`$ on the OPE side of the sum rule Eq. (8) of course vanishes by construction. In the limit that we neglect $`m_{u,d}^2`$ and $`\alpha _sm_{u,d}m_s`$ relative to $`m_s^2`$, moreover, the $`D=2`$ terms in the OPE representation of $`\mathrm{\Pi }_{V+A;ij}^{(J)}`$ become simply proportional to $`m_s^2`$. Were the OPE representations of both the $`L+T`$ and longitudinal contributions above to be well converged at scale $`m_\tau ^2`$, Eq. (8) would thus allow a determination of $`m_s`$ in terms of the difference of experimental non-strange and strange decay number distributions.
The perturbative series for the integrated $`D=2`$ longitudinal contribution in Eq. (8), however, turns out not to be convergent at the scale $`s_0=m_\tau ^2`$ , creating a serious problem for the analysis in the absence of an experimental separation of transverse and longitudinal spectral contributions. This separation is straightforward at low $`s`$ but experimentally problematic above $`1\mathrm{GeV}^2`$.<sup>§</sup><sup>§</sup>§In Ref. , an attempt was made to circumvent this problem by assuming the validity, even in the region of non-convergence, of a relation between the integrated longitudinal OPE vector and axial vector $`D=2`$ contributions valid in the region of convergence of the OPE representations of both. If true, this would allow the longitudinal strange axial integral to be obtained from the longitudinal strange vector integral. The latter can be obtained using the model strange scalar spectral function of Ref. . Using appropriately-weighted FESR’s for the strange pseudoscalar channel, we have now been able to test this assumption, and demonstrate that it is, in fact, incorrect. Our inability to treat the OPE representation of the longitudinal contributions in a reliable manner thus creates difficult-to-quantify uncertainties for any FESR involving significant longitudinal spectral contributions. Existing analyses are included in this category since, for example, the central value for the difference of non-strange and strange spectral integrals from the analysis of Refs. ,
$$\mathrm{\Delta }^{00}\frac{R_\tau ^{ud}}{|V_{ud}|^2}\frac{R_\tau ^{us}}{|V_{us}|^2}=0.394\pm 0.137,$$
(9)
corresponds to $`L+T`$, longitudinal and higher dimension condensate contributions which are $`0.184`$, $`0.155`$ and $`0.055`$, respectively.
Another practical problem is the close cancellation between the rescaled $`us`$ and $`ud`$ spectral integrals for the sum rules above, based on the kinematic weights, $`w_{L+T}`$ and $`w_L`$. In the analysis of Refs. , for example, the cancellation is to the $`10\%`$ level, making the results very sensitive to both small variations in the input parameters and the sizeable experimental errors ($`2030\%`$) on the strange decay number distribution above the $`K^{}`$ region. Two features of the analysis of Refs. illustrate the former sensitivity. First, Refs. employ $`|V_{us}|=0.2218\pm 0.0016`$, c.f. the PDG98 value $`0.2196\pm 0.0023`$. Though compatible within errors, the squares of the two central values differ by $`2\%`$; use of the PDG98 value decreases the flavor-breaking difference, $`\mathrm{\Delta }^{00}`$, by $`17\%`$. Since one cannot reliably employ the OPE representation of the longitudinal contributions, moreover, the longitudinal spectral contribution (which is dominated, at the $`80\%`$ level, by the $`K`$ pole term) must be subtracted; the shift in the inferred $`L+T`$ contribution (used to determine $`m_s`$) is thus even larger ($`36\%`$). Similarly, use of the PDG98 value $`f_K=113.0\pm 1.0\mathrm{MeV}`$ in place of the ALEPH determination, $`f_K=111.5\pm 2.5\mathrm{MeV}`$ lowers the inferred $`L+T`$ contribution to $`\mathrm{\Delta }^{00}`$ by a further $`12\%`$. The combined impact on the central value for $`m_s`$ is thus extremely large, though the two central values are, of course, compatible within the (large) errors quoted in Refs. . The relative size of the residual statistical errors as a fraction of the resulting $`\mathrm{\Delta }^{00}`$ is, of course, also significantly increased by such a decrease in $`\mathrm{\Delta }^{00}`$. It is thus highly desirable to choose, in place of the kinematic weights, weights which produce a less close cancellation between the $`ud`$ and $`us`$ spectral integrals. The easiest way to accomplish this goal is to choose weight functions which fall off more rapidly through the region of the excited strange resonances. This has the happy consequence of also suppressing contributions from the region where both the errors on the strange spectral distribution are large and the transverse/longitudinal separation is experimentally difficult.
The final difficulty to be dealt with is theoretical. Suppose we are able to solve the longitudinal/transverse separation problem, and thus work with FESR’s involving only the $`L+T`$ part of the flavour breaking difference,
$$\mathrm{\Pi }(q^2)\mathrm{\Pi }_{ud,V+A}^{(1+0)}\mathrm{\Pi }_{us,V+A}^{(1+0)}.$$
(10)
The leading ($`D=2`$) $`m_s`$-dependent terms in the OPE representation of $`\mathrm{\Pi }`$ are
$`\left[\mathrm{\Pi }(Q^2)\right]_{D=2}`$ $`=`$ $`{\displaystyle \frac{3}{2\pi }}{\displaystyle \frac{m_s^2(Q^2)}{Q^2}}\left[1+{\displaystyle \frac{7}{3}}a(Q^2)+(19.9332)a(Q^2)^2+\mathrm{}\right]`$ (11)
$``$ $`{\displaystyle \frac{3}{2\pi }}{\displaystyle \frac{m_s^2(Q^2)}{Q^2}}{\displaystyle \underset{k=0}{}}g_ka(Q^2)^k,`$ (12)
with $`a(Q^2)=\alpha _s(Q^2)/\pi `$ and $`m_s(Q^2)`$ the running coupling and running strange quark mass, both at scale $`\mu ^2=Q^2=s`$, in the $`\overline{MS}`$ scheme. The ratio of $`𝒪(a)`$ and $`𝒪(a^2)`$ coefficients in Eq. (12) is rather large ($`8.5`$), signalling potentially slow convergence (with $`\alpha _s(m_\tau ^2)=0.334`$ , the ratio of the $`𝒪(a^2)`$ and $`𝒪(a)`$ terms is $`0.90`$ at $`\mu ^2=m_\tau ^2`$, and $`>1`$ for $`\mu ^2`$ below $`2.2\mathrm{GeV}`$.) In recent analyses , this potential problem is brought under (apparent) control using the method of “contour improvement” . In this method, the logarithms in $`\mathrm{\Pi }`$ are first summed (as has already been done in Eq. (12)) by choosing the renormalization scale equal to $`Q^2`$ at each point on the circle $`|s|=s_0`$. The integrals
$$A_k^{[w_{L+T}]}(s_0)=\frac{1}{2\pi i}_{|s|=s_0}𝑑s\left[\frac{m(Q^2)^2}{Q^2}\right]a(Q^2)^kw_{L+T}(y);y=s/s_0$$
(13)
are then evaluated numerically, using the known $`4`$-loop forms for the running mass and coupling. The OPE side of the $`L+T`$ part of the conventional $`\tau `$ decay sum rule then reduces to a linear combination of the $`A_k^{[w_{L+T}]}(m_\tau ^2)`$, $`k=0,1,2`$, with the index $`k`$ giving the “contour-improved order”. Both the convergence and the residual scale dependence of the resulting truncated series are significantly improved by this procedure . Since, relative to an expansion in terms of $`a(\mu ^2)`$, for some fixed scale $`\mu ^2`$, contour improvement represents a resummation of the perturbative series, it is possible that this improvement is physically meaningful.
Unfortunately, it turns out that the apparent improvement is not a general one, but rather the result of an accidental suppression of the $`k=2`$ integral. To see this, let us, for illustrative purposes, imagine that the unknown coefficients, $`g_k`$, for $`k3`$, in Eq. (12) grow geometrically, i.e., $`g_k=(19.9332)\left[\frac{19.9332}{7/3}\right]^{k2},k3`$.Note that Refs. employ a form of the $`L+T`$ FESR in which the OPE integral has been partially integrated once in order to re-express it in terms of the difference of $`L+T`$ $`ud`$ and $`us`$ Adler functions. The contour-improved series for the Adler function version differs term-by-term from that based on the direct correlator difference. Though the agreement of the sums of the two versions to second order is excellent, the reader should bear in mind that the relative size of the terms of different order is not the same in the two cases. We then evaluate $`A_k^{[w_{L+T}^N]}(s_0)`$ for $`k=0,\mathrm{},10`$ and $`s_0=m_\tau ^2`$, where $`w_{L+T}^N(y)=w_{L+T}(y)[1y]^N`$, $`N=0,1,2`$, are the “spectral weights” employed in the analyses of Refs. . The results of this exercise, rescaled in each case by the corresponding $`k=0`$ value, are displayed in Table I. In columns 2-4 we see the apparently favorable convergence of the $`k=0,1,2`$ terms already discussed. The results of the remaining columns, however, show that the smallness of the $`k=2`$ term is not the result of a favorable resummation (which would lead also to improved convergence for the remainder of the series) but rather a consequence of the fact that $`A_k^{[w_{L+T}^N]}(m_\tau ^2)`$ has a zero as a function of $`k`$ rather close to $`k=2`$. The magnitudes of the $`k3`$ terms are such that truncation of the series at $`k=2`$ would produce a significant theoretical error, one much larger in magnitude than the size of the $`k=2`$ term.One should bear in mind that, were one to work with the Adler function version of the $`L+T`$ FESR, the assumption of geometric growth of the coefficients of the Adler function difference is not the same as the assumption of geometric growth of the coefficients of the correlator difference itself. The potential convergence problem, however, may also be demonstrated to exist in the former case. The contour improved analysis employing FESR’s based on the spectral weights thus has potentially significant theoretical uncertainties.
In light of the problems discussed above for those FESR’s based on the spectral weights, $`w_{L+T}^N`$, our goal in the next section will be to construct alternate weights which lead to FESR’s which bring these problems under control.
## III The construction of alternate weight functions
We begin our search for an alternate choice of weight function by attempting to understand the source of the potential slow convergence of the contour-improved series noted above. The goal will be to find a weight such that, even were the unknown $`g_k`$, $`k3`$, to grow geometrically, as assumed above, the tail of the contour-improved series would be small relative to the known terms, in contrast to the behavior shown in Table I for the series corresponding to the spectral weights, $`w_{L+T}^N`$. If we succeed in doing so, the reliability of the standard approach, in which the truncation error is taken to be given by the size of the last known term (in this case, $`k=2`$), will, of course, be improved regardless of the actual behavior of the unknown $`g_k`$. We will then attempt to simultaneously impose conditions which reduce the impact of the experimental errors.
To study the source of the slow convergence of the contour-improved series, it is useful to consider the behavior of the factor $`f_k(Q^2)m(Q^2)^2a(Q^2)^kg_k`$, appearing in the integrand of $`g_kA_k^{[w]}(s_0)`$, on the contour $`|s|=s_0`$. Let $`w(y)`$, $`y=s/s_0`$, be any analytic function real on the real $`s`$ axis, and $`Q^2=s_0\mathrm{exp}(i\varphi )`$ ($`\varphi =0,\pi `$ thus correspond to timelike and spacelike points, respectively). One then has
$$g_kA_k^{[w]}(s_0)=\frac{1}{\pi }_0^\pi 𝑑\varphi \mathrm{Re}\left[f_k(Q^2)w\left(\mathrm{exp}(i\varphi )\right)\right].$$
(14)
The behavior of $`\mathrm{Re}(f_k)`$ and $`\mathrm{Im}(f_k)`$ as a function of $`\varphi `$, for $`s_0=m_\tau ^2`$ and $`k=0,\mathrm{},10`$, is shown in Figure 1. We observe that both $`\mathrm{Re}(f_k)`$ and $`\mathrm{Im}(f_k)`$ have zeroes on the circle $`|s|=m_\tau ^2`$, and that these zeroes move with the order $`k`$. Moreover, while $`\mathrm{Re}(f_k)`$ (slowly) decreases with increasing $`k`$ for all angles $`\varphi `$, the magnitude of $`\mathrm{Im}(f_k)`$ is sizeable in the region $`\varphi \pi /2`$ even for $`k5`$. This slow convergence in the backwards (spacelike) direction is the origin of the slow convergence of the $`k3`$ tails of the integrated series shown in Table I, since the factor $`(1y)^{N+2}`$ entering the weight $`w_{L+T}^N`$ has maximum modulus at the spacelike point on the contour, and is more and more sharply peaked in the backward direction as $`N`$ increases. In addition, the behavior of $`\mathrm{Re}(f_2)`$ and $`\mathrm{Im}(f_2)`$ happens to be just such that, combined with the changes of sign of the real and imaginary parts of $`w_{L+T}^N`$, there is a very strong cancellation in the integral over $`\varphi `$ (particularly so for the case $`N=0`$). This strong cancellation is the origin of the “accidental” suppression of the magnitude of the $`k=2`$ term. As we have already seen in Table I, it is potentially dangerous to use weights for which the integrals $`A_k^{[w]}(s_0)`$ are small for a particular $`k`$ (or for a small number of values of $`k`$) only due to such cancellations. Higher order contributions can then easily be large again, thereby spoiling the seemingly good convergence of the first few terms of the contour-improved series.
The behavior of the $`\mathrm{Re}(f_k)`$ and $`\mathrm{Im}(f_k)`$ displayed in Figure 1 allows one not only to understand the origin of the potential convergence problem but also to construct alternate sum rules which avoid it. From Figure 1 it is evident that convergence can be improved by avoiding weights which are large in the spacelike direction. The results of Ref. also indicate that, for the FESR framework to be reliable at scales $`m_\tau ^2`$, it is necessary for the weight function to have a zero at $`s=s_0`$ ($`y=1`$).<sup>\**</sup><sup>\**</sup>\**Such a zero suppresses contributions from the OPE representation in the region near the timelike real axis where, at scales $`m_\tau ^2`$ and below, data shows that it breaks down . We have found two approaches useful for implementing these constraints. The first involves the use of polynomials with “shepherd” zeros, i.e., zeros either on, or near, the regions of the contour one wishes to suppress. The second involves the construction of weights, $`w_p`$, with $`\mathrm{Im}(w_p)`$ peaked on the contour at angles $`\varphi \pi /2`$, thereby avoiding large contributions from $`\mathrm{Im}(f_k)`$, $`k>1`$ (see Figure 1). A convenient and effective choice is to take $`\mathrm{Im}(w_p)`$ to have a Gaussian form on the contour. Choosing the width of the Gaussian to be $`10^{}`$ and the center to be $`\varphi =\varphi _p`$, good convergence of the $`k3`$ tail of the integrated series can be obtained for any $`20^{}\varphi _p90^{}`$. Technically, these profiles can be well represented using polynomials of degree $`K20`$
$$w_p(y)=\underset{i=0}{\overset{K}{}}a_iy^i.$$
(15)
The coefficients $`a_i`$ are determined, upon normalizing $`\mathrm{Im}(w_p)`$ such that $`w_p(0)=1`$, by the Fourier integrals
$$a_0=1,a_k=\frac{2}{\pi }_0^\pi 𝑑\varphi \mathrm{Im}\left(w_p(\varphi )\right)\mathrm{sin}(k\varphi ),k=1\mathrm{}K.$$
(16)
To summarize: given the problems discussed above with those FESR’s involving the spectral weights, $`w_{L+T}^N(y)`$, we would like to find, if possible, an alternate weight choice, $`w(y)`$,
(1) such that $`w(y)`$ is strongly suppressed in the region above $`s1\mathrm{GeV}^2`$, in order to (a) reduce the degree of cancellation between the $`ud`$ and $`us`$ spectral integrals, (b) reduce the impact of the large experimental errors in the $`us`$ spectral distribution above the $`K^{}`$ region, and (c) minimize the role of the longitudinal subtraction which must, at present, be performed theoretically; and
(2) such that $`w(y)`$ emphasizes those regions of the contour $`|s|=s_0`$ for which the convergence of the $`D=2`$ series is favorable.
It is, of course, not a priori obvious that there exist $`w(y)`$ having the desired properties. We have, however, succeeded in constructing several polynomial weights which do.<sup>††</sup><sup>††</sup>††An important further restriction results from the observation that, in the FESR framework, higher dimension contributions are suppressed only by inverse powers of $`s_0`$; in order to avoid generating potentially large, and unknown, higher dimension contributions, therefore, the coefficients of the polynomials we construct should all be comparable in magnitude to the leading coefficient, $`a_0=1`$. We have chosen to implement this constraint by keeping all coefficients less than $`2`$ in magnitude. Since, as we will see below, the resulting weights do not contain $`w_{L+T}(y)`$ as a factor, the approach is less inclusive than the analysis employing $`w_{L+T}(y)`$ , but it has the advantage of being theoretically cleaner.
The strategy involving shepherd zeros can be implemented with the zeros either on or off the contour. The first weight we have constructed satisfying the criteria above has all zeros on the contour, and is given by
$$w_{10}(y)=[1y]^4[1+y]^2[1+y^2][1+y+y^2]=1yy^2+2y^5y^8y^9+y^{10}.$$
(17)
The absence of $`𝒪(y^3,y^4)`$ terms, which suppresses $`D=8,10`$ contributions, is an additional positive feature of this weight. The fourth order zero at $`y=1`$ and second order zero at $`y=1`$ provide the desired suppressions of the timelike and spacelike regions. An alternate family of weights still having a fourth order zero at $`y=1`$, but with the remaining zeros moved off the contour and at a distance $`r`$ from the origin, is
$$\widehat{w}(r,\mathrm{cos}\theta _1,\mathrm{cos}\theta _2,y)=[1y]^4\left[1+\frac{y}{r}\right]^2\left[1+2\frac{y}{r}\mathrm{cos}\theta _1+\frac{y^2}{r^2}\right]\left[1+2\frac{y}{r}\mathrm{cos}\theta _2+\frac{y^2}{r^2}\right]$$
(18)
($`\theta _1`$ and $`\theta _2`$ give the angular positions of the pairs of off-contour complex conjugate zeros corresponding to the last two factors, with respect to the spacelike direction). The choice $`(r,\mathrm{cos}\theta _1,\mathrm{cos}\theta _2)=(1.2,0.5,0.1)`$ produces a second solution to the constraints above, one whose biggest coefficient is $`a_1=4/3`$. We denote this solution by
$$\widehat{w}_{10}(y)=\widehat{w}(1.2,0.5,0.1,y).$$
(19)
In the approach based on weights which have imaginary parts with a Gaussian profile on the contour, we choose a basis of such weights having different centers, $`\varphi _p`$. As noted above, so long as all the $`\varphi _p`$ lie in the interval $`20^{}\varphi _p90^{}`$, all of the corresponding integrated $`D=2`$ perturbative series will be under control. We then form linear combinations of these weights having different $`\varphi _p`$ in such a way as to construct a new weight which not only retains this good convergence, but at the same time has a zero of sufficiently high order at $`y=1`$ to strongly suppress contributions to the spectral integral from the region $`y>0.5`$. The weight of this type which most successfully satisfies the criteria discussed above has a rapid high-$`s`$ falloff produced by a $`6^{th}`$ order zero at $`y=1`$, a largest coefficient $`a_4=2.087`$, and is given by
$`w_{20}(y)=(1y)^6[1`$ $`+4.2451y+9.4682y^2+14.4155y^3+16.4589y^4+14.6598y^5`$ (22)
$`+10.2818y^6+5.5567y^7+2.1157y^8+0.3520y^90.2065y^{10}`$
$`0.2154y^{11}0.1040y^{12}0.03040y^{13}0.0045y^{14}].`$
The (vastly) improved convergence of the $`k3`$ tail of the integrated $`D=2`$ series for the weights $`w_{10}`$, $`\widehat{w}_{10}`$ and $`w_{20}`$ is displayed in Table II. The entries, as in Table I, have been rescaled by the corresponding $`k=0`$ value, and hence correspond to the ratios, $`g_kA_k^{[w]}(m_\tau ^2)/A_0^{[w]}`$. The results also show that an estimate of the truncation error given by the magnitude of the $`k=2`$ term is, for the new weights, almost certainly a very conservative one. We will demonstrate, in the next section, that the suppression of the high-$`s`$ region of the spectrum produced by the new weights is also sufficient to significantly reduce the impact of the experimental errors.
## IV Numerical Analysis and Results
In performing the numerical analysis of the FESR’s constructed above, we employ the ALEPH data for the nonstrange and strange number distributions<sup>‡‡</sup><sup>‡‡</sup>‡‡The 1998 tabulation of the nonstrange data receives a small overall normalization correction as a result of the shift in $`R_\tau ^{us}`$ between the preliminary 1998 and final 1999 analyses. We thank Shaomin Chen for bringing this point to our attention. and PDG98 values for $`f_K`$, $`f_\pi `$, $`|V_{ud}|`$ and $`|V_{us}|`$. As noted above, the weights have been chosen in such a way that, although theoretical input is required in order to subtract the longitudinal contributions to the experimental number distributions, and hence obtain the $`L+T`$ spectral functions, the effect of this subtraction on the final value of $`m_s`$ is negligible. We will quantify this statement below. Once the $`L+T`$ spectral function has been determined, it is a straightforward matter to evaluate the weighted $`L+T`$ spectral integrals. The choice of steeply falling weights ensures that the strange spectral integrals are dominated by the $`K`$ and $`K^{}`$ contributions, for which the experimental errors are much smaller than those of the rest of the strange number distribution. This plays a major role in reducing the impact of experimental errors on the final extracted value of $`m_s`$. To get a realistic determination of these errors it is important to separate correlated and uncorrelated errors, and also to take into account the strong correlations between the spectral integrals involving different weights.
The nature of the longitudinal subtraction differs significantly in the low-$`s`$ and high-$`s`$ ($`>1\mathrm{GeV}`$) regions. For low $`s`$, the $`\pi `$ and $`K`$ pole subtractions are experimentally unambiguous. For high $`s`$ (the resonance region), the longitudinal contributions are proportional to $`(m_s\pm m_u)^2`$, $`(m_d\pm m_u)^2`$, for $`us`$, $`ud`$, respectively, and hence dominated by the $`us`$ contributions. The longitudinal $`us`$ vector contribution is inferred from the strange scalar spectral function of Ref. . This procedure is consistent provided the value of $`m_s`$ resulting from the present analysis is compatible with that from the strange scalar channel , which it turns out to be. The longitudinal $`us`$ axial vector contribution is similarly inferred from the spectral function of the strange pseudoscalar channel. The latter is obtained by fixing the excited resonance decay constants of a sum-of-resonances spectral ansatz through matching of the hadronic and OPE sides of a family of “pinch-weighted” FESR’s, in analogy to the analysis of Ref .<sup>\**</sup><sup>\**</sup>\**The corresponding procedure works very well in the isovector vector channel, where the results can be checked against the well-known experimental spectral function . A similar statement is true even in channels with strongly attractive interactions near threshold, for which the spectral function will be poorly represented near threshold by the tail of a Breit-Wigner resonance form with “conventional” $`s`$-dependent width. For example, using the value of $`m_s`$ obtained from the strange scalar channel analysis as input and redoing the strange scalar channel analysis, using now a sum-of-resonances spectral ansatz in place of the more realistic ansatz of Ref. , one finds that the ansatz of Ref. is well-reproduced in the region of the dominant $`K_0^{}(1430)`$ peak. One can also use this approach to check the self-consistency between the assumed longitudinal contributions and the output $`m_s`$ value in kinematic-weight-based analysis of Ref. . It turns out that the high-$`s`$ longitudinal contributions assumed are more than a factor of $`2`$ smaller than would be expected based on the extracted value of $`m_s`$. If one employs the PDG98 values for $`|V_{us}|`$ and $`f_K`$, as discussed above, however, the assumed longitudinal contribution becomes compatible within the errors assigned to it in Ref. . The input value of $`m_s`$ required for this analysis should, in principle, be determined iteratively. We have, however, employed as input the value of $`m_s`$ obtained from the strange scalar analysis of Ref. , $`m_s(1\mathrm{GeV})=159\pm 11\mathrm{MeV}`$. This turns out to be consistent with our final result for $`m_s`$. Moreover, for the steeply-falling weights employed in our analysis, the sum of the high-$`s`$ $`V`$ and $`A`$ longitudinal subtractions is at the $`<0.1\%`$ level of the $`us`$ spectral integral, and hence at the $`<1\%`$ level in the $`ud`$-$`us`$ difference. As such, even were our evaluation to be in error by $`100\%`$, the effect on $`m_s`$ would be completely negligible on the scale of the other errors present in the analysis.
On the OPE side, we retain contributions up to and including $`D=8`$. The leading $`D=2`$ term was given above.
The $`D=4`$ contribution is
$`\left[\mathrm{\Pi }(Q^2)\right]_{(D=4)}`$ $`=`$ $`{\displaystyle \frac{2}{Q^4}}[(m_{\mathrm{}}<\overline{\mathrm{}}\mathrm{}>I_s)(1a(Q^2){\displaystyle \frac{13}{3}}a(Q^2)^2)`$ (24)
$`+{\displaystyle \frac{3}{7\pi ^2}}m_s^4(Q^2)({\displaystyle \frac{1}{a(Q^2)}}{\displaystyle \frac{7}{12}})],`$
where $`I_s`$ is the usual RG invariant modification of the non-normal-order strange quark condensate , $`m_{\mathrm{}}`$ is the average of the light $`u`$, $`d`$ masses, and $`<\overline{\mathrm{}}\mathrm{}>`$ is the light ($`u,d`$) condensate. We use the quark mass ratios determined from the ChPT analyses of Ref. , the GMO relation $`2m_{\mathrm{}}<\overline{\mathrm{}}\mathrm{}>=f_\pi ^2m_\pi ^2`$, and the range of values $`0.7<\overline{s}s/\overline{\mathrm{}}\mathrm{}<1`$ for the ratio of condensates. The contour integrals are performed as described below.
For the $`D=6`$ contribution we employ a rescaled version of the vacuum saturation approximation (VSA). From the results of Ref. , one finds
$$\left[\mathrm{\Pi }(Q^2)\right]_{(D=6)}=\frac{64\pi \rho \alpha _s}{81Q^6}\left[<\overline{\mathrm{}}\mathrm{}>^2<\overline{s}s>^2\right],$$
(25)
where $`\rho `$ represents a multiplicative rescaling of the VSA estimate. The analogous rescaling has been determined empirically for the isovector vector channel and the isospin-breaking vector $`38`$ correlator, and found to be $`5`$ in both cases . For the weights employed in our analysis, it turns out that the integrated $`D=6`$ contributions are very small. We are, therefore, able to employ the very conservative estimate $`\rho =5\pm 5`$ for the degree of VSA violation without significantly affecting the overall theoretical error. The combination $`\rho \alpha _s<\overline{q}q>^2`$ in Eq. (25) is to be understood as an effective RG-invariant combination for the evaluation of the OPE contour integrals.
Finally, for the $`D=8`$ contribution, we assume
$$\left[\mathrm{\Pi }(Q^2)\right]_{(D=8)}=\frac{C_8}{Q^8}.$$
(26)
For $`w_{10}`$ this term does not contribute to the integrated OPE; for $`w_{20}`$ and $`\widehat{w}_{10}`$, the value of the effective RG-invariant condensate combination, $`C_8`$, is to be determined as part of the analysis.
As noted above, the OPE contour integrals (for all $`D`$) are performed using the contour improvement prescription. Four-loop versions of the running mass and coupling are employed. To be specific, we have solved analytically for the running mass and coupling using the 4-loop truncated versions of the $`\beta `$ and $`\gamma `$ functions, with the value determined in nonstrange hadronic $`\tau `$ decays, $`\alpha _s(m_\tau ^2)=0.334\pm 0.022`$ , as input. Following conventional practice, we take the error associated with the truncation of the perturbative series for the Wilson coefficient of the $`D=2`$ term at $`𝒪(a^2)`$ to be equal to the value of the last ($`𝒪(a^2)`$) contribution retained. In light of the discussion above we consider this to represent an extremely conservative estimate.
From the point of view of uncertainties on the OPE side, the $`w_{10}`$ sum rule is favored over the $`\widehat{w}_{10}`$ and $`w_{20}`$ sum rules for three reasons: (1) it has no $`D=8,10`$ contributions, (2) it has the smallest truncation error, and (3) it has the smallest errors associated with uncertainties in the input values of the $`D=4`$ and $`D=6`$ condensates.<sup>\*†</sup><sup>\*†</sup>\*†Combining the errors associated with truncation, the condensate input values, and the uncertainty on $`\alpha _s(m_\tau ^2)`$ in quadrature, the resulting errors on $`m_s`$ are $`7.7\%`$, $`8.2\%`$ and $`8.4\%`$ for $`w_{10}`$, $`\widehat{w}_{10}`$ and $`w_{20}`$, respectively. In Table III we display, as a function of $`s_0`$, the extracted values of $`m_s(1\mathrm{GeV}^2)`$ obtained from the $`w_{10}`$ sum rule, analyzed neglecting contributions of dimension 12 and higher. Central values have been used for all input on the OPE side and for the experimental spectral data. For the analysis to be self-consistent, the extracted value of $`m_s`$ should be independent of $`s_0`$. This will be true for $`s_0`$ sufficiently large that the $`D12`$ contributions are negligible. As $`s_0`$ is decreased, the extracted $`m_s`$ values should eventually deviate from a constant, signalling the growth of the higher dimension terms. From the Table we see that the range $`2.75\mathrm{GeV}^2<s_0<3.15\mathrm{GeV}^2`$ provides an extremely good window of stability. In view of the falloff begining around $`s_02.55\mathrm{GeV}^2`$, we will work in the range $`s_02.55\mathrm{GeV}^2`$ in the discussions which follow. It is worth stressing that the central values obtained from $`w_{20}`$ and $`\widehat{w}_{10}`$ sum rules, though having slightly larger theoretical errors, are nonetheless completely consistent with those above: in the window $`2.55\mathrm{GeV}^2s_03.15\mathrm{GeV}^2`$, one finds that the range of solutions for $`m_s(1\mathrm{GeV}^2)`$ lies between $`156`$ and $`161\mathrm{MeV}`$ for $`w_{20}`$, $`158`$ and $`164\mathrm{MeV}`$ for $`\widehat{w}_{10}`$, and, as we saw already in Table III, $`159`$ and $`163\mathrm{MeV}`$ for $`w_{10}`$. In contrast, the $`w_{L+T}`$ sum rule, for which the longitudinal subtraction is important, and the $`D=2`$ convergence is not well under control, yields a range between $`161`$ and $`184`$ (with, moreover, inconsistent solutions for $`C_8`$).
From the point of view of the impact of the errors present in existing experimental data, the theoretically favored $`w_{10}`$ weight is, unfortunately, no longer the favored one. The reason is that, although the impact of the errors in the high-$`s`$ region of the $`us`$ spectrum has been strongly suppressed by the rapid falloff of the weights employed, the $`ud`$-$`us`$ cancellation is still rather close (e.g., at $`s_0=m_\tau ^2`$, to the level of $`6.0\%`$ for $`w_{10}`$, $`6.8\%`$ for $`\widehat{w}_{10}`$ and $`8.6\%`$ for $`w_{20}`$, to be compared with $`3.7\%`$, $`6.5\%`$ and $`9.3\%`$ for the $`w_{L+T}^N`$, $`N=0,1,2`$.) Although the dominant errors (those from the $`K^{}`$ region of the $`us`$ spectrum) are reasonably small, they are still large enough that the relative size of the residual statistical error grows very rapidly with the increase in the degree of cancellation. Thus, e.g., at $`s_0=m_\tau ^2`$, the statistical error represents $`42\%`$, $`36\%`$, $`26\%`$, $`77\%`$, $`38\%`$ and $`23\%`$ of the $`ud`$-$`us`$ spectral difference for the $`w_{10}`$, $`\widehat{w}_{10}`$, $`w_{20}`$, $`w_{L+T}^0`$, $`w_{L+T}^1`$, and $`w_{L+T}^2`$ sum rules, respectively.<sup>\*‡</sup><sup>\*‡</sup>\*‡Because of the high degree of cancellation, reducing $`s_0`$, which increases the degree of suppression of the (already small) high-$`s`$ $`us`$ contributions, still has a non-trivial effect; e.g., the relative statistical error for the $`w_{20}`$ sum rule is reduced from $`26\%`$ to $`19\%`$ when $`s_0`$ is lowered from $`m_\tau ^2`$ to $`2.55\mathrm{GeV}^2`$. The present experimental situation is, therefore, such that the errors on our final result for $`m_s`$ are minimized by working with $`w_{20}`$, rather than $`w_{10}`$.
Working with the $`w_{20}`$ sum rule in the window specified above we find, for our best fit,
$$m_s(1\mathrm{GeV}^2)=158.6\pm 18.7\pm 16.3\pm 13.3\mathrm{MeV},$$
(27)
which is equivalent to
$$m_s(4\mathrm{GeV}^2)=115.1\pm 13.6\pm 11.8\pm 9.7\mathrm{MeV},$$
(28)
where in both of Eqs. (27) and (28) the first error is statistical, the second is due to the uncertainty on $`|V_{us}|`$, and the third theoretical. The theoretical error has been obtained by combining the following in quadrature (where we quote the numerical values corresponding to Eq. (27) to be specific): $`\pm 5.2\mathrm{MeV}`$, associated with the error on $`\alpha _s(m_\tau ^2)`$; $`\pm 3.6\mathrm{MeV}`$, associated with the uncertainty in $`<\overline{s}s>/<\overline{\mathrm{}}\mathrm{}>`$; $`\pm 1.6\mathrm{MeV}`$, associated with the variation of $`m_s`$ within the window $`2.55\mathrm{GeV}^2s_0m_\tau ^2`$; $`\pm 0.6\mathrm{MeV}`$, associated with the uncertainty in the VSA-violating parameter, $`\rho `$; and $`\pm 11.6\mathrm{MeV}`$, associated with truncation of the $`D=2`$ series. The latter obviously remains the dominant source of theoretical error, despite the significant improvement produced by the use of the new weights. Figure 2 displays the quality of the match between the OPE and spectral integral sides of the $`w_{20}`$ sum rule corresponding to the fit above; the agreement in the previously-established stability window, $`s_0>2.55\mathrm{GeV}^2`$, is obviously excellent. The divergence of the OPE and spectral integral curves below $`s_02.55\mathrm{GeV}^2`$ is precisely what one would expect based on the observation above that, for the $`w_{10}`$ sum rule, $`D>10`$ contributions, not included in the truncated OPE representation, begin to become important in this region.
The result of Eqs. (27) and (28) is in good agreement with the strange scalar channel results of Refs. and , the strange pseudoscalar channel result of Ref. , and the recent hadronic $`\tau `$ decay analysis of Ref. , but, we believe, has signficantly reduced theoretical and experimental errors. In particular, the statistical error has, at this point, been reduced almost to the level of that associated with the uncertainty in $`|V_{us}|`$.
Improvements in the accuracy of the experimental $`us`$ spectral data, in particular in the $`K^{}`$ region, could lead to a significant improvement in the size of the statistical error. Such an improvement should be possible using BaBar data. Reduced uncertainties in our knowledge of $`|V_{us}|`$ would also be helpful. On the theoretical side, while significant improvements in the accuracy of the spectral data would allow one to move from the $`w_{20}`$ to the $`w_{10}`$ sum rule, the decrease in the theoretical uncertainty that would result from this shift would be only $`1.3\mathrm{MeV}`$. Far more likely to lead to a significant improvement in the size of the theoretical error would be a computation of the $`𝒪(a^3)`$ coefficient in the $`D=2`$ contribution to the flavor-breaking correlator difference, $`\mathrm{\Pi }`$.
## ACKNOWLEDGMENTS
The authors would like to thank A. Höcker and S. Chen for providing detailed information on the ALEPH nonstrange and strange spectral distributions, S. Chen for pointing out the normalization correction to the 1998 nonstrange data necessitated by the results of the 1999 strange data analysis, and G. Colangelo for his collaboration at an early stage of this work. KM acknowledges the ongoing support of the Natural Sciences and Engineering Research Council of Canada, and the hospitality of the Special Research Centre for the Subatomic Structure of Matter at the University of Adelaide, where much of this work was performed, and JK the partial support of the Schweizerischer Nationalfonds and the EEC-TMR program, Contract No. CT 98-0169.
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# Confinement, asymptotic freedom and renormalons in type 0 string duals
## 1 Introduction and summary
The attempts to generalize the Maldacena conjecture to four-dimensional non-supersymmetric YM theories have followed different directions. The approach proposed by Witten is to start from the duality between type II strings and a five-dimensional super-YM; compactifying the fifth dimension on a circle with different boundary conditions for bosons and fermions one breaks supersymmetry and is left with an effective four-dimensional non-supersymmetric theory . A conceptually similar possibility is to reduce the supersymmetry to $`N=1`$ or even $`N=0`$ with the addition of mass deformations .
A different strategy is to start directly from a string theory with worldsheet supersymmetry but no spacetime supersymmetry. Polyakov has advocated this approach and has suggested the use of non-critical strings , and Klebanov and Tseytlin have adapted this suggestion to critical type 0B string theory. On the gravity side one considers a configuration of $`N`$ electric D3-brane in the type 0B theory, and the dual field theory is conjectured to be $`SU(N)`$ four-dimensional YM plus six adjoint scalars. The six scalars are the price to pay for working in the critical dimension.
In this paper we follow the latter approach and we study the renormalization group (RG) flow derived from the equations of motion of type 0B gravity. Klebanov and Tseytlin have found that in the UV one recovers asymptotic freedom,
$$\frac{1}{g_{\mathrm{YM}}^2(E)}\mathrm{ln}\frac{E}{\mathrm{\Lambda }}\mathrm{const}.\mathrm{ln}\mathrm{ln}\frac{E}{\mathrm{\Lambda }}.$$
(1)
The numerical value of the beta function coefficients are not under control because of various corrections discussed below, but it is very encouraging to see the correct logarithmic behaviour. Minahan has instead presented confining IR solutions. Further related work appeared in refs. -. It is therefore natural to ask whether asymptotic freedom and confinement are connected by a single RG flow.
In this paper we perform a detailed analytic and numerical investigation of the RG flow, and we find the following results.
* The UV solution found in refs. is just a member of a one-parameter family of solutions. The free parameter $`\kappa `$, which is related to subleading terms in the UV, determines the IR behaviour of the solution. For $`\kappa `$ smaller than a critical value $`\kappa _c`$ (that we can determine analytically) the solution flows in the IR toward a confining solution. The case considered in corresponds instead to $`\kappa >\kappa _c`$, and it was found in that in this case the solution flows in the IR toward a conformal fixed point at infinite coupling.
* We will then show, both numerically and analytically, the existence of RG trajectories that connect asymptotic freedom in the UV with confinement in the IR.
* We will find that in the UV, beside the logarithmic terms shown in eq. (1), type 0B gravity predicts also power corrections,
$$\frac{1}{g_{\mathrm{YM}}^2(E)}\mathrm{ln}\frac{E}{\mathrm{\Lambda }}c_1\mathrm{ln}\mathrm{ln}\frac{E}{\mathrm{\Lambda }}+c_2F_1(E)\frac{\mathrm{\Lambda }^2}{E^2}+c_3F_2(E)\frac{\mathrm{\Lambda }^4}{E^4}+c_4\frac{\mathrm{\Lambda }^4}{E^4}+O(\frac{1}{E^6})$$
(2)
where the $`c_i`$ are constants to be defined below. These power corrections match exactly the non-perturbative contributions coming from renormalons. The term $`F_1(E)/E^2`$ on the field theory side comes from the first UV renormalon, the term $`F_2(E)/E^4`$ from the second UV renormalon, and similarly we get the contributions of all UV renormalons. The term $`c_4/E^4`$ is the leading IR renormalon, and we correctly find no IR renormalon contribution $`1/E^2`$. Thus, we get the right positions of all renormalon singularities in the Borel plane.
* Even the energy dependence of the prefactors, that is $`F_1(E)`$ for the first UV renormalon and a constant for the leading IR renormalon, matches what is known from the field theory side. In particular, we will predict the energy dependence
$$F_1(E)\mathrm{sin}\left(\beta _T\mathrm{ln}\frac{E}{\mathrm{\Lambda }}+\alpha _T\right)$$
(3)
with $`\alpha _T,\beta _T`$ constants. Expanding at first order in $`\beta _T`$ we get a result consistent with what is obtained in QCD summing diagrams with a single bubble chain. The above closed form is a prediction of the dual string theory for the resummation of multiple chain bubble graphs, which would be very interesting to check from the field theory side.
The above results hold at lowest order in $`\alpha ^{}`$ and string loop corrections, and are presented in sects. 3-6. The fact that some features of the UV limit of the YM theory can be obtained from the gravity side is surprising, since the $`SU(N)`$ YM and the gravity descriptions are duals to each other and are based on opposite expansions, in $`g^2N`$ and $`1/\sqrt{g^2N}`$, respectively. It is therefore important to study the corrections to the gravity predictions in the UV limit. We will discuss them in sect. 7. It has been found in ref. that, if the $`\alpha ^{}`$ corrections can be expressed in terms of the Weyl tensor, then the only price to pay in the UV is that the numerical values of the beta function coefficients cannot be predicted, but the logarithmic running of the coupling is unaltered. Furthermore we will find that, under the same conditions, the position of the renormalon singularities in the Borel plane is independent of the $`\alpha ^{}`$ corrections.
## 2 General strategy for computing the RG flow
In this section we discuss the general strategy and the assumptions used to extract the RG flow in the dual theory from the gravity equations of motion of type 0 strings. The low energy effective action of type 0B strings is written in terms of a (closed string) tachyon field $`T`$, the dilaton $`\varphi `$, the metric, and a doubled set of RR fields. The terms involving the tachyon in the effective action have been discussed by Klebanov and Tseytlin , and are
$$d^{10}x\sqrt{G}e^{2\varphi }\left(\frac{1}{2}G^{mn}_mT_nT+V(T)\right)+d^{10}x\sqrt{G}\frac{1}{2}f(T)|F_5|^2.$$
(4)
Because of world-sheet supersymmetry, the tachyon potential $`V(T)`$ has only even powers of $`T`$. We write (to have the same notations of ) $`V(T)=2g(T)`$, and
$$g(T)=\frac{1}{4}m^2T^2\lambda T^4+O(T^6)=\frac{1}{2}T^2+O(T^4).$$
(5)
The tachyon mass is $`m^2=2/\alpha ^{}`$ and we have set $`\alpha ^{}=1`$. The function $`f(T)`$ describes the coupling of the tachyon to the 5-form field strength $`F_5`$, and is given by
$$f(T)=1+T+\frac{1}{2}T^2+O(T^3)$$
(6)
and the coupling to the other RR fields can be found in ref. . Note that, due to the doubling of RR fields, the 5-form field strength is not constrained to be self-dual, so that $`|F_5|^2`$ is non vanishing. D-branes in this theory have been studied in . Another consequence of the doubling of the RR fields is that there are two kinds of D-branes, electrically and magnetically charged. It is then interesting to consider the solution of the gravity equations of motion corresponding to a large number $`N`$ of coinciding electrically charged 3-branes, since from eq. (4) we see that the coupling of the 3-brane to the 5-form RR field strength can stabilize the closed string tachyon. The tachyon from open strings which ends on parallel D-branes is instead eliminated by the GSO projection . In the Einstein frame, the ansatz for the 3-brane metric can be conveniently written in terms of two functions $`\xi `$ and $`\eta `$ of a transverse variable $`\rho `$,
$$ds_E^2=e^{\frac{\xi }{2}5\eta }d\rho ^2+e^{\frac{\xi }{2}}dx_{(4)}^2+e^{\frac{\xi }{2}\eta }d\mathrm{\Omega }_{(5)}^2.$$
(7)
For the 4-form RR field $`C_{(4)}`$ one takes $`C_{0123}(\rho )`$ as the only non-vanishing component, so that $`F_5=dC_{(4)}`$ has only the component $`(F_5)_{0123\rho }`$; similarly one takes $`T=T(\rho ),\varphi =\varphi (\rho )`$. The equation of motion for the RR field can be explicitly integrated and gives an integration constant $`Q`$ which is just the RR charge, and is proportional to $`N`$. The remaining equations of motion then separate into four dynamical equations
$`\ddot{\varphi }+{\displaystyle \frac{1}{2}}g(T)e^{\frac{\varphi +\xi }{2}5\eta }=0,`$ (8)
$`\ddot{\xi }+{\displaystyle \frac{1}{2}}g(T)e^{\frac{\varphi +\xi }{2}5\eta }+{\displaystyle \frac{2Q^2}{f(T)}}e^{2\xi }=0,`$ (9)
$`\ddot{\eta }+{\displaystyle \frac{1}{2}}g(T)e^{\frac{\varphi +\xi }{2}5\eta }+8e^{4\eta }=0,`$ (10)
$`\ddot{T}+{\displaystyle \frac{2Q^2f^{}(T)}{f^2(T)}}e^{2\xi }+2g^{}(T)e^{\frac{\varphi +\xi }{2}5\eta }=0,`$ (11)
and (due to the invariance under reparametrizations of $`\rho `$) a constraint on the initial values, conserved by the dynamical equations,
$$\frac{\dot{\varphi }^2}{2}+\frac{\dot{\xi }^2}{2}5\dot{\eta }^2+\frac{\dot{T}^2}{4}+g(T)e^{\frac{\varphi +\xi }{2}5\eta }+20e^{4\eta }Q^2\frac{e^{2\xi }}{f(T)}=0.$$
(12)
The dot is the derivative with respect to $`\rho `$. The explicit form of $`f(T)`$ plays an important role in the UV regime, where the tachyon is stabilized at its minimum, $`f^{}(T_0)=0`$ .
There are uncertainties due to the fact that neither $`f(T)`$ nor $`g(T)`$ are known in closed form; actually the latter is not even unambiguously defined, since, using the equations of motions, a term $`^2T`$ in the effective action can be traded for a term $`m^2T`$ so that, for instance, a higher derivative term $`(^2T)^4`$ in the effective action can be reabsorbed into a $`T^4`$ term. In the following we will at first set $`g(T)=(1/2)T^2`$ and $`f(T)=1+T+T^2/2`$ (the minimum of $`f(T)`$ is in this case at $`T_0=1`$). We will then discuss how our results depend on these choices.
Our aim is to solve eqs. (8)–(12), and to translate the solutions into a renormalization group flow of the coupling $`g_{\mathrm{YM}}^2`$ of the dual theory. The first issue is how to read $`g_{\mathrm{YM}}^2`$ from the gravity side. We follow the approach of and we read it from the quark-antiquark potential, computed from the Wilson loops as in refs. . This means that we take
$$g_{\mathrm{YM}}^2\mathrm{exp}\{\varphi (\rho )/2\}$$
(13)
as our definition of the coupling. The exact proportionality factor will not concern us here since, as found in refs. and as we will discuss below, the gravity solution is in any case subject to a number of uncertainties, due to $`\alpha ^{}`$ corrections and to the exact form of $`f(T),g(T)`$, that affect the numerical value of the proportionality coefficient in eq. (13).
The definition (13) is apparently in contradiction with what one would obtain from the D-brane effective action, which seems to give $`g_{\mathrm{YM}}^2\mathrm{exp}\{\varphi \}`$ rather than $`\mathrm{exp}\{\varphi /2\}`$. However, as discussed in , in the type 0 theory there exists also a tachyon tadpole on the D-brane, so that the effective action of the D-brane is proportional not simply to $`e^\varphi `$ but rather to $`k(T)e^\varphi `$, with $`k(T)=1+T/4+O(T^2)`$. Since in the UV the tachyon runs toward its minimum at $`T_0=1`$, higher powers in $`k(T)`$ cannot be neglected and the closed form of $`k(T)`$ is needed to draw conclusions. It is clear, however, that this is a point which deserves further investigations.
Eq. (13) provides us with the dependence of the coupling on the transverse coordinate $`\rho `$. To obtain an RG flow, the second step is to connect $`\rho `$ with a physical energy scale. This can be done as follows. From eq. (7) we see that a dilatation of the 4-dimensional coordinates, $`x\lambda x`$, can be reabsorbed into a rescaling of $`\xi `$ such that $`\mathrm{exp}\{\xi /2\}`$ scales like an energy squared (while rescaling $`\eta `$ and $`\rho `$ so to keep the other terms invariant). We will then define the energy scale from $`\mathrm{exp}\{\xi (\rho )/2\}E^2`$, or
$$\mathrm{ln}\frac{E}{\mathrm{\Lambda }}=\frac{1}{4}\xi (\rho ).$$
(14)
$`\mathrm{\Lambda }`$ is a scale that we will determine later. In the UV region the solutions will approach $`AdS_5\times S^5`$ and, writing the metric of $`AdS_5`$ as $`du^2/u^2+u^2dx_{(4)}^2`$, we see that $`\mathrm{exp}(\xi /2)u^2`$ and then the prescription (14) reduces to the by now standard identification $`uE`$. Far from the UV there is a certain arbitrariness in this definition, that we regard as corresponding to the arbitrariness in the choice of the renormalization scheme on the gauge theory side. Of course the details of the beta function are not independent of these choices, and only universal properties, like the existence of zeros and how the zeros are approached are what really matters.
Solving the equations of motion we get $`\varphi (\rho )`$ and $`\xi (\rho )`$. Combining eqs. (13) and (14) then provides the dependence of $`g_{\mathrm{YM}}^2`$ from the energy scale and therefore the RG flow.
## 3 The UV and IR asymptotics
An asymptotic solution valid in the UV limit has been presented in ref. and, as a more systematic expansion, in ref. . Direct inspection of eqs. (8)–(12) reveals however that there is actually a one-parameter family of solutions. It is convenient to introduce a new variable $`y`$ from
$$\rho e^y.$$
(15)
In the limit $`y\mathrm{}`$ the solution is
$`\varphi `$ $`=`$ $`2\mathrm{ln}y+15\mathrm{ln}2+{\displaystyle \frac{1}{y}}(39\mathrm{ln}y+\kappa )+O\left({\displaystyle \frac{\mathrm{ln}^2y}{y^2}}\right),`$ (16)
$`\xi `$ $`=`$ $`y+\mathrm{ln}2+{\displaystyle \frac{1}{y}}+{\displaystyle \frac{1}{2y^2}}(39\mathrm{ln}y+\kappa 104)+O\left({\displaystyle \frac{\mathrm{ln}^2y}{y^3}}\right),`$ (17)
$`\eta `$ $`=`$ $`{\displaystyle \frac{y}{2}}+\mathrm{ln}2+{\displaystyle \frac{1}{y}}+{\displaystyle \frac{1}{2y^2}}(39\mathrm{ln}y+\kappa 38)+O\left({\displaystyle \frac{\mathrm{ln}^2y}{y^3}}\right),`$ (18)
$`T`$ $`=`$ $`1+{\displaystyle \frac{8}{y}}+{\displaystyle \frac{4}{y^2}}(39\mathrm{ln}y+\kappa 20)+O\left({\displaystyle \frac{\mathrm{ln}^2y}{y^3}}\right).`$ (19)
Here and in the following we set $`Q=1`$; we see from the equations that the solution for generic $`Q`$ can be recovered from
$`\xi (y;Q)`$ $`=`$ $`\xi (y;Q=1)+\mathrm{ln}Q,`$
$`\varphi (y;Q)`$ $`=`$ $`\varphi (y;Q=1)\mathrm{ln}Q,`$ (20)
while $`\eta ,T`$ are independent of $`Q`$; $`\kappa `$ is a free parameter, and the solution found in ref. corresponds to $`\kappa =0`$. In spite of the fact that in the UV $`\kappa `$ only appears in terms which look quite subleading, we will find that its value is very important for determining the IR behaviour of the solution.<sup>1</sup><sup>1</sup>1This free parameter comes out because, inserting into the equations of motion an ansatz of the form of eqs. (16)–(19) with generic coefficients, one finds that in eq. (8) the terms of order $`(\mathrm{ln}y)/y^2`$ cancel automatically, and therefore impose no constraint on the coefficients of the solution. Looking only to the terms up to $`O(1/y)`$ one might think that this parameter can be removed by a conformal rescaling, but this is not true anymore when one considers also the terms $`O(1/y^2)`$ in $`\xi ,\eta `$. In the large $`y`$ limit $`\xi y\mathrm{}`$. Eq. (14) then shows that $`y\mathrm{}`$ (or $`\rho 0`$) is the UV region. At $`y=\mathrm{}`$ the metric reduces to $`AdS_5\times S^5`$.
Using the prescription discussed in sect. 2, one can now extract the RG flow in the UV, with the result
$$g_{\mathrm{YM}}^2(E)\frac{1}{\mathrm{ln}(E/\mathrm{\Lambda })\frac{39}{8}\mathrm{ln}\mathrm{ln}(E/\mathrm{\Lambda })}.$$
(21)
Note that the logarithmic terms are independent of $`\kappa `$. Eq. (21) should be compared to the running of the coupling in the proposed dual theory, that is four-dimensional $`SU(N)`$ YM theory with 6 scalars in the adjoint representation, which at the two-loop level is
$$g_{\mathrm{YM}}^2(E)=\frac{8\pi ^2}{b_1\left(\mathrm{ln}(E/\mathrm{\Lambda })+\frac{b_2}{2b_1^2}\mathrm{ln}\mathrm{ln}(E/\mathrm{\Lambda })\right)}$$
(22)
with
$$b_1=\frac{8}{3}N,\frac{b_2}{2b_1^2}=\frac{3}{16}.$$
(23)
We see that eq. (21) qualitatively reproduces the logarithmic running of the coupling in the UV. The precise value of the beta function coefficients cannot be reliably estimated, since the gravity prediction is affected by $`\alpha ^{}`$ corrections (see refs. and sect. 7). Furthermore, with a generic tachyon potential $`g(T)`$, and a generic function $`f(T)`$ with a minimum at $`T=T_0`$, the factor $`39/8`$ in eq. (21) is modified as
$$\frac{39}{8}\frac{7}{8}+\left(\frac{g^{}(T_0)}{g(T_0)}\right)^2,$$
(24)
so it is clear that the values of the numerical coefficients are not under control.
In the opposite limit, $`\rho \mathrm{}`$, it is again possible to find asymptotic solutions. In this case however there is a much larger variety of solutions. In particular Minahan has considered the generic behaviour
$`\varphi `$ $``$ $`\varphi _1\rho +\varphi _0`$
$`\xi `$ $``$ $`\xi _1\rho +\xi _0`$ (25)
$`\eta `$ $``$ $`\eta _1\rho +\eta _0`$
$`T`$ $``$ $`t_1\rho +t_0,`$
At large $`\rho `$ this is a solution of the equations of motion if
$$\xi _1,\eta _1>0,5\eta _1\frac{1}{2}\varphi _1\frac{1}{2}\xi _1>0$$
(26)
and
$$\frac{1}{2}\varphi _1^2+\frac{1}{2}\xi _1^25\eta _1^2+\frac{1}{4}t_1^2=0.$$
(27)
The latter equation follows from the constraint equation (12), while the inequalities (26) ensure that all exponentials in eqs. (8)–(12) are suppressed and therefore that the equations of motions are satisfied.
In the limit $`\rho \mathrm{}`$ we now have $`\xi +\mathrm{}`$ and eq. (14) tells that this is the IR limit. To investigate confinement one can compute the quark-antiquark potential as in ref. . One considers a Wilson loop on the boundary $`\rho =0`$, with edges along the directions $`x,t`$, and looks for the classical string world-sheet $`\rho (x,t)`$ which has the Wilson loop as its boundary. The Nambu-Goto action is
$$S_{NG}=\frac{1}{2\pi \alpha ^{}}𝑑\sigma 𝑑\tau \sqrt{det(G_{MN}_\alpha X^M_\beta X^N)},$$
(28)
where $`G_{MN}`$ is the string frame metric, which in ten dimensions is related to the Einstein frame metric $`G_{MN}^{(E)}`$ by $`G_{MN}=e^{\varphi /2}G_{MN}^{(E)}`$, and $`G_{MN}^{(E)}`$ for the 3-brane solution is given in eq. (7). Inserting $`G_{MN}`$ into eq. (28) one finds that the static potential $`V(L)`$ of a $`q\overline{q}`$ pair separated by a distance $`L`$ along the direction $`x`$ is obtained from
$$V(L)=\frac{1}{2\pi \alpha ^{}}_0^L𝑑x\left[e^{\varphi 5\eta }\left(\frac{\rho }{x}\right)^2+e^{\varphi \xi }\right]^{1/2},$$
(29)
after subtracting the divergent part corresponding to the energy of two separated massive quarks . The above result holds for generic $`Q`$. The $`Q`$ dependence can be extracted using eq. (3) and gives
$$V(L)=\frac{1}{2\pi \alpha ^{}\sqrt{Q}}_0^L𝑑x\left[[e^{\varphi 5\eta }]_{Q=1}\left(\frac{\rho }{x}\right)^2+\frac{1}{Q}[e^{\varphi \xi }]_{Q=1}\right]^{1/2}.$$
(30)
The result depends crucially on the function
$$\stackrel{~}{f}^2(\rho )[e^{\varphi \xi }]_{Q=1}\stackrel{\rho \mathrm{}}{}e^{(\varphi _1\xi _1)\rho +(\varphi _0\xi _0)}.$$
(31)
(We define $`\varphi _0,\xi _0`$ as the coefficient of the solution with $`Q=1`$). If $`\varphi _1<\xi _1`$, $`\stackrel{~}{f}`$ vanishes at $`\rho =\mathrm{}`$, and the classical string configuration will go all the way to the region where $`\rho \mathrm{}`$, where there is no cost in energy in separating the quarks further; for $`\varphi _1=\xi _1`$, the classical string configuration again goes to the minimum of $`\stackrel{~}{f}(\rho )`$ at $`\rho \mathrm{}`$, but now $`\stackrel{~}{f}(\mathrm{})=\mathrm{exp}\{(\varphi _0\xi _0)\}`$ is non-vanishing; the $`q\overline{q}`$ potential at large $`L`$ then becomes
$$V(L)\sigma L,$$
(32)
with
$$\sigma =\frac{1}{2\pi \alpha ^{}Q}e^{(\varphi _0\xi _0)/2}$$
(33)
If instead $`\varphi _1>\xi _1`$, then the function $`\stackrel{~}{f}(\rho )`$ diverges at $`\rho \mathrm{}`$. On the other hand, it also diverges in the UV, where $`(\varphi \xi )y+\mathrm{}`$, see eqs. (16)–(17). This means that it has a minimum at some value $`\rho _{\mathrm{min}}`$, where the classical string configuration goes. Then we have again a linear $`q\overline{q}`$ potential with
$$\sigma =\frac{1}{2\pi \alpha ^{}Q}\stackrel{~}{f}(\rho _{\mathrm{min}};Q=1).$$
(34)
This gives the relation between the ‘QCD’ string tension $`\sigma `$, the type 0 string tension $`1/(2\pi \alpha ^{})`$ and the RR charge $`Q`$. In the case of $`SU(3)`$ YM, the experiment gives $`\sigma (1\mathrm{fm})^2`$ and eqs. (33) or (34) fix the value of $`\alpha ^{}`$ for the dual string theory.
Summarizing, the condition for confinement is
$$\varphi _1\xi _1.$$
(35)
A very different IR solution has been found by Klebanov and Tseytlin . In this case, approaching the IR limit, the Einstein frame metric becomes again asymptotic to $`AdS_5\times S^5`$, while the dilaton diverges and the tachyon goes to zero. Therefore the dual theory flows toward a conformally invariant point with infinite coupling.
It is at this point natural to ask whether the one-parameter family of UV solutions discussed before is smoothly connected to any of these IR solutions. This is the issue that we will address in the next section numerically and in sect. 5 analytically.
## 4 Numerical integration
The numerical study of eqs. (8)–(11) is not as straightforward as one might hope; actually we found that, if we start from the UV with initial conditions that reproduce the asymptotic behaviour given by eqs. (16)–(19), the numerical integration runs almost immediately into divergencies. The reason for this numerical problem will become clear at the end of this section, but we have found that starting instead from the IR with a solution of the type (25) and integrating toward the UV poses no numerical problem.
However, this class of IR solutions from which we are starting the integration is characterized by more free parameters than the solution (16)-(19) to which we would like to connect in the UV. This means that we have to scan the parameter space of the IR solutions in order to find those very particular solutions (if any) that match to the desired UV behaviour. This can be done as follows. We start by studying an IR ($`\rho 1`$) solution of the asymptotic form
$$\varphi (\rho )=\rho ,\xi (\rho )=\rho ,\eta (\rho )=\frac{1}{\sqrt{5}}\rho +\eta _0,T(\rho )=0,$$
(36)
keeping $`\eta _0`$ as the only free parameter. This is a solution of the type (25), where we have chosen $`\varphi _1=\xi _1`$ in order to have a confining solution, see eq. (35), and we have arbitrarily set $`\xi _1=1`$. We also specialize to solutions such that in the IR the tachyon goes to zero.<sup>2</sup><sup>2</sup>2We have also studied the case of a tachyon linearly growing in the IR, i.e. $`t_10`$ in eq. (25), and we have found that in the UV the behaviour of the solution is essentially the same as in the case $`t_1=0`$. However, the case $`t_1=0`$ is in a sense more solid because our ignorance on higher orders in $`T`$ of the functions $`f(T),g(T)`$ becomes irrelevant in the IR. Eq. (27) then fixes $`\eta _1=1/\sqrt{5}`$. The inequalities (26) are satisfied. The invariance under constant shifts in $`\rho `$ allows to fix, e.g., $`\xi _0=0`$. We also set $`\varphi _0=0`$; from eq. (33) we see that a different value of $`\varphi _0`$ reflects itself on a quantitatively different value of the string tension $`\sigma `$, but we expect that this does not give rise to qualitative differences in the solution (as we have indeed checked numerically, setting $`\eta _0=0`$ and using $`\xi _0`$ as a free parameter).
The results of the numerical integration are as follows. For $`\eta _0`$ smaller than a critical value $`\eta _c`$, we find that $`\eta (\rho )\mathrm{}`$ at a finite value of $`\rho =\rho _0`$, which therefore sets a limit to the range of $`\rho `$. The fact that $`\rho `$ cannot range from $`\mathrm{}`$ to $`+\mathrm{}`$ is to be expected. In the usual Dp-brane solutions $`\rho `$ is indeed defined only in a semi-infinite range which, after a constant shift of $`\rho `$, can always be taken to be $`0\rho <\mathrm{}`$, and $`\rho =0`$ is the D-brane horizon. However, in our case we find that $`\xi (\rho _0)`$ is still finite. Since $`\xi `$ is related to the energy scale, $`\xi =4\mathrm{ln}(E/\mathrm{\Lambda })`$, this means that in these solutions there is no region that we can interpret as an UV limit of the dual theory (or, in the gravity language, these solutions have no horizon). Therefore the solutions that start in the IR with $`\eta _0<\eta _c`$ are not connected to any of the UV solutions given by eqs. (16)–(19).
If instead $`\eta _0>\eta _c`$, the situation is reversed and $`\xi \mathrm{}`$ at a finite value $`\rho _0`$ while $`\eta (\rho _0)`$ and $`\varphi (\rho _0)`$ are still finite. This means that at $`\rho _0`$ the energy scale $`E\mathrm{}`$, so now the solution reaches the UV regime. However, $`\varphi (\rho _0)`$ is finite and therefore $`g_{\mathrm{YM}}^2`$ does not go to zero in the UV. Then, again, this IR solution does not match to the desired UV solution.
The situation is however different for $`\eta _0=\eta _c`$. In this case $`\varphi ,\xi `$ and $`\eta `$ all go toward $`\mathrm{}`$ at the same value $`\rho =\rho _0`$, while $`T1`$; therefore, there are chances of matching this solution with the UV solution of eqs. (16)–(19).
It is convenient to introduce the variable $`y`$ in analogy to eq. (15), so that the UV corresponds to $`y\mathrm{}`$. In the analytic computation of sect. (3) the UV limit was at $`\rho =0`$ while now it is at $`\rho =\rho _0`$, and to have the same convention we must make a constant shift in $`\rho `$ by $`\rho _0`$, so we define $`y`$ from
$$\rho \rho _0=e^y.$$
(37)
With repeated runs we have located accurately the value of $`\eta _c`$ and the corresponding $`\rho _0`$, obtaining $`\eta _c=0.4962700(1),\rho _0=0.2076889(1)`$. Locating $`\eta _c`$ with great precision is of course necessary if we want to follow the solution deep into the UV regime. With a precision $`O(10^7)`$ on $`\eta _c`$ we estimate that we can reliably follow the solution in the UV region up to $`y7\mathrm{ln}10=O(10)`$.
The solution is shown in figs. (1)-(4) (solid lines), plotted against $`y`$. The functions $`\xi ,\eta `$ and $`T`$ in the UV are very well fitted by the leading terms of eqs. (17), (18) and (19), suggesting that we have succeeded in matching the confining IR solution to this solution. For the dilaton, however, the situation is more subtle. Fig. (5) shows $`e^{\varphi /2}1/g_{\mathrm{YM}}^2`$ plotted against $`\xi /4`$, i.e. against $`\mathrm{ln}(E/\mathrm{\Lambda })`$. In the IR the coupling is strong and $`g_{\mathrm{YM}}^21/E^2`$, while in the UV regime $`1/g_{\mathrm{YM}}^2`$ scales linearly with $`\mathrm{ln}E`$.
Thus, first of all we see that we indeed succeeded in matching a confining IR solution with an asymptotically free UV regime, and now we want to understand whether this UV solution is related to eq. (16). If one compares with the data shown in fig. (2) one finds that, numerically, the expansion given by eq. (16) misses badly. It is instructive to understand the reason. In the UV region the numerical result is very well fitted by
$$e^{\varphi /2}A+by+O(\mathrm{ln}y),$$
(38)
and the fit gives $`b0.005`$. Thus, the solution in the UV is well reproduced by
$$\varphi 2\mathrm{ln}\left(A+by+O(\mathrm{ln}y)\right).$$
(39)
At asymptotically large values of $`y`$, this is the same as
$$\varphi 2\mathrm{ln}y2\mathrm{ln}b+O(\mathrm{ln}y/y)+O(1/y).$$
(40)
This is just the leading behaviour predicted by eq. (16), which also predicts $`b=2^{15/2}0.0055`$, in excellent agreement with the value from the fit.
However, the expansion of $`\mathrm{ln}(A+by)`$ is valid only if $`byA`$, and here $`A`$ is of order one while $`b0.0055`$ is quite small. So, while eq. (39) reproduces the data very well, its expansion, and therefore eq. (16), is only valid for $`y1/b181`$, which is well beyond the point where we can push the numerical integration<sup>3</sup><sup>3</sup>3This also explains why we failed to integrate the solution starting from the UV asymptotics: we gave initial conditions that would have reproduced eqs. (16)–(19) in a region $`y10`$ where this is not a good approximation to the solution. On the other hand, in the numerical integration it is very difficult to start from much higher values of $`y`$, because one is confronted with exponentially small terms in the equations of motion.. Then, it should be clear that our numerical solution in the UV is indeed nothing but a member of the family of solutions given in eqs. (16)–(19), and that the analytic expansion for $`\varphi `$ given in eq. (16) is only valid for $`y181`$, rather than for $`y1`$. However, to clear up any doubt, in the next section we will present an improved analytic solution of the equations that is valid for $`y1`$, rather than only for $`y181`$, that in the region $`y181`$ reproduces eq. (16), including all the subleading terms written there, and that in the region $`y<10`$ where the numerical integration is possible reproduces very well the data. An unexpected bonus from this analysis will be that we will also find terms $`\mathrm{exp}\{\mathrm{const}.y\}`$, that we will relate in sect. (6) to renormalon singularities.
The physics of the solution is illustrated by figs. (5)-(7). Fig. (5), as already discussed, shows $`\mathrm{exp}\{\varphi /2\}`$, which is proportional to the inverse coupling $`1/g_{\mathrm{YM}}^2`$, as a function of $`\mathrm{ln}(E/\mathrm{\Lambda })`$. This plot can also be used to define the constant $`\mathrm{\Lambda }`$ which fixes the energy scale. We define it as the value of $`E`$ when asymptotic freedom sets in, so by definition in fig. (5) the change of regime takes place at $`\mathrm{ln}(E/\mathrm{\Lambda })=0`$. On the right we have asymptotic freedom, since $`1/g_{\mathrm{YM}}^2\mathrm{ln}(E/\mathrm{\Lambda })`$, while on the left side of the plot we have strong coupling and confinement.
Fig. (6) shows the beta functions $`\beta (g_{\mathrm{YM}})`$ for our solution. It has no zero except for the perturbative one at $`g_{\mathrm{YM}}=0`$, in contrast with the two-loop result, which has a second zero at $`1/g_{\mathrm{YM}}^2=N/(4\pi )^2`$, of course well outside the limit of validity of the perturbative computation, since there the ’t Hooft coupling is $`Ng_{\mathrm{YM}}^2=(4\pi )^21`$. The numerical result in the intermediate region matches well the analytic results at weak and strong coupling<sup>4</sup><sup>4</sup>4With the vertical scale used in the figure, needed to show the matching of the numerical result with the IR analytic behaviour, the UV behaviour $`\beta (g)g^3`$ cannot be distinguished from the horizontal axis, but the numerical result indeed matches with $`g^3`$ in the UV. Note also that, to produce this plot, we have set to unit the proportionality constant in eq. (13). Once one has the correct proportionality constant, the correct figure is obtained with a rescaling of the units of the axes..
Another interesting quantity is the radius of the 5-sphere (in the Einstein frame, since in this frame the UV metric approaches $`AdS_5\times S^5`$) $`R_{(5)}`$, which is given by eq. (7),
$$R_{(5)}^2=e^{\frac{\xi }{2}\eta }.$$
(41)
If the adjoint scalar fields become massive, the radius of the 5-sphere shrinks to zero. We see from fig. (7) that $`R_{(5)}^2`$ diverges in the IR limit. At the transition between the IR and UV region it bounces and finally settles to the constant value $`1/\sqrt{2}`$ predicted by eqs. (17) and (18). We see that it never vanishes, and therefore the 6 adjoint scalars remain massless.
## 5 Analytic solution
In this section we find an analytic expression for the RG flow which reproduces very well the solution from the IR to the UV regions, and even allows to compute exponentially small corrections in the UV limit.
The results of the previous section suggest the following improved UV ansatz
$`e^{\varphi /2}`$ $`=`$ $`A+by+u(y)`$ (42)
$`\xi (y)`$ $`=`$ $`y+\mathrm{ln}2+\zeta (y)`$ (43)
$`\eta (y)`$ $`=`$ $`{\displaystyle \frac{y}{2}}+\mathrm{ln}2+\omega (y)`$ (44)
$`T(y)`$ $`=`$ $`1+t(y).`$ (45)
where $`u,\zeta ,\omega `$ and $`t`$ are small, in a sense that will be clarified below. Here $`b=2^{15/2}1/181`$, and $`A`$ is a constant of order one, yet to be fixed ($`A>0`$ otherwise the solution would not make sense when $`by<A`$). At first, we linearize the equations of motion with respect to $`u,\zeta ,\omega `$ and $`t`$ (we will discuss later the non-linear terms), and we get
$`\zeta ^{\prime \prime }+\zeta ^{}2\zeta `$ $``$ $`{\displaystyle \frac{2b}{A+by}}\left[1+{\displaystyle \frac{\zeta }{2}}5\omega 2t{\displaystyle \frac{u}{A+by}}\right],`$ (46)
$`\omega ^{\prime \prime }+\omega ^{}2\omega `$ $``$ $`{\displaystyle \frac{2b}{A+by}}\left[1+{\displaystyle \frac{\zeta }{2}}5\omega 2t{\displaystyle \frac{u}{A+by}}\right],`$ (47)
$`t^{\prime \prime }+t^{}+2t`$ $``$ $`+{\displaystyle \frac{16b}{A+by}}\left[1+{\displaystyle \frac{\zeta }{2}}5\omega t{\displaystyle \frac{u}{A+by}}\right],`$ (48)
$$u^{\prime \prime }+\left(1\frac{2b}{A+by}\right)u^{}+\frac{b^2}{(A+by)^2}u\frac{b^2}{A+by}+b(\frac{\zeta }{2}5\omega 2t).$$
(49)
where the prime is $`d/dy`$. We can now take advantage of the fact that $`b`$ is numerically small, and we treat it as a small expansion parameter. Instead $`A+by`$ in the UV is at least of order one. Then, to lowest order in $`b`$ a particular solution of eqs. (46)–(48) is given simply by
$$\zeta (y)=\omega (y)=\frac{1}{8}t(y)=\frac{b}{A+by}\left[1+O\left(\frac{b}{A+by}\right)\right].$$
(50)
In fact, in the region where $`A+by=O(1)`$, $`\zeta ^{}`$ is smaller by a factor $`b`$ compared to $`\zeta `$ while for $`y>>1/b`$ we have $`1\zeta 1/y\zeta ^{}`$, and similar for $`\omega ,t`$. Substituting eq. (50) into eq. (49) we obtain the particular solution
$$u(y)=\frac{39}{2}b\mathrm{ln}(A+by)\left[1+O\left(\frac{b}{A+by}\right)\right],$$
(51)
To obtain the most general solution, we first tentatively assume that the terms $`\zeta ,\omega ,t`$ and $`u/(A+by)`$ on the right-hand side of eqs. (46)–(49) can be neglected compared to 1. Then the most general solution is obtained adding to the above particular solution the solution of the homogeneous equations
$`\zeta ^{\prime \prime }+\zeta ^{}2\zeta `$ $``$ $`0,\omega ^{\prime \prime }+\omega ^{}2\omega 0,`$
$`t^{\prime \prime }+t^{}+2t`$ $``$ $`0,u^{\prime \prime }+u^{}=0.`$ (52)
This would give
$`\zeta (y)`$ $`\stackrel{\mathrm{?}}{=}`$ $`{\displaystyle \frac{b}{A+by}}\left[1+O\left({\displaystyle \frac{b}{A+by}}\right)\right]+B_\xi e^{2y},`$ (53)
$`\omega (y)`$ $`\stackrel{\mathrm{?}}{=}`$ $`{\displaystyle \frac{b}{A+by}}\left[1+O\left({\displaystyle \frac{b}{A+by}}\right)\right]+B_\eta e^{2y},`$ (54)
$`t(y)`$ $`\stackrel{\mathrm{?}}{=}`$ $`{\displaystyle \frac{8b}{A+by}}\left[1+O\left({\displaystyle \frac{b}{A+by}}\right)\right]+B_Te^{y/2}\mathrm{sin}\left({\displaystyle \frac{\sqrt{7}}{2}}y+\alpha _T\right),`$ (55)
$`u(y)`$ $`\stackrel{\mathrm{?}}{=}`$ $`{\displaystyle \frac{39}{2}}b\mathrm{ln}(A+by)\left[1+O\left({\displaystyle \frac{b}{A+by}}\right)\right]+B_ue^y,`$ (56)
where $`B_\xi ,B_\eta ,B_T,B_u`$ and $`\alpha _T`$ are arbitrary coefficients of the homogeneous solutions. Note that the homogeneous equation for the tachyon has two complex solutions $`t=\mathrm{exp}\{(1\pm i\sqrt{7})y/2\}`$, and then the solution is an oscillating function, as indeed we expected from fig. (1).<sup>5</sup><sup>5</sup>5We discarded exponentially growing solutions of eq. (52), as well as the constant solution of $`u^{\prime \prime }+u^{}=0`$, since the constant part of $`u(y)`$ has already been taken care of by the constant $`A`$, eq. (42).
However, if we now insert back $`\zeta ,\omega ,t`$ and $`u/(A+by)`$ on the right-hand side of eqs. (46)–(49) to check whether the approach is self-consistent, we see that the tachyon gives a contribution $`bB_Te^{y/2}\mathrm{sin}((\sqrt{7}/2)y+\alpha _T)`$ to the right-hand side of eqs. (46)–(49), so that, for instance, the term $`e^{2y}`$ in the solution for $`\zeta `$ dominates over this contribution only if $`be^{y/2}e^{2y}`$ or $`y\mathrm{ln}(1/b)`$. In the opposite limit $`y\mathrm{ln}(1/b)`$ (which with our value of $`b`$ means $`y3`$) before considering terms $`B_\xi e^{2y}`$ we must also include the contribution of the exponentially small term coming from $`t`$ to the right-hand side of eqs. (46)–(49), and we must also go beyond the linear approximation: e.g., if we want to get correctly all the terms up to $`e^{2y}`$ we must include the contributions from terms $`t^2,t^3,t^4`$ and $`u^2`$. The correct result then turns out to be of the form
$`e^{\varphi (y)/2}`$ $`=`$ $`A+by{\displaystyle \frac{39}{2}}b\mathrm{ln}(A+by)\left[1+O\left({\displaystyle \frac{b}{A+by}}\right)\right]+`$
$`+`$ $`bB_Te^{y/2}\mathrm{sin}\left({\displaystyle \frac{\sqrt{7}}{2}}y+\alpha _T\right)\left[1+O\left(bg_1(y)\right)\right]+`$
$`+`$ $`\left[B_u+bB_T^2g_2(y)\right]e^y+O\left(e^{3y/2}\right),`$
with $`g_1(y),g_2(y)`$ non-singular and oscillating functions of $`y`$. In this work, we are not particularly interested in their explicit form (for the comparison with the numerical data these terms play a very marginal role), so we computed explicitly only the term $`e^{y/2}`$, at lowest order in $`b`$, but it is important conceptually to understand the structure of the exponentially small corrections, as we will see in sect. (6).
A similar structure can be obtained for $`\xi ,\eta ,T`$,
$`\xi (y)`$ $`=`$ $`y+\mathrm{ln}2+{\displaystyle \frac{b}{A+by}}\left[1+O\left({\displaystyle \frac{b}{A+by}}\right)\right]`$
$`{\displaystyle \frac{bB_T}{A+by}}e^{y/2}\mathrm{sin}\left({\displaystyle \frac{\sqrt{7}}{2}}y+\alpha _T\right)\left[1+O\left({\displaystyle \frac{b}{A+by}}\right)\right]+O(e^y),`$
$`\eta (y)`$ $`=`$ $`{\displaystyle \frac{y}{2}}+\mathrm{ln}2+{\displaystyle \frac{b}{A+by}}\left[1+O\left({\displaystyle \frac{b}{A+by}}\right)\right]`$
$`{\displaystyle \frac{bB_T}{A+by}}e^{y/2}\mathrm{sin}\left({\displaystyle \frac{\sqrt{7}}{2}}y+\alpha _T\right)\left[1+O\left({\displaystyle \frac{b}{A+by}}\right)\right]+O(e^y),`$
$`T(y)`$ $`=`$ $`1+{\displaystyle \frac{8b}{A+by}}\left[1+O\left({\displaystyle \frac{b}{A+by}}\right)\right]+`$
$`+B_Te^{y/2}\mathrm{sin}\left({\displaystyle \frac{\sqrt{7}}{2}}y+\alpha _T\right)\left[1+O\left({\displaystyle \frac{b}{A+by}}\right)\right]+O\left(e^y\right).`$
We can now compare with the numerical solution and extract $`A,B_T,B_u,\alpha _T`$ from a fit. We find $`A1.1,B_T0.8,B_u0.6,\alpha _T0.6`$. This analytic solution is shown in figs. (1-4) (black dashed line) and we see that it reproduces very well the numerical solution already for $`y>1`$.
Expanding eqs. (5)–(5) in the UV we find that they reproduce eqs. (16)–(19), including all subleading terms. The coefficient $`\kappa `$ is related to $`A`$ by
$$\kappa =39\mathrm{ln}b\frac{2A}{b}=\frac{585}{2}\mathrm{ln}2256\sqrt{2}A,$$
(61)
As we noted above, our solution makes sense only for $`A>0`$, since otherwise $`\mathrm{ln}(A+by)`$ becomes imaginary at intermediate values of $`y`$. Eq. (61) then shows that there exists a critical value of $`\kappa `$,
$$\kappa _c=\frac{585}{2}\mathrm{ln}2$$
(62)
and that our solution (5)-(5) can produce, in the deep UV region $`y1/b`$, only solutions of the form (16)-(19) with $`\kappa <\kappa _c`$. The solution considered by Klebanov and Tseytlin has instead $`\kappa =0`$ and it does not belong to this class. Indeed, ref. finds that in the IR it connects to another conformal fixed point at infinite coupling, rather than to a confining solution.
We see therefore that $`\kappa `$ plays a crucial role in determining the IR behaviour of the RG flow. In the language of Wilson’s renormalization group, this means that $`\kappa `$ is a parameter that determines to which universality class the UV theory belongs, and we have found that there are (at least) two universality classes: for $`\kappa <\kappa _c`$ we are in the domain of attraction of a confining IR fixed point, while for $`\kappa >\kappa _c`$ we are in the domain of attraction of a different fixed point.
It is also important to understand how the result changes if we change the form of the function $`f(T)`$. We find that, as long as $`f(T)`$ has a minimum at a finite value $`T_0`$, the qualitative behaviour of the solution does not change. However, it is especially important to known whether the exponential terms $`e^{y/2}`$ are affected by the form of $`f(T)`$. Linearizing eq. (11) for $`f(T)`$ generic, and writing $`T=T_0+t`$, the left-hand side of eq. (48) becomes $`t^{\prime \prime }+t^{}+ct`$, with $`c=f^{\prime \prime }(T_0)/(2f^2(T_0))`$; for $`f(T)=1+T+T^2/2`$, $`c=2`$. The homogeneous equation now has the solutions $`t=e^{\gamma y}`$ with
$$\gamma =\frac{1}{2}(1\pm i\sqrt{4c1}).$$
(63)
We see that, as long as $`c1/4`$, Re $`\gamma `$ is unchanged and we still have solutions of the form $`e^{y/2}\mathrm{sin}(\nu y)`$, but the frequency $`\nu `$ depends on the explicit form of $`f(T)`$. Therefore the value that we found, $`\nu =\sqrt{7}/2`$, does not have a special significance, being related to the unknown form of $`f(T)`$, while the factor $`e^{y/2}`$ is universal as long as $`f(T)`$ obeys
$$2f^{\prime \prime }(T_0)f^2(T_0).$$
(64)
The function $`g(T)`$ instead does not enter the homogeneous equation, so from this point of view its form is irrelevant.
Finally, it is straightforward to compute analytically the corrections to the solution in the IR regime. The corrections are exponentially small in $`\rho `$,
$`\varphi =`$ $`\varphi _1\rho +\varphi _0`$ $`{\displaystyle \underset{i=1}{}}C_i^{(\varphi )}\mathrm{exp}\{\gamma _i^{(\varphi )}\rho \},`$
$`\xi =`$ $`\xi _1\rho +\xi _0`$ $`{\displaystyle \underset{i=1}{}}C_i^{(\xi )}\mathrm{exp}\{\gamma _i^{(\xi )}\rho \},`$ (65)
$`\eta =`$ $`\eta _1\rho +\eta _0`$ $`{\displaystyle \underset{i=1}{}}C_i^{(\eta )}\mathrm{exp}\{\gamma _i^{(\eta )}\rho \},`$
$`T=`$ $`{\displaystyle \underset{i=1}{}}C_i^{(T)}\mathrm{exp}\{\gamma _i^{(T)}\rho \}.`$
We have computed the coefficients $`C_i,\gamma _i`$ for $`i=1,2,3`$. We get, for $`\varphi _1=\xi _1=\sqrt{5}\eta _1`$,
$`\gamma _1^{(\varphi )}=(\sqrt{5}+3)\varphi _1\gamma _2^{(\varphi )}`$ $`=`$ $`2(\sqrt{5}+1)\varphi _1\gamma _3^{(\varphi )}={\displaystyle \frac{9+3\sqrt{5}}{\sqrt{5}}}\varphi _1`$
$`\gamma _1^{(\xi )}=2\varphi _1\gamma _2^{(\xi )}`$ $`=`$ $`4\varphi _1\gamma _3^{(\xi )}=(\sqrt{5}+3)\varphi _1`$
$`\gamma _1^{(\eta )}={\displaystyle \frac{4}{\sqrt{5}}}\varphi _1\gamma _2^{(\eta )}`$ $`=`$ $`{\displaystyle \frac{8}{\sqrt{5}}}\varphi _1\gamma _3^{(\eta )}=(\sqrt{5}+3)\varphi _1`$
$`\gamma _1^{(T)}=2\varphi _1\gamma _2^{(T)}`$ $`=`$ $`(\sqrt{5}+1)\varphi _1\gamma _3^{(T)}=4\varphi _1`$
$`C_1^{(\varphi )}`$ $`=`$ $`{\displaystyle \frac{1}{16(\sqrt{5}+3)^2\varphi _1^6}}e^{\frac{1}{2}\varphi _0\frac{7}{2}\xi _05\eta _0}C_2^{(\varphi )}={\displaystyle \frac{1}{16(\sqrt{5}+1)^4\varphi _1^8}}e^{\varphi _03\xi _010\eta _0}`$
$`C_3^{(\varphi )}`$ $`=`$ $`{\displaystyle \frac{125}{32(9+3\sqrt{5})^2\varphi _1^8}}e^{\frac{1}{2}\varphi _0\frac{7}{2}\xi _09\eta _0}C_1^{(\xi )}={\displaystyle \frac{1}{2\varphi _1^2}}e^{2\xi _0}`$
$`C_2^{(\xi )}`$ $`=`$ $`{\displaystyle \frac{3}{16\varphi _1^4}}e^{4\xi _0}C_3^{(\xi )}={\displaystyle \frac{1}{(\sqrt{5}+3)^2\varphi _1^6}}\left({\displaystyle \frac{2}{(\sqrt{5}+1)^2}}{\displaystyle \frac{1}{16}}\right)e^{\frac{1}{2}\varphi _0\frac{7}{2}\xi _05\eta _0}`$
$`C_1^{(\eta )}`$ $`=`$ $`{\displaystyle \frac{5}{2\varphi _1^2}}e^{4\eta _0}C_2^{(\eta )}={\displaystyle \frac{25}{4\varphi _1^4}}e^{8\eta _0}`$
$`C_3^{(\eta )}`$ $`=`$ $`{\displaystyle \frac{1}{16(\sqrt{5}+3)^2\varphi _1^6}}e^{\frac{1}{2}\varphi _0\frac{7}{2}\xi _05\eta _0}C_1^{(T)}={\displaystyle \frac{1}{2\varphi _1^2}}e^{2\xi _0}`$
$`C_2^{(T)}`$ $`=`$ $`{\displaystyle \frac{1}{(\sqrt{5}+1)^2\varphi _1^4}}e^{\frac{1}{2}\varphi _0\frac{3}{2}\xi _05\eta _0}C_3^{(T)}={\displaystyle \frac{3}{16\varphi _1^4}}e^{4\xi _0}`$
This analytic solution is shown in figs. (1-4) as a blue dotted line, and we see that it reproduces very well the numerical solution in the IR region. Having included the the first three subleading exponential, we see that the IR analytic solution has an overlap with the analytic UV solution, so that the two expansions together describe analytically the solution in the whole region.
## 6 The renormalon singularities
A striking feature of the UV solution is the appearance of the exponentially small terms $`e^{y/2},e^y`$, etc.; in the UV, at leading order, $`y\xi =4\mathrm{ln}E`$, so (extracting an overall proportionality factor $`4b`$) eq. (5) gives
$$\frac{1}{g_{\mathrm{YM}}^2(E)}\mathrm{ln}\frac{E}{\mathrm{\Lambda }}\frac{39}{8}\mathrm{ln}\mathrm{ln}\frac{E}{\mathrm{\Lambda }}+\frac{B_T}{4}F_1(E)\frac{\mathrm{\Lambda }^2}{E^2}+B_T^2F_2(E)\frac{\mathrm{\Lambda }^4}{E^4}+C\frac{\mathrm{\Lambda }^4}{E^4}+O(\frac{1}{E^6})$$
(66)
with
$$F_1(E)=\mathrm{sin}\left(\beta _T\mathrm{ln}\frac{E}{\mathrm{\Lambda }}+\alpha _T\right)\left[1+O(b)\right],$$
(67)
and $`C=B_u/(4b)`$ (however, we expect that the numerical values of these coefficients will be affected by our uncertainties, as it happens to the factor $`39/8`$). $`F_2(E)`$ can be determined adding the terms $`t^2`$ to eqs. (46) and (49). The positive constant $`\beta _T`$ is given by
$$\beta _T=2\left[\frac{2f^{\prime \prime }(T_0)}{f^2(T_0)}1\right]^{1/2},$$
(68)
and depends on the explicit form of $`f(T)`$, so its numerical value is not under control. The exponential correction to the relation $`y\xi `$ gives an $`O(b)`$ contribution to (67).
First of all, we read from eq. (66) that there are power corrections, and they are exactly of the form that is expected from the effect of renormalons.
Let us recall that, given a divergent series
$$G(g_{\mathrm{YM}}^2)=\underset{n=1}{\overset{\mathrm{}}{}}G_ng_{\mathrm{YM}}^{2n},$$
(69)
with $`|G_n|`$ bounded by $`r^nn!`$ at large $`n`$, one defines its Borel transform $`B(z)`$ as
$$B(z)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{G_{n+1}}{n!}z^n,$$
(70)
inside its convergence radius, $`|zr|<1`$, and by analytic continuation elsewhere. Then
$$G(g_{\mathrm{YM}}^2)=_0^{\mathrm{}}𝑑ze^{z/g_{\mathrm{YM}}^2}B(z)$$
(71)
(where the integration contour is on the positive real axis) can provide a resummation of the perturbative series, depending on the structure of the singularities of $`B(z)`$ in the complex plane (Borel plane), and on the convergence properties of the integral at infinity. For asymptotically free non abelian gauge theories a number of singularities in the Borel plane have been identified as follows . IR renormalons give singularities on the positive semiaxis, at
$$\overline{z}=\frac{2}{\beta _0}k,k=2,3,\mathrm{}$$
(72)
($`\beta _0>0`$ for asymptotic freedom; in our case $`\beta _0=\frac{8}{3}N/(8\pi ^2)`$, see eq. (23)). Since these singularities are on the positive semiaxis, we must supplement eq. (71) with a prescription for dealing with them, and this produces non-perturbative contributions of order
$$\mathrm{exp}\{\frac{2k}{\beta _0g_{\mathrm{YM}}^2}\}\frac{1}{E^{2k}},$$
(73)
times regular functions of the energy. Since the first IR renormalon is at $`k=2`$, it gives a contribution $`1/E^4`$. UV renormalons are instead on the negative semiaxis, at
$$\overline{z}=\frac{2}{\beta _0}k,k=1,2,3,\mathrm{}$$
(74)
Even if these singularities are not on the integration contour, they limit the convergence radius of the expansion for $`B(z)`$ and are responsible for effects of the form
$$\mathrm{exp}\{\frac{|\overline{z}|}{g_{\mathrm{YM}}^2}\}\frac{1}{E^{2k}}$$
(75)
However, now $`k`$ starts from 1, so that the leading effect is $`1/E^2`$. There are further singularities on the positive semiaxis due to instantons, but their effect is of order $`\mathrm{exp}(8\pi ^2/g_{\mathrm{YM}}^2)`$, which for $`SU(N)`$ with 6 adjoint scalars is $`(1/E)^{8N/3}`$ (or $`(1/E)^{16N/3}`$ if the first non-vanishing contribution comes from an instanton-antiinstanton pair) and therefore is quite suppressed for large $`N`$.
Comparing our result (66) with the above situation, we see that the $`1/E^2`$ term that we have found matches with the contribution of the first UV renormalon. The fact that we get the correct location of the renormalon singularity in the Borel plane is non-trivial, and can be traced back to the fact that the linearization of eq. (11) gave eq. (48), which has the homogeneous solution $`t(y)=e^{\gamma y}`$ with Re $`\gamma =1/2`$, as long as eq. (64) holds. From the point of view of the gravity computation, we could have obtained $`1/E^\alpha `$ with $`\alpha `$ a priori any real number, not necessarily integer nor rational. We therefore regard this agreement as a successful and non trivial test of the conjectured duality between type 0B theory and non supersymmetric YM theory.
The term $`B_T^2F_2(E)/E^4`$ comes from the iteration of the term $`B_TF_1(E)/E^2`$, as we see from the coefficient $`B_T^2`$, so it clearly corresponds to the UV renormalon with $`k=2`$. Iterating further the contribution of the tachyon in the solution of the equations of motion, we find all the UV renormalons at their correct locations in the Borel plane.
The term $`C/E^4`$ instead has a prefactor which is just equal to the constant $`C`$ and is not related to $`F_1(E)`$ or $`B_T`$, so it has a different physical origin. Thus, we identify it with the $`k=2`$ IR renormalon. Again, it is very remarkable that its position in the Borel plane matches exactly the position of the first ($`k=2`$) IR renormalon; furthermore, we correctly find no term corresponding to an IR renormalon with $`k=1`$, which corresponds to the absence of gauge-invariant operators of dimension 2 in the OPE of current-current correlation functions.
Finally, even the energy dependence of the prefactors is consistent with what is known about renormalons from the analysis of Feynman graphs. In fact, in the approximation in which only a single chain of bubbles is considered, it is known that all UV renormalons are double poles, while the $`k=2`$ IR renormalon is a single pole, and all others IR renormalons are again double poles . For a single pole
$$B(z)\frac{1}{z\frac{4}{\beta _0}},$$
(76)
and then eq. (71) gives a contribution $`\mathrm{exp}\{4/(\beta _0g_{\mathrm{YM}}^2)\}1/E^4`$ without any further energy dependence, in full agreement with our interpretation of $`C/E^4`$ as the contribution of the $`k=2`$ IR renormalon. For a double pole on the positive semiaxis, instead, the non-perturbative contribution is of order
$`G(g_{\mathrm{YM}}^2)`$ $``$ $`{\displaystyle _0^{\mathrm{}}}𝑑ze^{z/g_{\mathrm{YM}}^2}{\displaystyle \frac{1}{(z\overline{z})^2}}`$ (77)
$``$ $`{\displaystyle \frac{1}{g_{\mathrm{YM}}^2}}{\displaystyle _0^{\mathrm{}}}𝑑ze^{z/g_{\mathrm{YM}}^2}{\displaystyle \frac{1}{(z\overline{z})}}{\displaystyle \frac{1}{g_{\mathrm{YM}}^2}}e^{\overline{z}/g_{\mathrm{YM}}^2},`$
where in the second line we have integrated by parts and kept only the non-analytic term. This provides a $`\mathrm{ln}E`$ term in the prefactor. Similarly, there are logarithmic terms from the double poles on the negative semiaxis . Thus in the single chain approximation the first UV renormalon is of order
$$\frac{\beta _0\mathrm{ln}E}{E^2},$$
(78)
and in our approach is reproduced by the first term of the expansion of eq. (67) in powers of $`\beta _T`$ (which suggests that $`\beta _T`$ is proportional to $`\beta _0`$). Higher order terms are expected to come from graphs with multiple chains insertions, and are not well understood in QCD . In this sense, eq. (67) is also a remarkable prediction of the dual string theory, that would be very interesting to check from the YM side.
At low energies $`F_1(E)`$ oscillates faster and faster, see (fig. 8). This is physically very reasonable, since these oscillations signal the transition from asymptotic freedom to the confining regime. Thus, the form of $`F_1(E)`$ seems to capture both the correct single chain result and a physically correct behaviour when the IR is approached.
We conclude this section with a general comment. The rule concerning the position of renormalons in field theory is that we find on the positive semiaxis the renormalons corresponding to strongly coupled physics. For asymptotically free gauge theories IR renormalons are at $`z>0`$, while UV renormalons are at $`z<0`$. In QED or $`\lambda \varphi ^4`$ the coupling becomes strong in the UV and correspondingly the UV renormalons are at $`z>0`$ and IR renormalons are at $`z<0`$. When a singularity is at $`z>0`$ we must give a prescription for going around it; if we understand the physics associated to this singularity, we can give a physically motivated prescription; then we are left with a non-perturbative contribution, but no ambiguity on the resummation of the perturbative series. This is the case of instantons . If however we do not understand the physics behind the singularity, we do not know what prescription to use. Different prescriptions (e.g. going around a pole from above or from below) give answers that differ by terms of order $`\mathrm{exp}(2k/\beta _0g_{\mathrm{YM}}^2)`$, so renormalons on the positive semiaxis are usually taken as a source of an uncertainty of this order in the resummation of the perturbative series<sup>6</sup><sup>6</sup>6Actually, in QCD even UV renormalons can give rise to uncertainties on the perturbative series which from the Feynman graphs point of view are uncalculable, despite the fact that the UV physics is well understood ..
In our case, however, the coefficients $`B_T,C`$ in eq. (66) are fixed from the comparison with the numerical solution. Once we specify the IR solution, the values of $`B_T,C`$, etc. follow, and the renormalon contribution, rather than an uncertainty on the perturbative series, becomes a well defined non-perturbative contribution. The fact that IR physics fixes the renormalon contributions confirms the old conjecture of ’t Hooft that renormalon singularities are related to the quark confinement mechanism. It is interesting to note that instead $`B_T,C`$ cannot be fixed, even in principle, from the UV side, since there they just appear as integration constants of the homogeneous equations (52).
## 7 $`\alpha ^{}`$ and string loop corrections
We next consider the effect of $`\alpha ^{}`$ corrections on the solutions. For the type 0 theory, these have been discussed in ref. , where it is found that they have the same structure as in type II theories; in particular, world-sheet supersymmetry implies the vanishing of $`\alpha ^{}(\mathrm{Riemann})^2`$ corrections, and the first non-vanishing correction is $`\alpha ^3(\mathrm{Riemann})^4`$. It is well known that $`\alpha ^{}`$ corrections suffer from a certain degree of ambiguity, because not all the coefficients of the possible operators are fixed by the comparison with the string amplitudes; a related ambiguity is the fact that one can perform field redefinitions that mix different orders in $`\alpha ^{}`$, e.g. $`g_{\mu \nu }g_{\mu \nu }+\alpha ^{}R_{\mu \nu }`$. If one has an exactly conformal background these field redefinitions do not matter, but of course if one works at a finite order in the $`\alpha ^{}`$ expansion they make a difference.
An important observation is that we can choose the $`\alpha ^{}`$ correction such that they depend only on the Weyl tensor, rather than on the Riemann tensor, and then they have a much milder behaviour when we approach $`AdS_5`$. To be quantitative, consider eq. (8), which, including the first non-vanishing $`\alpha ^{}`$ correction, has the general form
$$\ddot{\varphi }+\frac{1}{2}g(T)e^{\frac{\varphi +\xi }{2}5\eta }=ce^{\frac{3\varphi +\xi }{2}5\eta }(\mathrm{Weyl})^4,$$
(79)
where $`c`$ is a constant, (Weyl)<sup>4</sup> denotes the appropriate contractions of four Weyl tensors $`C_{MNRS}`$, and again we have set $`\alpha ^{}=1`$. Writing $`\ddot{\varphi }=e^{2y}(\varphi ^{\prime \prime }+\varphi ^{})`$, where as before the dot is $`d/d\rho `$ and the prime is $`d/dy`$, with $`\rho =e^y`$, we have
$$\varphi ^{\prime \prime }+\varphi ^{}+\frac{1}{2}g(T)e^{\frac{\varphi +\xi }{2}5\eta 2y}=ce^{\frac{3\varphi +\xi }{2}5\eta 2y}(\mathrm{Weyl})^4.$$
(80)
We see that an ansatz of the form (16)-(19), or its improved form (42)-(45), is still a solution, but the numerical coefficients in front of the factors $`\mathrm{ln}2`$ change. In fact, consider the ansatz $`\varphi =2\mathrm{ln}y+a_1\mathrm{ln}2+O(\mathrm{ln}y/y),\xi =y+a_2\mathrm{ln}2+O(1/y),\eta =y/2+a_3\mathrm{ln}2+O(1/y),T=1+O(1/y)`$. We have checked that on this ansatz all components of the Weyl tensor $`C_{MNRS}`$ are $`O(1/y)`$, so that
$$(\mathrm{Weyl})^4\frac{1}{y^4}.$$
(81)
Substituting into eq. (80) we get (setting again $`g(T)=T^2/2`$)
$$\frac{2}{y}+\frac{1}{4y}e^{(\frac{a_1+a_2}{2}5a_3)\mathrm{ln}2}+O(\frac{1}{y^2})=\frac{\mathrm{const}.}{y}e^{(\frac{3a_1+a_2}{2}5a_3)\mathrm{ln}2}+O(\frac{1}{y^2}),$$
(82)
and similarly for the equations for $`\xi ,\eta ,T`$. Thus, the leading $`y`$ dependence cancels, and we can adjust the coefficients $`a_1,a_2,a_3`$ so that the equations are satisfied. We see that the numerical values of the coefficients $`a_1,a_2,a_3`$ are affected by the inclusion of the (Weyl)<sup>4</sup> term. The logarithmic behaviour in the UV is not spoiled, but the numerical values of the beta function coefficients change .
We now want to see whether the terms related to renormalons are affected. If we use again the improved ansatz (42)-(45), where now $`b`$ and the coefficients of $`\mathrm{ln}2`$ in $`\xi ,\eta `$ are modified according to the discussion above, we see that, in the limit $`by1`$, where $`b/(A+by)1/y`$, the net effect of the addition of the (Weyl)<sup>4</sup> term is to change the numerical factors inside the square brackets on the right-hand side of eqs. (46)–(49). However, the associated homogeneous equations are still given by eq. (52), and are unaffected by the inclusion of the (Weyl)<sup>4</sup> term. Therefore, the position of the renormalon singularities in the Borel plane is unaffected by these $`\alpha ^3`$ corrections. It is clear that the argument goes through at all orders in $`\alpha ^{}`$ as long as the corrections can be written in terms of the Weyl tensor. As for the string loop corrections, in the UV they are small by definition, since $`e^\varphi 1`$. Thus, it appears that the result on the position of renormalon singularities is quite robust.
In the IR the situation is more complicated, first of all because unavoidably when $`g_{\mathrm{YM}}`$ becomes strong also the string loop corrections become important. This is different from the standard AdS/CFT correspondence, where the dilaton is constant and $`g_s=e^\varphi `$ can be taken to be small everywhere.
The role of $`\alpha ^{}`$ corrections in the IR has been discussed in ref. , and it has been found that they depend crucially on the numerical relations between the coefficients $`\xi _1,\varphi _1,\eta _1`$ that parametrize the IR solution (25). With the choice $`\xi _1=\varphi _1`$ that we used in sect. 4, the corrections are indeed large in the IR, while, for instance, for $`\varphi _1=3\xi _1,t_1=0`$ we still have a confining solution but now the $`\alpha ^{}`$ corrections in the IR are exponentially small in $`\rho `$. Therefore we have repeated the numerical analysis discussed in sect. 4, for the case $`\varphi _1=3\xi _1,\xi _1=1,t_1=0`$ and, from eq. (27), $`\eta _1=1`$. We have found that the lowest order solution for the dilaton and the tachyon shows no qualitative difference, see fig. (9).
The functions $`\eta ,\xi `$ are also qualitatively very similar to the previous case; the main difference is that, since now $`\xi _1/2\eta _1<0`$, the radius of the 5-sphere $`R_{(5)}^2=e^{\frac{\xi }{2}\eta }`$ goes to zero in the IR, see fig. (10) (using eq. (27), we see that this happens in general if $`\varphi _1^2>(3/2)\xi _1^2`$). This suggests that the six adjoint scalars become massive, but to clarify this point we need a better understanding of string loop corrections in the IR limit.
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# On Commuting and Non-Commuting Complexes
## 1 Introduction
Given a finite group $`G`$, one defines a non-commuting set to be a set of elements $`\{g_1,\mathrm{},g_n\}`$ such that $`g_i`$ does not commute with $`g_j`$ for $`ij`$. The sizes of maximal non-commuting sets in a group are interesting invariants of the group. In particular, the largest integer $`n`$ such that the group $`G`$ has a non-commuting set of order $`n`$, which is denoted by $`nc(G)`$, is known to be closely related to other invariants of $`G`$. For example if $`k(G)`$ is the size of the largest conjugacy class in $`G`$ then
$$k(G)4(nc(G))^2\text{(See }\text{[P]}\text{.)}$$
Also if we define $`cc(G)`$ to be the minimal number of abelian subgroups of $`G`$ that covers $`G`$, then I. M. Isaacs (see \[J\]) has shown that
$$nc(G)cc(G)(nc(G)!)^2.$$
Confirming a conjecture of P. Erdös, L. Pyber \[P\] has also shown that there is a positive constant $`c`$ such that
$$cc(G)|G:Z(G)|c^{nc(G)}$$
for all groups $`G`$. Another interesting place where the invariant $`nc(G)`$ appears is in the computation of the cohomology length of extra-special $`p`$-groups (See \[Y\].)
In this paper we study the topology of certain complexes associated to the poset of non-commuting sets in a group $`G`$. Let $`NC(G)`$ be the complex whose vertices are just the nontrivial elements of the group $`G`$ and whose faces are the noncommuting sets in $`G`$. The central elements form point components in this complex and are not so interesting. So, we look at the subcomplex $`BNC(G)`$ formed by non-central elements of $`G`$. We show
###### Result 1 3.2).
If $`G`$ is a nonabelian group, then $`BNC(G)`$ is simply-connected.
In general we also notice that $`BNC(G)`$ is equipped with a free $`Z(G)`$-action where $`Z(G)`$ is the center of $`G`$. It is also equipped in general with a $`/2`$-action whose fixed point set is exactly $`BNC_2(G)`$, the corresponding complex where we use only the elements of order 2 (involutions). Thus if $`G`$ is an odd order group or if $`G`$ has nontrivial center, then the Euler characteristic of $`BNC(G)`$ is not $`1`$ and so it is not contractible.
On a more refined level, we use a recent simplicial decomposition result of A. Björner, M. Wachs and V. Welker to show that there is a simplicial complex $`S`$, called the core of $`BNC(G)`$, so that the following decomposition formula holds:
###### Result 2 4.11).
If $`G`$ is a finite nonabelian group and $`S`$ is the core of $`BNC(G)`$, then
$$BNC(G)S\underset{FS}{}[Susp^{|F|}(LkF)]^{\gamma (F)}$$
where the $`F`$ are the faces of $`S`$, $`Lk`$ stands for the link of a face, $`Susp^k`$ stands for a $`k`$-fold suspension and $``$ stands for wedge product. The number $`\gamma (F)=_{[x]F}(m_x1)`$ where $`m_x`$ is the size of the centralizer class $`[x]`$.
It is clear from this decomposition formula that when the core of $`BNC(G)`$ is contractible, then $`BNC(G)`$ is a wedge of suspensions of spaces and hence has a trivial ring structure on its cohomology. This is true, for example, when $`G=\mathrm{\Sigma }_p`$, the symmetric group on $`p`$ letters, for some prime $`p`$.
The following is an important consequence of the above decomposition formula:
###### Result 3 4.14).
Let $`G`$ be a finite nonabelian group, and let $`S_s`$ denote the set of maximal non-commuting sets in $`G`$ of size $`s`$. Then, for $`s>1`$,
$$rk(H_{s1}(BNC(G)))\underset{FS_s}{}[\underset{xF}{}(1\frac{1}{m_x})]$$
where $`m_x`$ is the size of the centralizer class containing $`x`$.
In particular, if $`G`$ is an odd order group or if $`G`$ has a nontrivial center $`(|G|2)`$, then
$$\stackrel{~}{H}_{s1}(NC(G))0$$
whenever $`G`$ has a maximal non-commuting set of size $`s`$.
There is also a $`p`$-local version of this theorem which, in particular, gives that $`\stackrel{~}{H}_{s1}(NC_p(G))0`$ whenever $`G`$ has a maximal non-commuting $`p`$-set of size $`s`$ and $`p`$ is an odd prime. For $`p=2`$, the same is true under the condition $`2||Z(G)|`$ and $`|G|2`$. Observe that this result has a striking formal similarity (in terms of their conclusions) to the following theorem proved by Quillen (theorem 12.1 in \[Q\]):
###### Theorem 1.1 (Quillen).
If $`G`$ is a finite solvable group having no nontrivial normal $`p`$-subgroup, then
$$\stackrel{~}{H}_{s1}(A_p(G))0$$
whenever $`G`$ has a maximal elementary abelian $`p`$-group of rank $`s`$.
A nice consequence of Result 3 can be stated as follows: If $`G`$ is a group of odd order or a group with nontrivial center such that $`BNC(G)`$ is spherical, i.e., homotopy equivalent to a wedge of equal dimensional spheres, then all maximal non-commuting sets in $`G`$ have the same size.
A natural question to ask is: For which groups is $`BNC(G)`$ spherical? As a partial answer, we show that if $`G`$ is a group where the commutation relation is transitive, then $`BNC(G)`$ is spherical. We give examples of such groups and compute $`BNC(G)`$ for these groups.
Observe that one could also define a commuting complex $`C(G)`$ analogous to the way we defined $`NC(G)`$ by making the faces consist of commuting sets of elements instead of non-commuting sets. However, this definition does not provide us with new complexes. For example, $`C_p(G)`$, the commuting complex formed by the elements of prime order $`p`$, is easily shown to be $`G`$-homotopy equivalent to Quillen’s complex $`A_p(G)`$. However, the definition of $`C_p(G)`$ helps us to see a duality between $`NC_p(G)`$ and $`C_p(G)`$ where $`NC_p(G)`$ is the subcomplex of $`NC(G)`$ spanned by the vertices which correspond to elements of order $`p`$. Using a result of Quillen on $`A_p(G)`$, we obtain:
###### Result 4 5.2).
Let $`G`$ be a group and $`p`$ a prime with $`p||G|`$. Pick a Sylow $`p`$-subgroup $`P`$ of $`G`$ and define $`N`$ to be the subgroup generated by the normalizers $`N_G(H)`$ where $`H`$ runs over all the nontrivial subgroups of $`P`$.
Then $`NC_p(G)`$ is $`(|G:N|2)`$-connected. In fact $`NC_p(G)`$ is the $`|G:N|`$-fold join of a complex $`S`$ with itself where $`S`$ is “dual” to a path-component of $`A_p(G)`$.
Finally under suitable conditions, $`BNC(G)`$ is shellable and this yields some combinatorial identities. As an application we obtain
###### Result 5 4.24).
Let $`G`$ be a nonabelian group with a transitive commuting relation, i.e., if $`[g,h]=[h,k]=1`$, then $`[g,h]=1`$ for every noncentral $`g,h,kG`$. Then,
$$nc(G)(nc(G)1)+|G|(|G|m)2(nc(G)1)(|G||Z(G)|)0.$$
where $`m`$ denotes the number of conjugacy classes in $`G`$.
## 2 Background
We start the section with a discussion of complexes associated with posets of subgroups of a group $`G`$. For a complete account of these well known results, see chapter 6 in \[B\].
Given a finite poset $`(P,)`$, one can construct a simplicial complex $`|P|`$ out of it by defining the $`n`$-simplices of $`|P|`$ to be chains in $`P`$ of the form $`p_0<p_1<\mathrm{}<p_n`$. This is called the simplicial realization of the poset $`P`$.
Furthermore any map of posets $`f:(P_1,_1)(P_2,_2)`$ (map of posets means $`x_1_1x_2f(x_1)_2f(x_2)`$) yields a simplicial map between $`|P_1|`$ and $`|P_2|`$ and hence one has in general a covariant functor from the category of finite posets to the category of finite simplicial complexes and simplicial maps. Thus if a (finite) group $`G`$ acts on a poset $`P`$ via poset maps (we say $`P`$ is a $`G`$-poset) then $`G`$ will act on $`|P|`$ simplicially.
Brown, Quillen, Webb, Bouc, Thévenaz and many others constructed many finite $`G`$-simplicial complexes associated to a group $`G`$ and used them to study the group $`G`$ and its cohomology. In particular, the following posets of subgroups of $`G`$ have been studied extensively:
(a) The poset $`s_p(G)`$ of nontrivial $`p`$-subgroups of $`G`$.
(b) The poset $`a_p(G)`$ of nontrivial elementary abelian $`p`$-subgroups of $`G`$.
(c) The poset $`b_p(G)`$ of nontrivial $`p`$-radical subgroups of $`G`$. (Recall a $`p`$-radical subgroup of $`G`$ is a $`p`$-subgroup $`P`$ of $`G`$ such that $`PN_G(P)/P`$ has no nontrivial normal $`p`$-subgroups.)
$`G`$ acts on each of these posets by conjugation and thus from each of these $`G`$-posets, one gets a $`G`$-simplicial complex $`S_p(G)`$, $`A_p(G)`$ and $`B_p(G)`$ respectively. $`S_p(G)`$ is usually called the Brown complex of $`G`$ and $`A_p(G)`$ is usually called the Quillen complex of $`G`$ where the dependence on the prime $`p`$ is understood. Notice again that the trivial subgroup is not included in any of these posets, since if it were the resulting complex would be a cone and hence trivially contractible.
It was shown via work of Quillen and Thévenaz, that $`S_p(G)`$ and $`A_p(G)`$ are $`G`$-homotopy equivalent and via work of Bouc and Thévenaz that $`B_p(G)`$ and $`S_p(G)`$ are $`G`$-homotopy equivalent. Thus, in a sense, these three $`G`$-complexes capture the same information.
Recall the following elementary yet very important lemma (see \[B\]):
###### Lemma 2.1.
If $`f_0,f_1:P_1P_2`$ are two maps of posets such that $`f_0(x)f_1(x)`$ for all $`xP_1`$ then the simplicial maps induced by $`f_0`$ and $`f_1`$ from $`|P_1|`$ to $`|P_2|`$ are homotopic.
Using this, Quillen made the following observation, if $`P_0`$ is a nontrivial normal $`p`$-subgroup of $`G`$, then we may define a poset map $`f:s_p(G)s_p(G)`$ by $`f(P)=P_0P`$ and by the lemma above, $`f`$ would be homotopic to the identity map, but on the other hand since $`f(P)`$ contains $`P_0`$ always, again by the lemma, $`f`$ is also homotopic to a constant map. Thus we see that $`S_p(G)`$ is contractible in this case. Quillen then conjectured
###### Conjecture 2.2 (Quillen).
If $`G`$ is a finite group, $`S_p(G)`$ is contractible if and only if $`G`$ has a nontrivial normal $`p`$-subgroup.
He proved his conjecture in the case that $`G`$ is solvable but the general conjecture remains open. Notice though that if $`S_p(G)`$ is $`G`$-homotopy equivalent to a point space then this does imply $`G`$ contains a nontrivial normal $`p`$-subgroup since in this case $`S_p(G)^G`$ is homotopy equivalent to a point which means in particular $`S_p(G)^G`$ is not empty, yielding a nontrivial normal $`p`$-subgroup.
The purpose of this paper is to introduce some simplicial complexes associated to elements of a group rather than subgroups of a group and use these to give a different perspective on some of the complexes above.
For this purpose we make the following definitions:
###### Definition 2.3.
Let $`G`$ be a group. Define a simplicial complex $`C(G)`$ by declaring a $`n`$-simplex in this complex to be a collection $`[g_0,g_1,\mathrm{},g_n]`$ of distinct nontrivial elements of $`G`$, which pairwise commute.
Similarly define a simplicial complex $`NC(G)`$ by declaring a $`n`$-simplex to be a collection $`[g_0,g_1,\mathrm{},g_n]`$ of nontrivial elements of $`G`$, which pairwise do not commute.
It is trivial to verify that the above definition, does indeed define complexes on which $`G`$ acts simplicially by conjugation.
Usually when one studies simplicial group actions, it is nice to have admissible actions, i.e., actions where if an element of $`G`$ fixes a simplex, it actually fixes it pointwise. Although $`C(G)`$ and $`NC(G)`$ are not admissible in general, one can easily fix this by taking a barycentric subdivision. The resulting complex is of course $`G`$-homotopy equivalent to the original, however it now is the realization of a poset.
Thus if we let $`PC(G)`$ be the barycentric subdivision of $`C(G)`$, it corresponds to the realization of the poset consisting of subsets of nontrivial, pairwise commuting elements of $`G`$, ordered by inclusion. Similarly if we let $`PNC(G)`$ be the barycentric subdivision of $`NC(G)`$, it corresponds to the realization of the poset consisting of subsets of nontrivial, pairwise non-commuting elements of $`G`$, ordered by inclusion.
Depending on the situation, one uses either the barycentric subdivision or the original. For purposes of understanding the topology, the original is easier but for studying the $`G`$-action, the subdivision is easier.
Of course, we will want to work a prime at a time also, so we introduce the following $`p`$-local versions of $`C(G)`$ and $`NC(G)`$.
###### Definition 2.4.
Given a group $`G`$ and a prime $`p`$, let $`C_p(G)`$ be the subcomplex of $`C(G)`$ where the simplices consist of sets of nontrivial, pairwise commuting elements of order $`p`$.
Similarly let $`NC_p(G)`$ be the subcomplex of $`NC(G)`$ where the simplices consist of sets of nontrivial, pairwise non-commuting elements of order $`p`$
Of course the same comments about the $`G`$-action and the barycentric subdivision apply to these $`p`$-local versions.
Our first order of business is to see that the commuting complexes $`C(G)`$ and $`C_p(G)`$ are nothing new. We will find the following standard lemma useful for this purpose (see \[B\]):
###### Lemma 2.5.
If $`f`$ is a $`G`$-map between admissible $`G`$-simplicial complexes $`X`$ and $`Y`$ with the property that for all subgroups $`HG`$, f restricts to an ordinary homotopy equivalence between $`X^H`$ and $`Y^H`$ (recall $`X^H`$ is the subcomplex of $`X`$ which consists of elements fixed pointwise by $`H`$), then $`f`$ is a $`G`$-homotopy equivalence, i.e., there is a $`G`$-map $`g:YX`$ such that $`fg`$ and $`gf`$ are $`G`$-homotopic to identity maps.
###### Theorem 2.6.
Let $`G`$ be a finite group, then $`C(G)`$ is $`G`$-homotopy equivalent to the simplicial realization of the poset $`A(G)`$ of nontrivial abelian subgroups of $`G`$, ordered by inclusion and acted on by conjugation.
Furthermore if $`p`$ is a prime, then $`C_p(G)`$ is $`G`$-homotopy equivalent to $`A_p(G)`$ (and thus $`S_p(G)`$ and $`B_p(G)`$.)
###### Proof.
First we will show homotopy equivalence and remark on $`G`$-homotopy equivalence later.
We work with $`PC(G)`$, the barycentric subdivision. Notice that the associated poset of $`PC(G)`$ contains the poset $`A(G)`$ of nontrivial abelian subgroups of $`G`$ as a subposet, they are merely the commuting sets whose elements actually form an abelian subgroup (minus identity). Let $`i:A(G)PC(G)`$ denote this inclusion.
We now define a poset map $`r:PC(G)A(G)`$ as follows: If $`S`$ is a set of nontrivial, pairwise commuting elements of $`G`$, then $`<S>`$, the subgroup generated by $`S`$ will be a nontrivial abelian subgroup of $`G`$, thus we can set $`r(S)=\{<S>1\}`$. It is obvious that $`r`$ is indeed a poset map, and that $`Sr(S)`$ and so $`ir`$ is homotopic to the identity map of $`PC(G)`$ by lemma 2.1. Furthermore it is clear that $`ri=Id`$ and so $`r`$ is a deformation retraction of $`PC(G)`$ onto $`A(G)`$.
Thus $`PC(G)`$ is homotopy equivalent to $`A(G)`$. To see this is a $`G`$-homotopy equivalence, we just need to note that $`r`$ is indeed a $`G`$-map, and maps a commuting set invariant under conjugation by a subgroup $`H`$ into a subgroup invariant under conjugation by $`H`$ and thus induces a homotopy equivalence between $`PC(G)^H`$ and $`A(G)^H`$ for any subgroup $`H`$. Thus $`r`$ is indeed a $`G`$-homotopy equivalence by lemma 2.5.
The $`p`$-local version follows exactly in the same manner, once one notes that the subgroup generated by a commuting set of elements of order $`p`$ is an elementary abelian $`p`$-subgroup. ∎
Thus we see from theorem 2.6, that the commuting complexes at a prime $`p`$ are basically the $`A_p(G)`$ in disguise. However for the rest of the paper, we look at the non-commuting complexes and we will see that they are quite different, and in some sense dual to the commuting ones. However before doing that, we conclude this section by looking at a few more properties of the commuting complex.
Recall the following important proposition of Quillen \[Q\]:
###### Proposition 2.7.
If $`f:XY`$ is a map of posets, and $`yY`$ we define
$`\begin{array}{cc}\hfill f/y& =\{xX|f(x)y\}\hfill \\ \hfill y\backslash f& =\{xX|f(x)y\}.\hfill \end{array}`$
Then if $`f/y`$ is contractible for all $`yY`$ (respectively $`y\backslash f`$ is contractible for all $`yY`$) then $`f`$ is a homotopy equivalence between $`|X|`$ and $`|Y|`$.
Using this we will prove:
###### Proposition 2.8.
If $`G`$ is a group with nontrivial center then $`A(G)`$ and hence $`C(G)`$ is contractible. Moreover, $`A(G)`$ is homotopy equivalent to $`Nil(G)`$ where $`Nil(G)`$ is the poset of nontrivial nilpotent subgroups of $`G`$.
###### Proof.
If $`G`$ has a nontrivial center $`Z(G)`$, then $`A(G)`$ is conically contractible via $`AAZ(G)Z(G)`$ for any abelian subgroup $`A`$ of $`G`$. Thus $`C(G)`$ is also contractible as it is homotopy equivalent to $`A(G)`$.
Let $`i:A(G)Nil(G)`$ be the natural inclusion of posets. Take $`NNil(G)`$ and let us look at $`i/N=\{BA(G)|BN\}=A(N)`$. However since $`N`$ is nilpotent, it has a nontrivial center $`Z`$ and hence $`A(N)`$ is conically contractible. Thus by proposition 2.7 the result follows.
It is natural to ask if:
###### Conjecture 2.9.
$`C(G)`$ is contractible if and only if $`G`$ has a nontrivial center.
## 3 Non-commuting complexes
Fix a group $`G`$, let us look at $`NC(G)`$. The first thing we notice is that any nontrivial central element in $`G`$ gives us a point component in $`NC(G)`$ and hence is not so interesting. Thus we define:
###### Definition 3.1.
$`BNC(G)`$ is the subcomplex of $`NC(G)`$ consisting of those simplices of $`NC(G)`$ which are made out of noncentral elements. Thus $`BNC(G)`$ is empty if $`G`$ is an abelian group.
The first thing we will show is that if $`G`$ is a nonabelian group, (so that $`BNC(G)`$ is nonempty) then $`BNC(G)`$ is not only path-connected but it is simply-connected. Notice also that this means the general picture of $`NC(G)`$ is as a union of components, with at most one component of positive dimension and this is $`BNC(G)`$ and it is simply-connected. Also notice that $`BNC(G)`$ is invariant under the conjugation $`G`$-action, and the point components of $`NC(G)`$ are just fixed by the $`G`$-action as they correspond to central elements.
###### Theorem 3.2.
If $`G`$ is a nonabelian group, then $`BNC(G)`$ is a simply-connected $`G`$-simplicial complex.
###### Proof.
First we show that $`BNC(G)`$ is path-connected. Take any two vertices in $`BNC(G)`$, call them $`g_0`$ and $`g_1`$, then these are two noncentral elements of $`G`$, thus their centralizer groups $`C(g_0)`$ and $`C(g_1)`$ are proper subgroups of $`G`$.
It is easy to check that no group is the union of two proper subgroups for suppose $`G=HK`$ where $`H,K`$ are proper subgroups of $`G`$. Then we can find $`hGK`$ (it follows $`hH`$) and $`kGH`$ (hence $`kK`$). Then $`hk`$ is not in $`H`$ as $`hH`$ and $`kH`$. Similarly $`hkK`$ so $`hkHK=G`$ which is an obvious contradiction. Thus no group is the union of two proper subgroups.
Thus we conclude that $`C(g_0)C(g_1)G`$ and so we can find an element $`w`$ which does not commute with either $`g_0`$ or $`g_1`$ and so the vertices $`g_0`$ and $`g_1`$ are joined by an edge path $`[g_0,w]+[w,g_1]`$. (The $`+`$ stands for concatenation.) Thus we see $`BNC(G)`$ is path-connected. In fact, any two vertices of $`BNC(G)`$ can be connected by an edge path involving at most two edges of $`BNC(G)`$.
To show it is simply-connected, we argue by contradiction. If it was not simply connected, then there would be a simple edge loop which did not contract, i.e., did not bound a suitable union of 2-simplices. (a simple edge loop is formed by edges of the simplex and is of the form $`L=[e_0,e_1]+[e_1,e_2]+\mathrm{}+[e_{n1},e_n]`$ where all the $`e_i`$ are distinct except $`e_0=e_n`$.)
Take such a loop $`L`$ with minimal size $`n`$. (Notice $`n`$ is just the number of edges involved in the loop.)
Since we are in a simplicial complex, certainly $`n3`$.
Suppose $`n>5`$, then $`e_3`$ can be connected to $`e_0`$ by an edge path $`E`$ involving at most two edges by our previous comments. This edge path $`E`$ breaks our simple edge loop into two edge loops of smaller size which hence must contract since our loop was minimal. However, then it is clear that our loop contracts which is a contradiction so $`n5`$.
So we see $`3n5`$ so we have three cases to consider:
(a) $`n=3`$:
Here $`L=[e_0,e_1]+[e_1,e_2]+[e_2,e_3]`$ with $`e_3=e_0`$. But then it is easy to see $`\{e_0,e_1,e_2\}`$ is a set of pairwise non-commuting elements and so gives us a 2-simplex $`[e_0,e_1,e_2]`$ in $`BNC(G)`$ which bounds the loop which gives a contradiction.
(b) $`n=4`$:
Here $`L=[e_0,e_1]+[e_1,e_2]+[e_2,e_3]+[e_3,e_4]`$ with $`e_4=e_0`$. Thus $`L`$ forms a square. Notice by the simplicity of $`L`$, the diagonally opposite vertices in the square must not be joined by an edge in $`BNC(G)`$, i.e., they must commute, thus $`e_0`$ commutes with $`e_2`$ and $`e_1`$ commutes with $`e_3`$.
Since $`e_0`$ and $`e_1`$ do not commute, $`\{e_0,e_1,e_0e_1\}`$ is a set of mutually non-commuting elements and so forms a 2-simplex of $`BNC(G)`$. Since $`e_2`$ commutes with $`e_0`$ but not with $`e_1`$, it does not commute with $`e_0e_1`$ and thus $`\{e_0e_1,e_1,e_2\}`$ also is a 2-simplex in $`BNC(G)`$. Similar arguments show that $`\{e_0e_1,e_0,e_3\}`$ and $`\{e_0e_1,e_2,e_3\}`$ form 2-simplices in $`BNC(G)`$. The union of the four $`2`$-simplices mentioned in this paragraph, bound our loop giving us our contradiction. Thus we are reduced to the final case:
(c) $`n=5`$:
Here $`L=[e_0,e_1]+[e_1,e_2]+[e_2,e_3]+[e_3,e_4]+[e_4,e_5]`$ with $`e_5=e_0`$. Thus $`L`$ forms a pentagon, and again by simplicity of $`L`$, nonadjacent vertices in the pentagon cannot be joined by an edge in $`BNC(G)`$, thus they must commute. Similar arguments as for the $`n=4`$ case yield that $`[e_0,e_1,e_0e_1]`$, $`[e_0e_1,e_1,e_2]`$ and $`[e_0e_1,e_0,e_4]`$ are $`2`$-simplices in $`BNC(G)`$ and the union of these three simplices contract our loop $`L`$ into one of length four namely $`[e_0e_1,e_2]+[e_2,e_3]+[e_3,e_4]+[e_4,e_0e_1]`$ which by our previous cases, must contract thus yielding the final contradiction. ∎
From theorem 3.2, we see that $`BNC(G)`$ is simply-connected for any nonabelian group $`G`$. One might ask if it is contractible? The answer is no in general although there are groups $`G`$ where it is contractible. We look at these things next:
###### Proposition 3.3.
For a general finite nonabelian group $`G`$, the center $`Z(G)`$ of $`G`$ acts freely on $`BNC(G)`$ by left multiplication and hence $`|Z(G)|`$ divides the Euler characteristic of $`BNC(G)`$.
For a group of odd order, the simplicial map $`A`$ which maps a vertex $`g`$ to $`g^1`$ is a fixed point free map on $`BNC(G)`$ and on $`NC(G)`$ of order $`2`$. Thus the Euler characteristic of both $`BNC(G)`$ and $`NC(G)`$ is even in this case.
Thus if $`BNC(G)`$ is $`𝔽`$-acyclic for some field $`𝔽`$, then $`G`$ must have trivial center and be of even order.
###### Proof.
The remarks about Euler characteristics follow from the fact that if a finite group $`H`$ acts freely on a space where the Euler characteristic is defined, then $`|H|`$ must divide the Euler characteristic. So we will concentrate mainly on finding such actions.
First for the action of $`Z(G)`$. $`aZ(G)`$ acts by taking a simplex $`[g_0,\mathrm{},g_n]`$ to a simplex $`[ag_0,\mathrm{},ag_n]`$. Notice this is well-defined since $`ag_i`$ is noncentral if $`g_i`$ is noncentral and since $`ag_i`$ commutes with $`ag_j`$ if and only if $`g_i`$ commutes with $`g_j`$.
Furthermore if $`a`$ is not the identity element, this does not fix any simplex $`[g_0,\mathrm{},g_n]`$ since if the set $`\{g_0,\mathrm{},g_n\}`$ equals the set $`\{ag_0,\mathrm{},ag_n\}`$ then $`ag_0`$ is one of the $`g_j`$’s but $`ag_0`$ commutes with $`g_0`$ so it would have to be $`g_0`$ but $`ag_0=g_0`$ gives $`a=1`$, a contradiction.
Thus this action of nonidentity central $`a`$ does not fix any simplex and so we get a free action of $`Z(G)`$ on $`BNC(G)`$.
Now for the action of $`A`$, first note that $`A`$ is well-defined since if $`[g_0,\mathrm{},g_n]`$ is a set of mutually non-commuting elements of $`G`$, so is $`[g_0^1,\mathrm{},g_n^1]`$. Clearly $`AA=Id`$. Furthermore, if the two sets above are equal, then $`g_0^1`$ would have to be one of the $`g_j`$. But since $`g_0^1`$ commutes with $`g_0`$ it would have to be $`g_0`$, i.e., $`g_0`$ would have to have order $`2`$. Similarly, all the $`g_i`$ would have to have order 2. Thus in a group of odd order, $`A`$ would not fix any simplex of $`BNC(G)`$ or $`NC(G)`$, and hence would not have any fixed points.
The final comment is to recall that if $`BNC(G)`$ were $`𝔽`$-acyclic, its Euler characteristic would be 1 and hence the center of $`G`$ would be trivial and $`G`$ would have to have even order by the facts we have shown above. ∎
Proposition 3.3 shows that $`BNC(G)`$ is not contractible if $`G`$ is of odd order or if $`G`$ has a nontrivial center, for example if $`G`$ were nilpotent. There is a corresponding $`p`$-local version which we state next:
###### Proposition 3.4.
If $`2||Z(G)|`$ or if $`G`$ has odd order then $`NC_2(G)`$ is not contractible, in fact the Euler characteristic is even. $`NC_p(G)`$ is never contractible for any odd prime $`p`$, in fact it always has even Euler characteristic.
###### Proof.
Follows from the proof of proposition 3.3, once we note that left multiplication by a central element of order $`2`$ takes the subcomplex $`NC_2(G)`$ of $`NC(G)`$ to itself and that the map $`A`$ maps $`NC_p(G)`$ into itself, as the inverse of an element has the same order as the element. For odd primes $`p`$, $`A`$ is fixed point free on $`NC_p(G)`$ always as no elements of order 2 are involved in $`NC_p(G)`$. ∎
The fact that $`BNC(G)`$ can be contractible sometimes is seen in the next proposition.
###### Proposition 3.5.
If $`G`$ is a nonabelian group with a self-centralizing involution i.e., an element $`x`$ of order 2 such that $`C(x)=\{1,x\}`$, then $`BNC(G)`$ is contractible. In fact, $`BNC(G)=NC(G)`$ in this case.
Thus for example, $`NC(\mathrm{\Sigma }_3)=BNC(\mathrm{\Sigma }_3)`$ is contractible where $`\mathrm{\Sigma }_3`$ is the symmetric group on 3 letters.
###### Proof.
Since $`x`$ does not commute with any nontrivial element, the center of $`G`$ is trivial and $`NC(G)=BNC(G)`$. Furthermore, it is clear that $`BNC(G)`$ is a cone with $`x`$ as vertex and hence is contractible. ∎
To help show $`NC(G)`$ of a group is not contractible, we note the following observation which uses Smith theory. (See \[B\]).
###### Proposition 3.6.
If $`G`$ is a group and $`NC(G)`$ is $`𝔽_2`$-acyclic where $`𝔽_2`$ is the field with two elements, then $`NC_2(G)`$ is also $`𝔽_2`$-acyclic. Furthermore one always has $`\chi (NC(G))=\chi (NC_2(G))`$ mod $`2`$.
###### Proof.
We first recall that the map $`A`$ from proposition 3.3 has order 2 as a map of $`NC(G)`$. However it might have fixed points, in fact from the proof of that proposition, we see that $`A`$ fixes a simplex $`[g_0,\mathrm{},g_n]`$ of $`NC(G)`$ if and only if each element $`g_i`$ has order $`2`$ and it fixes the simplex pointwise. Thus the fixed point set of $`A`$ on $`NC(G)`$ is nothing other than $`NC_2(G)`$ the $`2`$-local non-commuting complex for $`G`$. Since $`A`$ has order 2, we can apply Smith Theory to finish the proof of the first statement of the proposition. The identity on the Euler characteristics follows once we note that under the action of $`A`$, the cells of $`NC(G)`$ break up into free orbits and cells which are fixed by $`A`$ and the fixed cells exactly form $`NC_2(G)`$. ∎
We observed earlier that if $`BNC(G)`$ is contractible then the center of $`G`$ is trivial, i.e. $`BNC(G)=NC(G)`$, hence, by proposition 3.6, $`NC_2(G)`$ is $`𝔽_2`$-acyclic (and in particular nonempty).
###### Definition 3.7.
Let $`G`$ be a finite group. We define $`nc(G)`$ to be the maximum size of a set of pairwise non-commuting elements in $`G`$. Thus $`nc(G)1`$ is the dimension of $`NC(G)`$.
We now compute a general class of examples, the Frobenius groups. Recall a group $`G`$ is a Frobenius group if it has a proper nontrivial subgroup $`H`$ with the property that $`HH^g=1`$ if $`gGH`$. $`H`$ is called the Frobenius complement of $`G`$. Frobenius showed the existence of a normal subgroup $`K`$ such that $`K=G_{gG}(H^g\{1\})`$. Thus $`G`$ is a split extension of $`K`$ by $`H`$, i.e., $`G=K\times _\varphi H`$, for some homomorphism $`\varphi :HAut(K)`$. This $`K`$ is called the Frobenius kernel of $`G`$.
We have:
###### Proposition 3.8.
If $`G`$ is a Frobenius group with Frobenius kernel $`K`$ and Frobenius complement $`H`$, then
$$NC(G)=NC(K)NC(H)^{|K|}$$
where $``$ stands for simplicial join and the superscript $`|K|`$ means that $`NC(K)`$ is joined repeatedly with $`|K|`$ many copies of $`NC(H)`$.
It also follows that $`nc(G)=nc(K)+|K|nc(H)`$.
Finally if both $`H`$ and $`K`$ are abelian, then $`NC(G)`$ is homotopy equivalent to a wedge of $`(|K|2)(|H|2)^{|K|}`$ many $`|K|`$-spheres.
###### Proof.
First, from the condition that $`HH^g=1`$ for $`gGH`$ we see that no nonidentity element of $`H`$ commutes with anything outside of $`H`$. Thus conjugating the picture, no nonidentity element of any $`H^g`$ commutes with anything outside $`H^g`$. Thus if $`H^{g_1},H^{g_2},\mathrm{},H^{g_m}`$ is a complete list of the conjugates of $`H`$ in $`G`$, we see easily that $`G\{1\}`$ is partitioned into the sets $`K\{1\},H^{g_1}\{1\},\mathrm{},H^{g_m}1`$ and two elements picked from different sets in this partition will not commute. Thus it follows that the non-commuting complex based on the elements of $`G\{1\}`$ decomposes as a join of the non-commuting complexes based on each set in the partition. To complete the picture one notices that each $`H^g`$ is isomorphic to $`H`$ and so contributes the same non-commuting complex as $`H`$ and furthermore since the conjugates of $`H`$ make up $`GK`$, a simple count gives that $`m=\frac{|G||K|}{|H|1}=\frac{|H||K||K|}{|H|1}=|K|`$.
The sentence about $`nc(G)`$ follows from the fact that if we define $`d(S)=\mathrm{dim}(S)+1`$ for a simplicial complex $`S`$, then $`d(S_1S_2)=d(S_1)+d(S_2)`$. Thus since $`d(NC(G))=nc(G)`$ this proves the stated formula concering $`nc(G)`$.
Finally when $`H`$ and $`K`$ are abelian, $`NC(H)`$ and $`NC(K)`$ are just a set of points namely the nonidentity elements in each group. The short exact sequence of the join together with the fact that $`NC(G)=BNC(G)`$ is simply-connected finishes the proof. ∎
###### Example 3.9.
$`A_4`$ is a Frobenius group with kernel $`/2\times /2`$ and complement $`/3`$. Thus proposition 3.8 shows $`NC(A_4)S^4S^4`$ and $`nc(A_4)=5`$.
###### Claim 3.10.
$`NC_2(A_5)`$ is a $`4`$-spherical complex homotopy equivalent to a bouquet of $`32`$ $`4`$-spheres. Thus it is $`3`$-connected and has odd Euler characteristic. Hence $`NC(A_5)`$ has odd Euler characteristic.
###### Proof.
The order of $`A_5`$ is $`60=435`$ of course. It is easy to check that there are five Sylow 2-subgroups $`P`$ which are elementary abelian of rank 2 and self-centralizing, i.e., $`C_G(P)=P`$, and are “disjoint”, i.e., any two Sylow subgroups intersect only at the identity element.
Thus the picture for the vertices of $`NC_2(A_5)`$ is as 5 sets $`\{S_i\}_{i=1}^5`$ of size 3. (Since each Sylow 2-group gives 3 involutions.) Now since the Sylow 2-groups are self-centralizing, this means that two involutions in two different Sylow 2-subgroups, do not commute and thus are joined by an edge in $`NC_2(A_5)`$. Thus we see easily that $`NC_2(A_5)`$ is the join $`S_1S_2S_3S_4S_5`$.
Using the short exact sequence for the join (see page 373, Exercise 3 in \[M\]), one calculates easily that $`NC_2(A_5)`$ has the homology of a bouquet of $`32`$ $`4`$-spheres. Since the join of two path connected spaces $`S_1S_2`$ and $`S_3S_4S_5`$ is simply-connected, it follows that $`NC_2(A_5)`$ is homotopy equivalent to a bouquet of $`32`$ $`4`$-spheres, and since it is obviously $`4`$-dimensional, this completes all but the last sentence of the claim. The final sentence follows from proposition 3.6 which says that $`NC(A_5)`$ has the same Euler characteristic as $`NC_2(A_5)`$ mod 2. ∎
At this stage, we would like to make a conjecture:
###### Conjecture 3.11.
If $`G`$ is a nonabelian simple group, then $`NC(G)=BNC(G)`$ has odd Euler characteristic.
Recall the following famous theorem of Feit and Thompson:
###### Theorem 3.12 (Odd order theorem).
Every group of odd order is solvable.
Notice, that if the conjecture is true, it would imply the odd order theorem. This is because it is easy to see that a minimal counterexample $`G`$ to the odd order theorem would have to be an odd order nonabelian simple group. The conjecture would then say $`BNC(G)`$ has an odd Euler characteristic and proposition 3.3 would say $`G`$ was even order which is a contradiction to the original assumption that $`G`$ has odd order.
## 4 General commuting structures
Before we further analyze the $`NC(G)`$ complexes introduced in the last section, we need to extend our considerations to general commuting structures.
###### Definition 4.1.
A commuting structure is a set $`S`$ together with a reflexive, symmetric relation $``$ on $`S`$. If $`x,yS`$ with $`xy`$ we say $`x`$ and $`y`$ commute.
###### Definition 4.2.
Given a commuting structure $`(S,)`$, the dual commuting structure $`(S,^{})`$ is defined by
(a) $`^{}`$ is reflexive and
(b)For $`xy`$, $`x^{}y`$ if and only if $`x`$ does not commute with $`y`$ in $`(S,)`$.
When it is understood we write a commuting structure as $`S`$ and its dual as $`S^{}`$. It is easy to see that $`S^{\prime \prime }=S`$ in general.
###### Definition 4.3.
Given a commuting structure $`S`$, $`C(S)`$ is the simplicial complex whose vertices are the elements of $`S`$ and such that $`[s_0,\mathrm{},s_n]`$ is a face of $`C(S)`$ if and only if $`\{s_0,\mathrm{},s_n\}`$ is a commuting set i.e., $`s_is_j`$ for all $`i`$ and $`j`$. Similarly we define $`NC(S)=C(S^{})`$ and refer to the elements in a face of $`NC(S)`$ as a non-commuting set in $`S`$.
Below are some examples of commuting structures which will be important in our considerations:
(a) The nontrivial elements of a group $`G`$ form a commuting structure which we denote also by $`G`$, where $`xy`$ if and only if $`x`$ and $`y`$ commute in the group $`G`$. In this case $`C(G)`$ and $`NC(G)`$ are the complexes considered in the previous sections.
(b) The noncentral elements of a group $`G`$ form a commuting structure $`GZ(G)`$ and $`NC(GZ(G))=BNC(G)`$. Similarly the elements of order $`p`$ in $`G`$ form a commuting structure $`G_p`$ and $`C(G_p)=C_p(G)`$, $`NC(G_p)=NC_p(G)`$.
(c) If $`V`$ is a vector space equipped with a bilinear map $`[,]:VVV`$. Then the nonzero elements of $`V`$ form a commuting structure also denoted by $`V`$, where $`v_1v_2`$ if $`[v_1,v_2]=0`$ or if $`v_1=v_2`$.
(d) In the situation in (c), we can also look at the set $`P(V)`$ of lines in $`V`$. The commuting structure on $`V`$ descends to give a well-defined commuting structure on $`P(V)`$, which we will call the projective commuting structure.
(e) If $`1CGQ1`$ is a central extension of groups with $`Q`$ abelian then one can form $`[,]:QC`$ by
$$[x,y]=\widehat{x}\widehat{y}\widehat{x}^1\widehat{y}^1$$
where $`x,yQ`$ and $`\widehat{x}`$ is a lift of $`x`$ in $`G`$ etc. It is easy to see this is well-defined, independent of the choice of lift and furthermore that the bracket $`[,]`$ is bilinear. Thus by (c), we get a commuting structure on the nontrivial elements of $`Q`$ via this bracket. We denote this commuting structure by $`(G;C)`$. Notice in general it is not the same as the commuting structure of the group $`G/C`$ which is abelian in this case. More generally even if $`Q`$ is not abelian, one can define a commuting structure on the nontrivial elements of $`Q`$ from the extension by declaring $`xy`$ if and only if $`\widehat{x}`$ and $`\widehat{y}`$ commute in $`G`$, we will denote this commuting structure by $`(G;C)`$ in general.
###### Example 4.4.
If $`P`$ is an extraspecial $`p`$-group of order $`p^3`$ then it has center $`Z`$ of order $`p`$ and $`P/Z`$ is an elementary abelian $`p`$-group of rank 2. Let $`Symp`$ denote the commuting structure obtained from a vector space of dimension two over $`𝔽_p`$ equipped with the symplectic alternating inner product $`[x,y]=1`$ where $`\{x,y\}`$ is a suitable basis. Then it is easy to check that $`(P;Z)=Symp`$.
One of the main results we will use in order to study the non-commuting complexes associated to commuting structures is a result of A. Björner, M. Wachs and V. Welker on “blowup” complexes which we describe next:
Let $`S`$ be a finite simplicial complex with vertex set $`[n]`$ (This means the vertices have been labelled 1,…,n). To each vertex $`1in`$, we assign a positive integer $`m_i`$. Let $`\overline{m}=(m_1,\mathrm{},m_n)`$, then the “blowup” complex $`S_{\overline{m}}`$ is defined as follows:
The vertices of $`S_{\overline{m}}`$ are of the form $`(i,j)`$ where $`1in,1jm_i`$. (one should picture $`m_i`$ vertices in $`S_{\overline{m}}`$ over vertex $`i`$ in $`S`$.)
The faces of $`S_{\overline{m}}`$ are exactly of the form $`[(i_0,j_0),\mathrm{},(i_n,j_n)]`$ where $`[i_0,\mathrm{},i_n]`$ is a face in $`S`$. (In particular $`i_li_k`$ for $`lk`$.)
The result of Björner, Wachs and Welker describes $`S_{\overline{m}}`$ up to homotopy equivalence, in terms of $`S`$ and its links. More precisely we have:
###### Theorem 4.5 (Björner, Wachs, Welker \[BWW\]).
For any connected simplicial complex $`S`$ with vertex set $`[n]`$ and given $`n`$-tuple of positive integers $`\overline{m}=(m_1,\mathrm{},m_n)`$ we have:
$$S_{\overline{m}}S\underset{FS}{}[Susp^{|F|}(Lk(F))]^{\gamma (F)}.$$
Here the $``$ stands for wedge of spaces, $`Susp^k`$ for $`k`$-fold suspension, $`Lk(F)`$ for the link of the face $`F`$ in $`S`$, $`|F|`$ for the number of vertices in the face $`F`$ (which is one more than the dimension of $`F`$) and $`\gamma (F)=_{iF}(m_i1)`$. Thus in the decomposition above, $`\gamma (F)`$ copies of $`Susp^{|F|}(Lk(F))`$ appear wedged together.
Note that in \[BWW\], the empty face is considered a face in any complex. In the above formulation we are not considering the empty face as a face and thus have separated out the $`S`$ term in the wedge decomposition.
For example, if we take a 1-simplex $`[1,2]`$ as our complex $`S`$ and use $`\overline{m}=(2,2)`$, then it is easy to see that $`S_{\overline{m}}`$ is a circle. On the other hand all links in $`S`$ are contractible except the link of the maximal face $`[1,2]`$ which is empty (a “$`(1)`$-sphere”). Thus everything in the right hand side of the formula is contractible except the 2-fold suspension of this $`(1)`$-sphere which gives a 1-sphere or circle as expected.
Also note that whenever some $`m_i=1`$, the corresponding $`\gamma (F)=0`$ and so that term drops our of the wedge decomposition. Thus if $`m_1=\mathrm{}=m_n=1`$, the decomposition gives us nothing as $`S_{\overline{m}}=S`$.
Now for some examples of where this theorem applies:
###### Corollary 4.6.
If $`G`$ is a finite group and $`Z(G)`$ is its center, then $`BNC(G)`$ is the blowup of $`(G;Z(G))`$ where each $`m_i=|Z(G)|`$. This is because everything in the same coset of the $`Z(G)`$ in $`G`$ commutes with each other and whether or not two elements from different cosets commute is decided in $`(G;Z(G))`$. Thus we conclude
$$BNC(G)NC(G;Z(G))\underset{FNC(G;Z(G))}{}[Susp^{|F|}(Lk(F))]^{(|Z(G)|1)^{|F|}}.$$
Before we do say more, let us look at another example of theorem 4.5 in our context. The proof is the same as that of example 4.6 and is left to the reader.
###### Corollary 4.7.
Let $`(V,[,])`$ be a vector space over a finite field $`𝐤`$ equipped with a nondegenerate bilinear form. Then $`NC(V)`$ is a blowup of $`NC(P(V))`$ and we have
$$NC(V)NC(P(V))\underset{FNC(P(V))}{}[Susp^{|F|}(Lk(F))]^{(|𝐤|2)^{|F|}}.$$
For example, it is easy to see that in $`Symp`$, if $`[x,y]=0`$ then $`x`$ is a scalar multiple of $`y`$. Thus in $`P(Symp)`$ two distinct elements do not commute and hence $`NC(P(Symp))`$ is a simplex. Thus all of its links are contractible except the link of the top face which is empty. This top face has $`(p^21)/(p1)=p+1`$ vertices. Thus following example 4.7, we see
$$NC(Symp)_{(p2)^{(p+1)}}S^p.$$
###### Remark 4.8.
If $`(V,[,])`$ is a vector space over a finite field equipped with a nondegenerate symmetric inner product corresponding to a quadratic form $`Q`$, then one can restrict the commuting structure induced on $`P(V)`$ to the subset $`S`$ consisting of singular points of $`Q`$, i.e., lines $`<x>`$ with $`Q(x)=0`$. ($`S`$ will be nonempty only if $`Q`$ is of hyperbolic type.) Among the non-commuting sets in $`S`$ are the subsets called ovoids defined by the property that every maximal singular subspace of $`V`$ contains exactly one element of the ovoid. (See \[G\]) Thus these ovoids are a special subcollection of the facets of $`NC(S)`$.
In general, we will find the following notion useful:
###### Definition 4.9.
If $`(S,)`$ is a commuting set and $`xS`$, we define the centralizer of $`x`$ to be
$$C(x)=\{yS|yx\}.$$
This allows us to define
###### Definition 4.10.
If $`(S,)`$ is a commuting set, we define an equivalence relation on the elements of $`S`$ by $`xy`$ if $`C(x)=C(y)`$. The equivalence classes are called the centralizer classes of $`(S,)`$. We then define the core of $`S`$ to be the commuting set $`(\overline{S},)`$ where the elements are the centralizer classes of $`S`$ and the classes $`[x]`$ and $`[y]`$ commute in $`\overline{S}`$ if and only if the representatives $`x`$, $`y`$ commute in $`S`$.
It is easy to see that $`NC(S)`$ is the blowup of $`NC(\overline{S})`$, where for the vertex $`[x]NC(\overline{S})`$ there are $`m_x`$ vertices above it in $`NC(S)`$, where $`m_x`$ is the size of the centralizer class $`[x]`$. Thus once again theorem 4.5 gives us:
###### Theorem 4.11.
Let $`S`$ be a commuting set and $`\overline{S}`$ be its core, and suppose $`NC(\overline{S})`$ is connected. Then
$$NC(S)NC(\overline{S})\underset{FNC(\overline{S})}{}[Susp^{|F|}(LkF)]^{\gamma (F)}$$
where $`\gamma (F)=_{[x]F}(m_x1)`$. Here again $`m_x`$ is the size of the centralizer class $`[x]`$.
###### Remark 4.12.
Note that if $`NC(S)`$ is connected, $`NC(\overline{S})`$ will automatically be connected as it is the image of $`NC(S)`$ under a continuous map.
With some abuse of notation, we will call $`NC(\overline{S})`$ the core of $`NC(S)`$.
Following the notation above, in a finite group $`G`$, the equivalence relation “has the same centralizer group” partitions $`G`$ into centralizer classes. The central elements form one class and the noncentral elements thus inherit a partition.
Notice if $`[x]`$ is the centralizer class containing $`x`$, then any other generator of the cyclic group $`<x>`$ is in the same class. Thus $`[x]`$ has at least $`\varphi (n)`$ elements where $`n`$ is the order of $`x`$, and $`\varphi `$ is Euler’s totient function. Thus if $`n>2`$, then $`[x]`$ contains at least two elements. Also notice that every thing in the coset $`xZ(G)`$ is also in $`[x]`$ so if $`Z(G)1`$ we can also conclude $`[x]`$ contains at least two elements.
###### Definition 4.13.
A non-commuting set $`S`$ in $`G`$ is a nonempty subset $`S`$, such that the elements of $`S`$ pairwise do not commute. A maximal non-commuting set $`S`$ is a non-commuting set which is not properly contained in any other non-commuting set of $`G`$.
In general, not all maximal non-commuting sets of a group $`G`$ have the same size.
One obtains the following immediate corollary of theorem 4.11 (Assume $`|G|>2`$ for the following results):
###### Corollary 4.14.
Let $`G`$ be a finite nonabelian group, and let $`S_s`$ denote the set of maximal non-commuting sets in $`G`$ of size $`s`$. Then, for $`s>1`$,
$$rk(H_{s1}(BNC(G)))\underset{FS_s}{}[\underset{xF}{}(1\frac{1}{m_x})].$$
where $`m_x`$ is the size of the centralizer class containing $`x`$.
In particular, if $`G`$ is an odd order group or if $`G`$ has a nontrivial center, then
$$\stackrel{~}{H}_{s1}(NC(G))0$$
whenever $`G`$ has a maximal non-commuting set of size $`s`$.
###### Proof.
Let $`X`$ be the non-commuting complex associated to the core of $`BNC(G)`$. Let us define a facet to be a face of a simplicial complex which is not contained in any bigger faces. Thus the facets of $`BNC(G)`$ consist exactly of maximal non-commuting sets of noncentral elements in $`G`$.
The first thing to notice is to every facet $`F`$ of $`BNC(G)`$, there corresponds a facet $`\overline{F}`$ of $`X`$, and further more this correspondence preserves the dimension of the facet (or equivalently the number of vertices in the facet).
Since the link of a facet is always empty (a $`(1)`$-sphere), in the wedge decomposition of theorem 4.11, we get the suspension of this empty link as a contribution. If the facet $`F`$ has $`n`$ vertices in it, then we suspend $`n`$ times to get a $`(n1)`$-sphere. Thus to every maximal non-commuting set of size $`s`$, we get a $`s1`$ sphere contribution from the corresponding facet in $`X`$. In fact we get $`\gamma (F)`$-many such spheres from the facet $`F`$. However above each facet $`\overline{F}`$ in $`X`$, there correspond $`_{[x]\overline{F}}(m_x)`$ many facets in $`BNC(G)`$. Thus in the sum over the facets of $`BNC(G)`$ stated in the theorem, we divide $`\gamma (F)`$ by $`_{xF}(m_x)`$ in order to count the contribution from the facet $`\overline{F}`$ in $`X`$ the correct number of times.
Notice that $`\frac{\gamma (F)}{_{xF}(m_x)}=_{xF}(1\frac{1}{m_x})\frac{1}{2^{|F|}}`$ if all the centralizer classes have size bigger than one. So if $`G`$ has the property that the size of the centralizer classes of noncentral elements is always strictly bigger than one, for example if $`G`$ is odd or if $`G`$ has a nontrivial center, then
$$2^srk(H_{s1}(BNC(G)))|S_s|$$
for all $`s`$. In particular $`H_{s1}(BNC(G))0`$ whenever $`G`$ has a maximal non-commuting set of size $`s`$. Also observe that $`\stackrel{~}{H}_0(NC(G))0`$ whenever $`G`$ has maximal non-commuting set of size $`1`$, i.e. a singleton consisting of a nontrivial central element (except for the trivial case when $`G`$ has order 2.) ∎
It is easy to see from this proof that a $`p`$-local version of corollary 4.14 is also true. We state here only the last part of this result for odd primes:
###### Corollary 4.15.
Let $`p`$ be an odd prime. Then,
$$\stackrel{~}{H}_{s1}(NC_p(G))0$$
whenever $`G`$ has a maximal non-commuting $`p`$-set of size $`s`$.
The same is true for $`p=2`$ under the additional condition $`2||Z(G)|`$.
###### Remark 4.16.
Notice that the conclusion of corollary 4.14 is consistent with the simple connectedness of $`BNC(G)`$, because there is no maximal non-commuting set of size $`2`$. To see this, observe that whenever there is a non-commuting set $`\{a,b\}`$ with two elements, we can form a bigger non-commuting set $`\{a,b,ab\}`$.
Also notice that this is no longer the case for the $`p`$-local case. For example, when $`G=D_8=a,b|a^2=b^2=c^2=1,[a,b]=c,ccentral`$, the complex $`BNC_2(G)`$ is a rectangle with vertices $`a,b,ac,bc`$ which is the inflated complex corresponding to maximal $`2`$-set $`\{a,b\}`$. In particular, $`BNC_2(G)`$ is not simply connected in general.
###### Remark 4.17.
One of the things that corollary 4.14 says is that if one wants to calculate $`nc(G)`$, the answer which is obviously the dimension of $`BNC(G)`$ plus one can also be determined by finding the highest nonvanishing homology of $`BNC(G)`$ in the case when $`Z(G)1`$ or $`G`$ is odd order. Thus the answer is determined already by the homotopy type of $`BNC(G)`$ in this situation. If $`Z(G)=1`$ and $`G`$ has even order, this is no longer true, for example $`NC(\mathrm{\Sigma }_3)`$ is contractible and so does not have any positive dimensional homology.
Sometimes one can show the non-commuting complex for the core of $`(S,)`$ is contractible as the next lemma shows:
###### Lemma 4.18.
If $`(S,)`$ has a centralizer class $`[x]`$ where $`[x]=C(x)`$ then if the core is $`\overline{S}`$, then $`NC(\overline{S})`$ is contractible.
###### Proof.
This is because it is easily seen that $`NC(\overline{S})`$ is a cone on the vertex $`[x]`$ as everything outside $`[x]`$ does not commute with $`x`$ as $`[x]=C(x)`$. ∎
Thus for example:
###### Example 4.19.
Let $`p`$ be a prime, then the core of $`BNC(\mathrm{\Sigma }_p)`$ is contractible.
###### Proof.
The cycle $`x=(1,2,\mathrm{},p)`$ in $`\mathrm{\Sigma }_p`$ has $`C(x)=<x>`$ by a simple calculation. Thus $`C(x)=[x]`$ and so the result follows from lemma 4.18. ∎
We now study an important special case:
###### Definition 4.20.
A $`TC`$-group $`G`$ is a group where the commuting relation on the noncentral elements is transitive. This is equivalent to the condition that all proper centralizer subgroups $`C(x)G`$ are abelian.
For example any minimal nonabelian group like $`S_3`$, $`A_4`$ or an extraspecial group of order $`p^3`$.
###### Corollary 4.21.
If $`(S,)`$ is a commuting set where $``$ is also transitive (i.e., $``$ is an equivalence relation), then $`NC(S)`$ is homotopy equivalent to a wedge of spheres of the same dimension. (We allow the “empty” wedge, i.e., we allow the case $`NC(S)`$ is homotopy equivalent to a point). The dimension of the spheres is equal to $`n1`$ where $`n`$ is the number of equivalence classes in $`(S,)`$ and the number of spheres appearing is $`_{i=1}^n(m_i1)`$ where $`m_i`$ is the size of equivalence class $`i`$.
Thus if $`G`$ is a $`TC`$-group, then $`BNC(G)`$ is homotopy equivalent to a wedge of spheres of dimension $`nc(G)1`$ and the number of spheres is given by a product as above where the $`m_i`$ are the orders of the distinct proper centralizer groups of $`G`$.
###### Proof.
If $``$ is an equivalence relation, it is easy to see that the centralizer classes are exactly the $``$ equivalence classes. Thus one sees that the non-commuting complex for the core, $`NC(\overline{S})`$, is a simplex. Thus in Theorem 4.11, all terms drop out except that corresponding to the maximum face in $`NC(\overline{S})`$ where the link is empty. This link is suspended to give a sphere of dimension equal to the number of equivalence classes minus one. The number of these spheres appearing in the wedge decomposition is the product $`_{i=1}^n(m_i1)`$ where $`m_i`$ is the size of equivalence class $`i`$ and the product is over all equivalence classes. (Thus this can be zero if one of the equivalence classes has size one, in which case the complex is contractible).
In the case of $`BNC(G)`$, for $`G`$ a $`TC`$-group, one just has to note that $`C(x)`$ is the centralizer class $`[x]`$ for any noncentral element $`x`$. ∎
###### Remark 4.22.
In Corollary 4.21, one did not actually have to use the general result of Bjorner, Wachs and Welker since it is easy to see that in this situation, $`NC(S)`$ is the join of each equivalence class as discrete sets and a simple count gives the result.
Now notice in the case that $`(S,)`$ has $``$ transitive, the above analysis shows that $`NC(S)`$ is a join of discrete sets. Thus it is easy to see that $`NC(S)`$ is shellable (This is because the facets of $`NC(S)`$ are just sets where we have chosen exactly one element from each of the $``$ equivalence classes. We can linearly order the equivalence classes and then lexicographically order the facets. It is easy to check that this is indeed a shelling.)
Given a shelling of a simplicial complex, there are many combinatorial equalities and inequalities which follow (See \[S\].) Since these are not so deep in the above general context, we will point out only the interpretation when applied to $`BNC(G)`$. Recall pure shellable just means shellable where all the facets have the same dimension.
###### Proposition 4.23.
If $`G`$ is a nonabelian group such that $`BNC(G)`$ is pure shellable, e.g., $`G`$ a $`TC`$-group like $`\mathrm{\Sigma }_3`$ or $`A_4`$, then if one sets $`nc_i`$ to be the number of non-commutative sets of non-central elements which have size $`i`$, one has:
$$(1)^jC(nc(G),j)+\underset{k=1}{\overset{j}{}}(1)^{jk}C(nc(G)k,jk)nc_k0$$
for all $`1jnc(G)`$ where $`C(n,k)`$ is the usual binomial coefficient.
###### Proof.
Follows from a direct interpretation of the inequalities in \[S\], page 4, Theorem 2.9. One warning about the notation in that paper is that $`|\sigma |`$ means the number of vertices in $`\sigma `$ and the empty face is considered a simplex in any complex. ∎
Using that $`nc_1=|G||Z(G)|`$, $`nc_2=\frac{|G|}{2}(|G|m)`$ where $`m`$ is the number of conjugacy classes in $`G`$, one gets for example from the inequality with $`j=2`$ above:
###### Corollary 4.24.
Let $`G`$ be a nonabelian group with a transitive commuting relation, i.e., if $`[g,h]=[h,k]=1`$, then $`[g,h]=1`$ for every noncentral $`g,h,kG`$. Then,
$$nc(G)(nc(G)1)+|G|(|G|m)2(nc(G)1)(|G||Z(G)|)0.$$
where $`m`$ denotes the number of conjugacy classes in $`G`$.
## 5 Duality
Let $`(S,)`$ be a finite set with a commuting relation. Suppose the commuting complex $`C(S)`$ breaks up as a disjoint union of path components $`C(S_1),\mathrm{},C(S_n)`$ where of course we are using $`S_i`$ to stand for the vertex set of component $`i`$.
Then notice in the corresponding non-commuting complex, $`NC(S)`$ we have $`NC(S)=NC(S_1)\mathrm{}NC(S_n)`$ where $``$ stands for the join operation as usual.
We state this simple but useful observation as the next lemma:
###### Lemma 5.1 (Duality).
Let $`(S,)`$ be a commuting set then if
$$C(S)=\underset{i=1}{\overset{n}{}}C(S_i)$$
where $``$ stands for disjoint union, then we have
$$NC(S)=_{i=1}^nNC(S_i)$$
where $``$ stands for join.
Thus in some sense “the less connected $`C(S)`$ is, the more connected $`NC(S)`$ is.”
We can apply this simple observation to say something about the complexes $`NC_p(G)`$ in general.
###### Theorem 5.2.
Let $`G`$ be a finite group and $`p`$ a prime such that $`p||G|`$. Let $`P`$ be a Sylow $`p`$-group of $`G`$ and define $`N`$ to be the subgroup generated by $`N_G(H)`$ as $`H`$ runs over all the nontrivial subgroups of $`P`$.
Then $`NC_p(G)`$ is $`(|G:N|2)`$-connected and in fact it is the $`|G:N|`$-fold join of some complex with itself.
###### Proof.
By Quillen \[Q\], if $`S_1,\mathrm{},S_n`$ are the components of $`A_p(G)`$, then under the $`G`$-action, $`G`$ acts transitively on the components with isotropy group $`N`$ under suitable choice of labelling. Thus the components are all simplicially equivalent and there are $`|G:N|`$ many of them.
However we have seen that $`A_p(G)`$ is $`G`$-homotopy equivalent to $`C_p(G)`$ and so we have the same picture for that complex. Thus $`C_p(G)`$ is the disjoint union of $`|G:N|`$ copies of some simplicial complex $`S`$. Thus by lemma 5.1, $`NC_p(G)`$ is the $`|G:N|`$-fold join of the dual of $`S`$ with itself. To finish the proof one just has to note that a $`k`$-fold join of nonempty spaces is always $`(k2)`$-connected. ∎
Dept. of Mathematics
University of Wisconsin-Madison,
Madison, WI 53706, U.S.A.
E-mail address: pakianat@math.wisc.edu
Dept. of Mathematics & Statistics
McMaster University
Hamilton, ON, Canada
L8S 4K1.
E-mail address: yalcine@math.mcmaster.ca
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# References
Recently the problem of bosonizing the pairing hamiltonian $`H_P`$ has been addressed in the formalism of the path integrals over Grassmann variables . In this context, using even elements of the Grassman algebra, a functional integral representation of the correlation functions has been given by performing a change of variables in the Berezin integral. It has thus been possible to express the pairing action and the ground state wave function through “collective variables”, namely specific linear combinations of the even elements of the algebra, and to obtain the ground state energy of $`H_P`$.
However, to solve the problem of the excited states of $`H_P`$ (with non zero seniority $`v`$) in the path integral framework has proved to be a quite difficult task to perform. Accordingly here we address the $`v0`$ problem resorting to the Grassmann algebra in the hamiltonian framework.
Of course the spectrum of $`H_P`$ has already been found long time ago within the quasi spin formalism , but still we believe it useful to address the same problem in this new context to shed light on the role of composite fields in nuclear physics and to provide explicit expressions for the eigenfunctions of $`H_P`$ with a non vanishing seniority.
To start with, we shortly recall the essential elements of the hamiltonian many-fermions problem in terms of Grassmann variables (the bosonic problem, likewise, can be formulated in terms of holomorphic variables) .
One exploits the isomorphism between the Fock space $`F`$ generated by the fermionic creation operators $`\widehat{a}_1^+,\mathrm{}\widehat{a}_{2\mathrm{\Omega }}^+`$ and the algebra $`𝒢^+`$ generated by the set of totally anticommuting objects $`\lambda _1^{},\mathrm{}\lambda _{2\mathrm{\Omega }}^{}`$ (odd elements of the algebra). The isomorphism is defined by mapping the vectors $`\widehat{a}_1^+\mathrm{}\widehat{a}_j^+|0>`$ onto the elements $`\lambda _1^{}\mathrm{}\lambda _j^{}`$. The image of a generic vector $`|\mathrm{\Psi }>F`$ under this mapping will be denoted by $`\mathrm{\Psi }(\lambda ^{})`$.
The scalar product in $`F`$ in terms of Grassmann variables reads then
$$<\mathrm{\Psi }_1|\mathrm{\Psi }_2>=𝑑\lambda _{2\mathrm{\Omega }}^{}𝑑\lambda _{2\mathrm{\Omega }}\mathrm{}𝑑\lambda _1^{}𝑑\lambda _1(\mathrm{\Psi }_1(\lambda ^{}))^{}\mu _{}(\lambda ^{}\lambda )\mathrm{\Psi }_2(\lambda ^{}),$$
(1)
the measure being
$$\mu _\pm (\lambda ^{}\lambda )=e^{\pm _i\lambda _i^{}\lambda _i}.$$
(2)
Actually to define the kernels of integral operators the measure with the plus sign is needed.
Indeed a linear operator in normal form, namely
$$\widehat{𝒪}=\underset{i_1\mathrm{}i_k}{}\underset{j_1\mathrm{}j_k}{}𝒪^{i_1\mathrm{}i_k,j_1\mathrm{}j_k}\widehat{a}_{i_1}^+\mathrm{}\widehat{a}_{i_k}^+\widehat{a}_{j_1}\mathrm{}\widehat{a}_{j_k},$$
(3)
can be expressed in terms of Grassmann variables according to
$$𝒪(\lambda ^{},\lambda )=\underset{i_1,\mathrm{}i_k,j_1\mathrm{}j_k}{}𝒪^{i_1,\mathrm{}i_k,j_1\mathrm{}j_k}\lambda _{i_1}^{}\mathrm{}\lambda _{i_k}^{}\lambda _{j_1}\mathrm{}\lambda _{j_k},$$
(4)
and its associated kernel reads
$$K_𝒪(\lambda ^{},\lambda )=𝒪(\lambda ^{},\lambda )\mu _+(\lambda ^{}\lambda ).$$
(5)
Then the action of $`𝒪`$ on a state $`\mathrm{\Psi }`$ is given, in integral form, as
$$(𝒪\mathrm{\Psi })(\lambda ^{})=[d\lambda ^{}d\lambda ^{}]K_𝒪(\lambda ^{},\lambda ^{})\mu _{}(\lambda ^{}\lambda ^{})\mathrm{\Psi }(\lambda ^{}).$$
(6)
We apply the formalism, shortly revisited above, to a system made of one kind of nucleons (e.g. neutrons) in a level specified by the angular momentum $`j`$, hence the indices labelling the $`\lambda `$ variables will correspond to the third components $`m`$ of $`j`$.
Now if only pairs of nucleons coupled to a total angular momentum with $`J_z=0`$ are considered, then the exponentials appearing in $`\mu _\pm `$ can be transformed as follows
$$[d\lambda ^{}]\mathrm{exp}\left(\pm \lambda _m^{}\lambda _m\pm \lambda _m^{}\lambda _m\right)f(\varphi ^{})=[d\varphi ^{}]\mathrm{exp}\left(\varphi _m^{}\varphi _m\right)f(\varphi ^{}),$$
(7)
in terms of the even Grassmann variables
$$\varphi _m=\lambda _m\lambda _m.$$
(8)
For $`N`$ (even) nucleons in a single particle level $`j`$, the pairing hamiltonian reads
$$\widehat{H}_P^{(j)}=g_P\widehat{A}_j^{}\widehat{A}_j,$$
(9)
where
$$\widehat{A}_j=\sqrt{\frac{\mathrm{\Omega }}{2}}\underset{m=j}{\overset{j}{}}jm,jm|00\widehat{a}_{jm}\widehat{a}_{jm}=\underset{m=1/2}{\overset{j}{}}(1)^{jm}\widehat{a}_{jm}\widehat{a}_{jm},$$
(10)
being $`\mathrm{\Omega }=(2j+1)/2`$ and $`jm,jm|00`$ the usual Clebsch Gordan coefficient. In the frame of the Grassmann algebra, we recast instead Eq.(9) as follows
$$H_P^{(j)}=g_PA_j^{}A_j$$
(11)
with
$$A_j=\underset{m=1/2}{\overset{j}{}}(1)^{jm}\lambda _m\lambda _m=\underset{m=1/2}{\overset{j}{}}(1)^{jm}\varphi _m.$$
(12)
We explore now whether the $`\varphi `$’s can be the building blocks of the eigenstates of $`H_P^{(j)}`$. Indeed this turns out to be the case, since, according to (6) and (7), the action of the hamiltonian over states set up with the $`\varphi `$’s is
$$H_P^{(j)}\psi (\varphi ^{})=[d\eta ^{}d\eta ]K_P^{(j)}(\varphi ^{},\eta )e^{{\scriptscriptstyle \eta ^{}\eta }}\psi (\eta ^{}),$$
(13)
the kernel reading
$$K_P^{(j)}(\varphi ^{},\eta )=H_P^{(j)}(\varphi ^{},\eta )e^{{\scriptscriptstyle \varphi ^{}\eta }}.$$
(14)
Hence the eigenvalues equation
$$H_P^{(j)}\psi (\varphi ^{})=[d\eta ^{}d\eta ]K_P^{(j)}(\varphi ^{},\eta )e^{{\scriptscriptstyle \eta ^{}\eta }}\psi (\eta ^{})=E^{(j)}\psi (\varphi ^{})$$
(15)
follows.
We then look for eigenstates of $`H_P^{(j)}`$ in the form of linear combinations of products of $`\varphi `$’s. More specifically, for fixed $`\mathrm{\Omega }`$ and for $`n=N/2`$ pairs, the wave function will be a superposition of $`\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{n}\right)`$ monomials, of the type
$$\psi (\varphi ^{})=\underset{i=1}{\overset{\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{n}\right)}{}}\beta _i[\varphi _{m_1}^{}\mathrm{}\varphi _{m_n}^{}]_i,$$
(16)
where the $`\beta _i`$ are complex coefficients.
In order to study the secular equation of $`H_P^{(j)}`$, we construct the matrix associated to the latter in the basis of the states (16). To start with, we consider one pair only. In this case, the wave function and the eigenvalue equation read, respectively,
$$\psi (\varphi ^{})=\underset{m=1/2}{\overset{j}{}}(1)^{jm}\beta _m\varphi _m^{}$$
(17)
and
$$(+1)\beta _m+\underset{p(m)=1/2}{\overset{j}{}}\beta _p=0$$
(18)
where $`=E/g_P`$ . Moreover the representation of $`H_P^{(j)}`$ in the basis (17) is
$$\left(\begin{array}{cccccc}+1& 1& 1& \mathrm{}& \mathrm{}& 1\\ 1& +1& 1& \mathrm{}& \mathrm{}& 1\\ 1& 1& +1& \mathrm{}& \mathrm{}& 1\\ 1& 1& \mathrm{}& +1& 1& 1\\ 1& 1& \mathrm{}& \mathrm{}& +1& 1\\ 1& 1& 1& \mathrm{}& \mathrm{}& +1\end{array}\right),$$
(19)
a matrix of dimension $`\mathrm{\Omega }`$ with the index $`m`$ labelling the rows and the columns. This matrix is obviously invariant for a permutation of the values of $`m`$. Elementary methods lead then to the characteristic equation
$$\left(+\mathrm{\Omega }\right)^{\mathrm{\Omega }1}=0.$$
We thus see that, of the $`\mathrm{\Omega }`$ expected real roots, two only are distinct, namely the lower one
$$_0(n=1)=\mathrm{\Omega }$$
with multiplicity $`\delta =1`$ and
$$_2(n=1)=0$$
with multiplicity $`\delta =\mathrm{\Omega }1`$.
In turn, the associated eigenvectors read, respectively,
$$\mathrm{\Psi }_0(n=1)=\frac{1}{\sqrt{\mathrm{\Omega }}}\underset{m=1/2}{\overset{j}{}}(1)^{jm}\varphi _m^{}$$
(20)
and
$$\mathrm{\Psi }_2(n=1)=𝒩\underset{m=1/2}{\overset{j1}{}}\beta _m\left[(1)^{jm}\varphi _m^{}\varphi _j^{}\right],$$
(21)
$`𝒩`$ being a normalization factor. The above findings agree with the physics of the hamiltonian $`H_P^{(j)}`$, which acts only in presence of $`J=0`$ pairs. Also the eigenvectors (20, 21) are orthogonal and the eigenvalues are those obtained in the usual quasi-spin framework.
Of significance is that the characteristic equation for the case with $`n=\mathrm{\Omega }1`$, namely
$$(+\mathrm{\Omega }1)\beta _m+\underset{p(m)=1/2}{\overset{j}{}}\beta _p=0,$$
(22)
has the same structure of (18). However the states read now
$$\psi (\varphi ^{})=\underset{m=1/2}{\overset{j}{}}(1)^{jm}\beta _m\varphi _{1/2}^{}\mathrm{}\varphi _{m1}^{}\varphi _{m+1}^{}\mathrm{}\varphi _j^{},$$
the index $`m`$ corresponding to the missing pair.
By extension, one finds that the case with $`n`$ pairs present is equivalent to the one with $`\mathrm{\Omega }n`$ pairs: indeed it leads to the same matrix but for the diagonal elements. Hence in the following we shall confine ourselves to explore the cases with $`n\frac{\mathrm{\Omega }}{2}`$ only.
Next we address the case $`n=2`$. Here the wave function reads
$$\psi (\varphi ^{})=\underset{m=1/2}{\overset{j}{}}\underset{n=1/2}{\overset{m1}{}}(1)^{jm}(1)^{jn}\beta _{mn}\varphi _m^{}\varphi _n^{},$$
(23)
entailing the secular equation
$$(+2)\beta _{mn}+\underset{p(n,m)=1/2}{\overset{j}{}}\left(\beta _{pm}+\beta _{pn}\right)=0.$$
(24)
As for the case $`n=1`$, the associated matrix can be written in $`\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{2}\right)!`$ equivalent ways (not all of them being distinct), each one corresponding to a different labelling of the states, but yielding the same determinant.
For example, for $`\mathrm{\Omega }=4`$ (the smallest $`\mathrm{\Omega }`$ for $`n=2`$) one finds the following $`6\times 6`$ matrix
$$\left(\begin{array}{cccccc}+2& 1& 1& 1& 1& 0\\ 1& +2& 1& 1& 0& 1\\ 1& 1& +2& 0& 1& 1\\ 1& 1& 0& +2& 1& 1\\ 1& 0& 1& 1& +2& 1\\ 0& 1& 1& 1& 1& +2\end{array}\right),$$
(25)
with the characteristic equation
$$(+6)(+2)^3^2=0.$$
Thus, in this case, out of 6 real roots only 3 turn out to be distinct, namely the lowest
$$_0(n=2)=6$$
with degeneracy $`\delta =1`$, the intermediate
$$_2(n=2)=2$$
with degeneracy $`\delta =3`$ and the highest
$$_4(n=2)=0$$
with degeneracy $`\delta =2`$. The associated orthogonal eigenvectors are found to be
$`\mathrm{\Psi }_0(n=2)=2𝒩_0(\varphi _{1/2}^{}\varphi _{3/2}^{}+\varphi _{1/2}^{}\varphi _{5/2}^{}\varphi _{1/2}^{}\varphi _{7/2}^{}`$
$`\varphi _{3/2}^{}\varphi _{5/2}^{}+\varphi _{3/2}^{}\varphi _{7/2}^{}\varphi _{5/2}^{}\varphi _{7/2}^{})`$ (26)
$`\mathrm{\Psi }_2(n=2)=𝒩_2[a_1(\varphi _{1/2}^{}\varphi _{3/2}^{}\varphi _{5/2}^{}\varphi _{7/2}^{})`$
$`+a_2\left(\varphi _{1/2}^{}\varphi _{5/2}^{}\varphi _{3/2}^{}\varphi _{7/2}^{}\right)`$
$`+a_3(\varphi _{1/2}^{}\varphi _{7/2}^{}\varphi _{3/2}^{}\varphi _{5/2}^{})]`$ (27)
$`\mathrm{\Psi }_4(n=2)=𝒩_4[(b_2b_3)(\varphi _{1/2}^{}\varphi _{3/2}^{}+\varphi _{5/2}^{}\varphi _{7/2}^{})`$
$`+b_2\left(\varphi _{1/2}^{}\varphi _{5/2}^{}+\varphi _{3/2}^{}\varphi _{7/2}^{}\right)`$
$`+b_3(\varphi _{1/2}^{}\varphi _{7/2}^{}+\varphi _{3/2}^{}\varphi _{5/2}^{})],`$ (28)
$`𝒩_0`$, $`𝒩_2`$, $`𝒩_4`$ being normalization factors and $`a_1,a_2,a_3,b_2`$ and $`b_3`$ free parameters.
Of significance is the correspondence between the number of parameters in the wave function (in eq. (26) the parameter has been absorbed in the normalization factor) and the roots degeneracy, namely the number of independent eigenfunctions of given energy, seniority and zero third component of the angular momentum which can be set up with the $`\varphi `$’s as building blocks.
A few comments are now in order. First, clearly, three distinct eigenvalues are just expected for two pairs, corresponding to zero, one and two pairs broken respectively. By extension, no matter what the value of $`\mathrm{\Omega }`$ is, the distinct roots of the secular equation will always be $`n+1`$. Their degeneracy can be found, as above stated, by inspecting the associated wave functions. However it can be understood on the basis of the following counting rule: the lowest eigenvalue is always unique
$$_0(n)\delta _0=1$$
and corresponds to the most collective eigenvector $`\mathrm{\Psi }_0`$ (see Eqs.20-26). The second eigenvector $`\mathrm{\Psi }_2`$ has a broken pair: since the states available for a pair are $`\mathrm{\Omega }`$, the broken pair can sit in any of them, except for the one corresponding to $`\mathrm{\Psi }_0`$: hence its degeneracy is
$$_2(n)\delta _2=\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{1}\right)1=\mathrm{\Omega }1.$$
The next eigenvector $`\mathrm{\Psi }_4`$ has two pairs broken. The associated degeneracy is obtained by subtracting, out of the states available for two pairs, the states already occupied by $`\mathrm{\Psi }_0`$ and $`\mathrm{\Psi }_2`$, namely
$$_4(n)\delta _4=\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{2}\right)1\left[\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{1}\right)1\right]=\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{2}\right)\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{1}\right).$$
In general the following degeneracy for the eigenvalues corresponding to seniority $`v`$ is found to hold (of course, a binomial coefficient with a negative lower index vanishes)
$$\delta _v=\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{\frac{v}{2}}\right)1\left[\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{1}\right)1\right]\left[\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{2}\right)\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{1}\right)\right]\mathrm{}\left[\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{\frac{v}{2}1}\right)\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{\frac{v}{2}2}\right)\right]$$
$$_v(n)\delta _v=\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{\frac{v}{2}}\right)\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{\frac{v}{2}1}\right),$$
(29)
$`v`$ being an even non-negative number referred to as seniority, to comply with the existing literature.
Finally we explore the general case. Here the structure of the symmetric matrix of dimension $`\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{n}\right)`$ associated with $`H_P^{(j)}`$ is easily found to be
$$\left(\begin{array}{ccc}+n& & 01\\ & \mathrm{}& \\ 01& & +n\end{array}\right),$$
(30)
where the symbol $`01`$ means that the upper (and the lower) triangle of the matrix is filled with zeros and ones. Indeed the matrix elements of $`H_P^{(j)}`$ turn out to be one, when the bra and the ket differ by the quantum state of one (out of $`n`$) pair, otherwise they vanish. The diagonal matrix elements simply count the number of pairs.
An elementary combinatorial analysis shows that the number of ones in each row (column) of the matrix is given by $`n(\mathrm{\Omega }n)`$. Indeed a non vanishing matrix element has the row specified by $`n`$ indices whereas, of the indices identifying the column, $`n1`$ should be extracted from the $`n`$ ones fixing the row in all the possible ways, which amounts to $`n`$ possibilities. The missing index should then be selected among the remaining $`\mathrm{\Omega }n`$ ones: hence the formula $`n(\mathrm{\Omega }n)`$ for the number of ones.
No matter where these are located <sup>*</sup><sup>*</sup>*As already noticed, any ordering of the states should of course lead to the same results. There are $`\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{n}\right)!`$ ordering alltogheter. Note however that for $`\mathrm{\Omega }`$ and $`n`$ large the number of matrix with $`n(\mathrm{\Omega }n)`$ ones in each row and column is larger than $`\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{n}\right)!`$., in the determinant of the matrix the first row (or column) can be replaced with the sum of all the rows (or columns). This leads to a new determinant with, e.g., the first row (or column) filled with ones and having the expression $`+n+n(\mathrm{\Omega }n)`$ factorized. Hence
1. the lowest eigenvalue is given by the integer
$$=n(1+\mathrm{\Omega }n)=nn(\mathrm{\Omega }n);$$
(31)
2. it is fixed uniquely by $`n`$ and by the number of ones present in each row (or column) of the matrix.
Next notice that, setting $`=0`$ in (30), for a given $`\mathrm{\Omega }`$, the number of pairs $`n(\frac{\mathrm{\Omega }}{2})`$, beyond fixing the dimension of the matrix and the number of ones in each row (column), also represents the common value of the diagonal elements. We have numerically checked that the determinant of such a matrix always vanishes. Hence it follows that $`=0`$ is always an eigenvalue of $`H_P^{(j)}`$ and actually the upper one. Indeed it corresponds to the situation where all the pairs are broken.
Let us now insert into the matrix (30) the lowest eigenvalue (31). We obtain
$$det\left(\begin{array}{ccc}n(\mathrm{\Omega }n)& & 01\\ & \mathrm{}& \\ 01& & n(\mathrm{\Omega }n)\end{array}\right)=0.$$
(32)
We extensively checked by direct computation that the above determinant continues to vanish under the replacement
$$n(\mathrm{\Omega }n)n(\mathrm{\Omega }n)+\frac{v}{2}\left(\mathrm{\Omega }\frac{v}{2}+1\right).$$
Hence the characteristic equation in the general case turns out to be
$$\underset{v/2=0}{\overset{n}{}}\left[+n(\mathrm{\Omega }n+1)\frac{v}{2}(\mathrm{\Omega }\frac{v}{2}+1)\right]^{\delta _v}=0$$
(33)
and the whole spectrum of the pairing hamiltonian follows.
The associated wave functions will be linear combinations of $`\left(\genfrac{}{}{0pt}{}{\mathrm{\Omega }}{n}\right)`$ monomials, each monomial being the product of $`n`$ $`\varphi `$’s, as usual.
In conclusion, we like to shortly address the problem of the angular momentum of our eigenstates, which belong to a definite value of energy, seniority and third component of the angular momentum, but not of the angular momentum. However, their building blocks, the monomials, can be expressed as superpositions of “collective” bosons $`\mathrm{\Phi }_J`$ of definite angular momentum, according to
$$\varphi _m=\sqrt{2}\underset{J}{}jm,jm|J0\mathrm{\Phi }_J.$$
(34)
It should be realized however that the nihilpotency of the variables $`\varphi `$’s induces a set of constraints to be fulfilled by the “collective” variables $`\mathrm{\Phi }_J`$. As an example, consider again the case $`\mathrm{\Omega }=4,n=2`$. Here the four constraints $`\varphi _m^2=0`$ are translated in the following four equations for the variables $`\mathrm{\Phi }_J`$ with $`J=0,2,4`$ and $`6`$:
$`\mathrm{\Phi }_0^2=\mathrm{\Phi }_2^2\mathrm{\Phi }_4^2\mathrm{\Phi }_6^2`$ (35)
$`77\sqrt{21}\mathrm{\Phi }_0\mathrm{\Phi }_2=88\mathrm{\Phi }_2^210\mathrm{\Phi }_4^2+98\mathrm{\Phi }_6^244\sqrt{33}\mathrm{\Phi }_2\mathrm{\Phi }_435\sqrt{21}\mathrm{\Phi }_4\mathrm{\Phi }_6`$ (36)
$`231\sqrt{77}\mathrm{\Phi }_0\mathrm{\Phi }_4=726\mathrm{\Phi }_2^2+432\mathrm{\Phi }_4^2+294\mathrm{\Phi }_6^2`$
$`20\sqrt{33}\mathrm{\Phi }_2\mathrm{\Phi }_4105\sqrt{77}\mathrm{\Phi }_2\mathrm{\Phi }_6+280\sqrt{21}\mathrm{\Phi }_4\mathrm{\Phi }_6`$ (37)
$`11\sqrt{33}\mathrm{\Phi }_0\mathrm{\Phi }_6=20\mathrm{\Phi }_4^220\mathrm{\Phi }_6^25\sqrt{33}\mathrm{\Phi }_2\mathrm{\Phi }_44\sqrt{77}\mathrm{\Phi }_2\mathrm{\Phi }_6+4\sqrt{21}\mathrm{\Phi }_4\mathrm{\Phi }_6.`$ (38)
Hence the $`v=0`$ eigenstate (26), when expressed through the $`\mathrm{\Phi }_J`$, reads
$$\mathrm{\Psi }_0(n=2)=𝒩_0\frac{1}{2}\left(3\mathrm{\Phi }_0^2+\mathrm{\Phi }_2^2+\mathrm{\Phi }_4^2+\mathrm{\Phi }_6^2\right)$$
(39)
which, exploiting (35), actually becomes
$$\mathrm{\Psi }_0(n=2)=2𝒩_0\mathrm{\Phi }_0^2.$$
(40)
The above is the product of two $`s`$-bosons, as the variable $`\mathrm{\Phi }_0`$ is often referred to, or of the wavefunctions of two unbroken pairs.
Analogously the $`v=2`$ states, exploiting the constraints (36-38), turns out to read
$`\mathrm{\Psi }_2(n=2)=𝒩_2{\displaystyle \frac{2}{\sqrt{3}\sqrt{7}\sqrt{11}}}[(4a_12a_2a_3)\sqrt{11}\mathrm{\Phi }_0\mathrm{\Phi }_2+`$
$`+(3a_12a_28a_3)\sqrt{3}\mathrm{\Phi }_0\mathrm{\Phi }_4+(2a_15a_2+2a_3)\sqrt{7}\mathrm{\Phi }_0\mathrm{\Phi }_6],`$ (41)
which shows that a particular choice of the parameters allows to select a state of good angular momentum $`(J=2,4,6)`$. Hence this state can be reduced to the product, e.g., of an $`s`$ and $`d`$ boson, the latter corresponding to the wave function of the broken pair.
More subtle is the $`v=4`$ case. Here the procedure above outlined leads to the following expression of the wavefunction
$`\mathrm{\Psi }_4(n=2)=𝒩_4`$
$`[b_2({\displaystyle \frac{4}{7}}\mathrm{\Phi }_2^2{\displaystyle \frac{5}{77}}\mathrm{\Phi }_4^2+{\displaystyle \frac{7}{11}}\mathrm{\Phi }_6^2+{\displaystyle \frac{32}{7\sqrt{33}}}\mathrm{\Phi }_2\mathrm{\Phi }_4+{\displaystyle \frac{36}{3\sqrt{77}}}\mathrm{\Phi }_2\mathrm{\Phi }_6+{\displaystyle \frac{8}{11\sqrt{21}}}\mathrm{\Phi }_4\mathrm{\Phi }_6)`$
$`+b_3({\displaystyle \frac{5}{7}}\mathrm{\Phi }_2^2{\displaystyle \frac{5}{7}}\mathrm{\Phi }_4^2{\displaystyle \frac{40}{7\sqrt{33}}}\mathrm{\Phi }_2\mathrm{\Phi }_4{\displaystyle \frac{4}{\sqrt{77}}}\mathrm{\Phi }_2\mathrm{\Phi }_6+{\displaystyle \frac{4}{\sqrt{21}}}\mathrm{\Phi }_4\mathrm{\Phi }_6)]`$ (42)
where, remarkably, the collective $`\mathrm{\Phi }_0`$ does not appear, as it should, since (42) is associated with two broken pairs.
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